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add vae, lcm, t2i model path
Browse files- lcm/lcm-lora-sdv1-5/.gitattributes +35 -0
- lcm/lcm-lora-sdv1-5/README.md +217 -0
- lcm/lcm-lora-sdv1-5/image.png +0 -0
- lcm/lcm-lora-sdv1-5/pytorch_lora_weights.safetensors +3 -0
- t2i/sd1.5/majicmixRealv6Fp16/feature_extractor/preprocessor_config.json +28 -0
- t2i/sd1.5/majicmixRealv6Fp16/model_index.json +29 -0
- t2i/sd1.5/majicmixRealv6Fp16/model_index.json.bk +33 -0
- t2i/sd1.5/majicmixRealv6Fp16/scheduler/scheduler_config.json +18 -0
- t2i/sd1.5/majicmixRealv6Fp16/text_encoder/config.json +25 -0
- t2i/sd1.5/majicmixRealv6Fp16/text_encoder/pytorch_model.bin +3 -0
- t2i/sd1.5/majicmixRealv6Fp16/tokenizer/merges.txt +0 -0
- t2i/sd1.5/majicmixRealv6Fp16/tokenizer/special_tokens_map.json +24 -0
- t2i/sd1.5/majicmixRealv6Fp16/tokenizer/tokenizer_config.json +34 -0
- t2i/sd1.5/majicmixRealv6Fp16/tokenizer/vocab.json +0 -0
- t2i/sd1.5/majicmixRealv6Fp16/unet/config.json +60 -0
- t2i/sd1.5/majicmixRealv6Fp16/unet/diffusion_pytorch_model.bin +3 -0
- t2i/sd1.5/majicmixRealv6Fp16/vae/config.json +30 -0
- t2i/sd1.5/majicmixRealv6Fp16/vae/diffusion_pytorch_model.bin +3 -0
- vae/sd-vae-ft-mse/.gitattributes +33 -0
- vae/sd-vae-ft-mse/README.md +83 -0
- vae/sd-vae-ft-mse/config.json +29 -0
- vae/sd-vae-ft-mse/diffusion_pytorch_model.bin +3 -0
- vae/sd-vae-ft-mse/diffusion_pytorch_model.safetensors +3 -0
lcm/lcm-lora-sdv1-5/.gitattributes
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lcm/lcm-lora-sdv1-5/README.md
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---
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library_name: diffusers
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base_model: runwayml/stable-diffusion-v1-5
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tags:
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- lora
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- text-to-image
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license: openrail++
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inference: false
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---
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# Latent Consistency Model (LCM) LoRA: SDv1-5
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Latent Consistency Model (LCM) LoRA was proposed in [LCM-LoRA: A universal Stable-Diffusion Acceleration Module](https://arxiv.org/abs/2311.05556)
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by *Simian Luo, Yiqin Tan, Suraj Patil, Daniel Gu et al.*
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It is a distilled consistency adapter for [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5) that allows
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to reduce the number of inference steps to only between **2 - 8 steps**.
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| Model | Params / M |
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|----------------------------------------------------------------------------|------------|
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| [**lcm-lora-sdv1-5**](https://huggingface.co/latent-consistency/lcm-lora-sdv1-5) | **67.5** |
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| [lcm-lora-ssd-1b](https://huggingface.co/latent-consistency/lcm-lora-ssd-1b) | 105 |
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| [lcm-lora-sdxl](https://huggingface.co/latent-consistency/lcm-lora-sdxl) | 197M |
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## Usage
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LCM-LoRA is supported in 🤗 Hugging Face Diffusers library from version v0.23.0 onwards. To run the model, first
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install the latest version of the Diffusers library as well as `peft`, `accelerate` and `transformers`.
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audio dataset from the Hugging Face Hub:
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```bash
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pip install --upgrade pip
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pip install --upgrade diffusers transformers accelerate peft
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```
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***Note: For detailed usage examples we recommend you to check out our official [LCM-LoRA docs](https://huggingface.co/docs/diffusers/main/en/using-diffusers/inference_with_lcm_lora)***
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### Text-to-Image
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The adapter can be loaded with SDv1-5 or deviratives. Here we use [`Lykon/dreamshaper-7`](https://huggingface.co/Lykon/dreamshaper-7). Next, the scheduler needs to be changed to [`LCMScheduler`](https://huggingface.co/docs/diffusers/v0.22.3/en/api/schedulers/lcm#diffusers.LCMScheduler) and we can reduce the number of inference steps to just 2 to 8 steps.
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Please make sure to either disable `guidance_scale` or use values between 1.0 and 2.0.
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```python
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import torch
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from diffusers import LCMScheduler, AutoPipelineForText2Image
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model_id = "Lykon/dreamshaper-7"
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adapter_id = "latent-consistency/lcm-lora-sdv1-5"
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pipe = AutoPipelineForText2Image.from_pretrained(model_id, torch_dtype=torch.float16, variant="fp16")
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pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
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pipe.to("cuda")
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# load and fuse lcm lora
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pipe.load_lora_weights(adapter_id)
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pipe.fuse_lora()
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prompt = "Self-portrait oil painting, a beautiful cyborg with golden hair, 8k"
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# disable guidance_scale by passing 0
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image = pipe(prompt=prompt, num_inference_steps=4, guidance_scale=0).images[0]
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```
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![](./image.png)
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### Image-to-Image
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LCM-LoRA can be applied to image-to-image tasks too. Let's look at how we can perform image-to-image generation with LCMs. For this example we'll use the [dreamshaper-7](https://huggingface.co/Lykon/dreamshaper-7) model and the LCM-LoRA for `stable-diffusion-v1-5 `.
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```python
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import torch
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from diffusers import AutoPipelineForImage2Image, LCMScheduler
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from diffusers.utils import make_image_grid, load_image
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pipe = AutoPipelineForImage2Image.from_pretrained(
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"Lykon/dreamshaper-7",
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torch_dtype=torch.float16,
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variant="fp16",
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).to("cuda")
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# set scheduler
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pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
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# load LCM-LoRA
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pipe.load_lora_weights("latent-consistency/lcm-lora-sdv1-5")
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pipe.fuse_lora()
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# prepare image
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url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
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init_image = load_image(url)
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prompt = "Astronauts in a jungle, cold color palette, muted colors, detailed, 8k"
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# pass prompt and image to pipeline
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generator = torch.manual_seed(0)
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image = pipe(
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prompt,
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image=init_image,
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num_inference_steps=4,
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guidance_scale=1,
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strength=0.6,
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generator=generator
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).images[0]
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make_image_grid([init_image, image], rows=1, cols=2)
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```
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![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_sdv1-5_i2i.png)
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### Inpainting
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LCM-LoRA can be used for inpainting as well.
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```python
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import torch
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from diffusers import AutoPipelineForInpainting, LCMScheduler
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from diffusers.utils import load_image, make_image_grid
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pipe = AutoPipelineForInpainting.from_pretrained(
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"runwayml/stable-diffusion-inpainting",
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torch_dtype=torch.float16,
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variant="fp16",
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).to("cuda")
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# set scheduler
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pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
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# load LCM-LoRA
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pipe.load_lora_weights("latent-consistency/lcm-lora-sdv1-5")
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pipe.fuse_lora()
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# load base and mask image
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init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
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mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
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# generator = torch.Generator("cuda").manual_seed(92)
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prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
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generator = torch.manual_seed(0)
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image = pipe(
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prompt=prompt,
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image=init_image,
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mask_image=mask_image,
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generator=generator,
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num_inference_steps=4,
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guidance_scale=4,
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).images[0]
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make_image_grid([init_image, mask_image, image], rows=1, cols=3)
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```
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![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_sdv1-5_inpainting.png)
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### ControlNet
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For this example, we'll use the SD-v1-5 model and the LCM-LoRA for SD-v1-5 with canny ControlNet.
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```python
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import torch
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import cv2
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import numpy as np
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from PIL import Image
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from diffusers import StableDiffusionControlNetPipeline, ControlNetModel, LCMScheduler
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from diffusers.utils import load_image
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image = load_image(
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"https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png"
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).resize((512, 512))
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image = np.array(image)
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low_threshold = 100
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high_threshold = 200
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image = cv2.Canny(image, low_threshold, high_threshold)
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image = image[:, :, None]
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image = np.concatenate([image, image, image], axis=2)
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canny_image = Image.fromarray(image)
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controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
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pipe = StableDiffusionControlNetPipeline.from_pretrained(
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"runwayml/stable-diffusion-v1-5",
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controlnet=controlnet,
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torch_dtype=torch.float16,
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safety_checker=None,
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variant="fp16"
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).to("cuda")
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# set scheduler
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pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
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# load LCM-LoRA
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pipe.load_lora_weights("latent-consistency/lcm-lora-sdv1-5")
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generator = torch.manual_seed(0)
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image = pipe(
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"the mona lisa",
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image=canny_image,
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num_inference_steps=4,
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guidance_scale=1.5,
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controlnet_conditioning_scale=0.8,
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cross_attention_kwargs={"scale": 1},
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generator=generator,
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).images[0]
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make_image_grid([canny_image, image], rows=1, cols=2)
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```
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![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_sdv1-5_controlnet.png)
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## Speed Benchmark
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TODO
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## Training
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TODO
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lcm/lcm-lora-sdv1-5/image.png
ADDED
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t2i/sd1.5/majicmixRealv6Fp16/feature_extractor/preprocessor_config.json
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t2i/sd1.5/majicmixRealv6Fp16/model_index.json
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t2i/sd1.5/majicmixRealv6Fp16/model_index.json.bk
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t2i/sd1.5/majicmixRealv6Fp16/scheduler/scheduler_config.json
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t2i/sd1.5/majicmixRealv6Fp16/text_encoder/config.json
ADDED
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t2i/sd1.5/majicmixRealv6Fp16/text_encoder/pytorch_model.bin
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t2i/sd1.5/majicmixRealv6Fp16/tokenizer/merges.txt
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The diff for this file is too large to render.
See raw diff
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t2i/sd1.5/majicmixRealv6Fp16/tokenizer/special_tokens_map.json
ADDED
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t2i/sd1.5/majicmixRealv6Fp16/tokenizer/tokenizer_config.json
ADDED
@@ -0,0 +1,34 @@
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t2i/sd1.5/majicmixRealv6Fp16/tokenizer/vocab.json
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t2i/sd1.5/majicmixRealv6Fp16/unet/config.json
ADDED
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t2i/sd1.5/majicmixRealv6Fp16/unet/diffusion_pytorch_model.bin
ADDED
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t2i/sd1.5/majicmixRealv6Fp16/vae/config.json
ADDED
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|
t2i/sd1.5/majicmixRealv6Fp16/vae/diffusion_pytorch_model.bin
ADDED
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version https://git-lfs.github.com/spec/v1
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oid sha256:e74bb1390c97f146afccbbb42e65c1d295d0f1e61f7886a240b8a4e5ad622328
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vae/sd-vae-ft-mse/.gitattributes
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*.7z filter=lfs diff=lfs merge=lfs -text
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diffusion_pytorch_model.safetensors filter=lfs diff=lfs merge=lfs -text
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vae/sd-vae-ft-mse/README.md
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---
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license: mit
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tags:
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- stable-diffusion
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- stable-diffusion-diffusers
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inference: false
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---
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# Improved Autoencoders
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## Utilizing
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These weights are intended to be used with the [🧨 diffusers library](https://github.com/huggingface/diffusers). If you are looking for the model to use with the original [CompVis Stable Diffusion codebase](https://github.com/CompVis/stable-diffusion), [come here](https://huggingface.co/stabilityai/sd-vae-ft-mse-original).
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#### How to use with 🧨 diffusers
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You can integrate this fine-tuned VAE decoder to your existing `diffusers` workflows, by including a `vae` argument to the `StableDiffusionPipeline`
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```py
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from diffusers.models import AutoencoderKL
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from diffusers import StableDiffusionPipeline
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model = "CompVis/stable-diffusion-v1-4"
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vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse")
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pipe = StableDiffusionPipeline.from_pretrained(model, vae=vae)
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```
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## Decoder Finetuning
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We publish two kl-f8 autoencoder versions, finetuned from the original [kl-f8 autoencoder](https://github.com/CompVis/latent-diffusion#pretrained-autoencoding-models) on a 1:1 ratio of [LAION-Aesthetics](https://laion.ai/blog/laion-aesthetics/) and LAION-Humans, an unreleased subset containing only SFW images of humans. The intent was to fine-tune on the Stable Diffusion training set (the autoencoder was originally trained on OpenImages) but also enrich the dataset with images of humans to improve the reconstruction of faces.
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The first, _ft-EMA_, was resumed from the original checkpoint, trained for 313198 steps and uses EMA weights. It uses the same loss configuration as the original checkpoint (L1 + LPIPS).
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The second, _ft-MSE_, was resumed from _ft-EMA_ and uses EMA weights and was trained for another 280k steps using a different loss, with more emphasis
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on MSE reconstruction (MSE + 0.1 * LPIPS). It produces somewhat ``smoother'' outputs. The batch size for both versions was 192 (16 A100s, batch size 12 per GPU).
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To keep compatibility with existing models, only the decoder part was finetuned; the checkpoints can be used as a drop-in replacement for the existing autoencoder.
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_Original kl-f8 VAE vs f8-ft-EMA vs f8-ft-MSE_
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## Evaluation
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### COCO 2017 (256x256, val, 5000 images)
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| Model | train steps | rFID | PSNR | SSIM | PSIM | Link | Comments
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|----------|---------|------|--------------|---------------|---------------|-----------------------------------------------------------------------------------|-------------------------------------------------------------------------------------------------|
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| | | | | | | | |
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| original | 246803 | 4.99 | 23.4 +/- 3.8 | 0.69 +/- 0.14 | 1.01 +/- 0.28 | https://ommer-lab.com/files/latent-diffusion/kl-f8.zip | as used in SD |
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| ft-EMA | 560001 | 4.42 | 23.8 +/- 3.9 | 0.69 +/- 0.13 | 0.96 +/- 0.27 | https://huggingface.co/stabilityai/sd-vae-ft-ema-original/resolve/main/vae-ft-ema-560000-ema-pruned.ckpt | slightly better overall, with EMA |
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| ft-MSE | 840001 | 4.70 | 24.5 +/- 3.7 | 0.71 +/- 0.13 | 0.92 +/- 0.27 | https://huggingface.co/stabilityai/sd-vae-ft-mse-original/resolve/main/vae-ft-mse-840000-ema-pruned.ckpt | resumed with EMA from ft-EMA, emphasis on MSE (rec. loss = MSE + 0.1 * LPIPS), smoother outputs |
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### LAION-Aesthetics 5+ (256x256, subset, 10000 images)
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| Model | train steps | rFID | PSNR | SSIM | PSIM | Link | Comments
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|----------|-----------|------|--------------|---------------|---------------|-----------------------------------------------------------------------------------|-------------------------------------------------------------------------------------------------|
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| | | | | | | | |
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| original | 246803 | 2.61 | 26.0 +/- 4.4 | 0.81 +/- 0.12 | 0.75 +/- 0.36 | https://ommer-lab.com/files/latent-diffusion/kl-f8.zip | as used in SD |
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| ft-EMA | 560001 | 1.77 | 26.7 +/- 4.8 | 0.82 +/- 0.12 | 0.67 +/- 0.34 | https://huggingface.co/stabilityai/sd-vae-ft-ema-original/resolve/main/vae-ft-ema-560000-ema-pruned.ckpt | slightly better overall, with EMA |
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| ft-MSE | 840001 | 1.88 | 27.3 +/- 4.7 | 0.83 +/- 0.11 | 0.65 +/- 0.34 | https://huggingface.co/stabilityai/sd-vae-ft-mse-original/resolve/main/vae-ft-mse-840000-ema-pruned.ckpt | resumed with EMA from ft-EMA, emphasis on MSE (rec. loss = MSE + 0.1 * LPIPS), smoother outputs |
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### Visual
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_Visualization of reconstructions on 256x256 images from the COCO2017 validation dataset._
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<p align="center">
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<br>
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<b>
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256x256: ft-EMA (left), ft-MSE (middle), original (right)</b>
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00025_merged.png />
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00011_merged.png />
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00037_merged.png />
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00043_merged.png />
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00053_merged.png />
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</p>
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<p align="center">
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<img src=https://huggingface.co/stabilityai/stable-diffusion-decoder-finetune/resolve/main/eval/ae-decoder-tuning-reconstructions/merged/00029_merged.png />
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</p>
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vae/sd-vae-ft-mse/config.json
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{
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"_class_name": "AutoencoderKL",
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"_diffusers_version": "0.4.2",
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"act_fn": "silu",
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"block_out_channels": [
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128,
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256,
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512,
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512
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],
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"down_block_types": [
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"DownEncoderBlock2D",
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"DownEncoderBlock2D",
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"DownEncoderBlock2D",
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"DownEncoderBlock2D"
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],
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"in_channels": 3,
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"latent_channels": 4,
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"layers_per_block": 2,
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"norm_num_groups": 32,
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"out_channels": 3,
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"sample_size": 256,
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"up_block_types": [
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"UpDecoderBlock2D",
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"UpDecoderBlock2D",
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"UpDecoderBlock2D",
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"UpDecoderBlock2D"
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]
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}
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vae/sd-vae-ft-mse/diffusion_pytorch_model.bin
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vae/sd-vae-ft-mse/diffusion_pytorch_model.safetensors
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