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Parent(s):
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Browse files- en/4284.html.txt +85 -0
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en/4284.html.txt
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Opera is a freeware web browser for Microsoft Windows, Android, iOS, macOS, and Linux operating systems[7][8] developed by Opera Software.[9][10][11] Opera is a Chromium-based browser using the Blink layout engine. It differentiates itself because of a distinct user interface and other features.
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Opera was conceived at Telenor as a research project in 1994 and was bought by Opera Software in 1995. It was a commercial software for the first ten years and had its own proprietary Presto layout engine. The Presto versions of Opera received many awards, but Presto development ended after a milestone transition to Chromium in 2013.
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There are also three mobile versions called Opera Mobile, Opera Touch,[12] and Opera Mini. A gaming browser called Opera GX was launched on 11 June 2019.[13]
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Opera began in 1994 as a research project at Telenor, the largest Norwegian telecommunications company.
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In 1995, it branched out into a separate company named Opera Software.[14]
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Opera was first publicly released in 1996 with version 2.10,[15] which only ran on Microsoft Windows 95 at that time.[16]
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In an attempt to capitalize on the emerging market for Internet-connected handheld devices, a project to port Opera to mobile device platforms was started in 1998.[16]
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Opera 4.0, released in 2000,[15] included a new cross-platform core that facilitated the creation of editions of Opera for multiple operating systems and platforms.[17]
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Up to this point, Opera was trialware and had to be purchased after the trial period ended. Version 5.0 (released in 2000) saw the end of this requirement. Instead, Opera became ad-sponsored, displaying advertisements to users who had not paid for it.[18] Later versions of Opera gave the user the choice of seeing banner ads or targeted text advertisements from Google.
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With version 8.5 (released in 2005) the advertisements were completely removed and the primary financial support for the browser came through revenue from Google (which is by contract Opera's default search engine).[19]
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Among the new features introduced in version 9.1 (released in 2006) was fraud protection using technology from GeoTrust, a digital certificate provider, and PhishTank, an organization that tracks known phishing web sites.[20] This feature was further improved and expanded in version 9.5, when GeoTrust was replaced with Netcraft, and malware protection from Haute Secure was added.[21]
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In 2006, Opera Software ASA was released as well as Internet Channel and Nintendo DS Browser for Nintendo's DS and Wii gaming systems.[22][23][24][25]
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A new JavaScript engine called Carakan, after the Javanese alphabet, was introduced with version 10.50.[26] According to Opera Software, Carakan made Opera 10.50 more than seven times faster in SunSpider than Opera 10.10.[27][28][29]
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On the 16th of December 2010, Opera 11 was released, featuring extensions,[30] tab stacking (where dragging one tab over another allows creating a group of tabs), visual mouse gestures and changes to the address bar.[31]
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Opera 12 was released on 14 June 2012.[32]
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On 12 February 2013, Opera Software announced that it would drop its own Presto layout engine in favour of WebKit as implemented by Google's Chrome browser, using code from the Chromium project. Opera Software planned as well to contribute code to WebKit.[33] On 3 April 2013, Google announced that it would fork components from WebKit to form a new layout engine known as Blink; the same day, Opera Software confirmed that it would follow Google in implementing the Blink layout engine.[34]
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On 28 May 2013, a beta release of Opera 15 was made available,[35] the first version is based on the Chromium project.[36][37] Many distinctive Opera features of the previous versions were dropped, and Opera Mail was separated into a standalone application derived from Opera 12.[38]
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In November 2016, the original Norwegian owner of Opera sold his stake in the business to a Chinese consortium under the name Golden Brick Capital Private Equity Fund I Limited Partnership for $600 million.[39][40][41] An earlier deal was not approved by regulators.[42]
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In January 2017, the source code of Opera 12.15 (one of the last few versions that was still based on the Presto layout engine) was leaked.[43]
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To demonstrate how radically different a browser could look, Opera Neon, dubbed a "concept browser", was released in January 2017. PC World compared it to demo models that automakers and hardware vendors release to show their visions of the future. Instead of a Speed Dial (also explained in the following chapter "Features"), it displays the frequently accessed websites in resemblance to a desktop with computer icons scattered all over it in an artistic formation.[44][45]
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Opera has originated features later adopted by other web browsers, including: Speed Dial, pop-up blocking, re-opening recently closed pages, private browsing, and tabbed browsing.[46][47] Opera includes a bookmarks bar and a download manager. Opera also has "Speed Dial", which allows the user to add an unlimited number of pages shown in thumbnail form in a page displayed when a new tab is opened. Speed Dial allows the user to more easily navigate to the selected web pages.[46][47][48]
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It is possible to control some aspects of the browser using the keyboard shortcuts.[49] Page zooming allows text, images and other content such as Adobe Flash Player, Java platform and Scalable Vector Graphics to be increased or decreased in size to help those with impaired vision.[50]
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Opera Software claims that when the Opera Turbo mode is enabled, the compression servers compress requested web pages (except HTTPS pages) by up to 50%, depending upon the content, before sending them to the users.[51] This process reduces the amount of data transferred and is particularly useful for crowded or slow network connections, making web pages load faster or when there are costs dependent on the total amount of data usage.[51] This technique is also used in Opera Mini for mobile devices[52] and smartwatches.[53]
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One security feature is the option to delete private data, such as HTTP cookies, browsing history, items in cache and passwords with the click of a button.[54] This lets users erase personal data after browsing from a shared computer.
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When visiting a site, Opera displays a security badge in the address bar which shows details about the website, including security certificates.[55] Opera's fraud and malware protection warns the user about suspicious web pages and is enabled by default. It checks the requested page against several databases of known phishing and malware websites, called blacklists.[55]
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In January 2007, Asa Dotzler of the competing Mozilla Corporation accused Opera Software of downplaying information about security vulnerabilities in Opera, (which were fixed in December 2006). Dotzler claimed that users were not clearly informed of security vulnerabilities that were present in the previous version of Opera and thus they would not realize that they needed to upgrade to the latest version or else risk being exploited by hackers.[56] Opera Software responded to these accusations on the next day.[57]
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In 2016, a free virtual private network (VPN) service was implemented in the browser.[58] Opera said that this would allow encrypted access to websites otherwise blocked, and provide security on public WiFi networks.[59] Either VPN or Turbo can be enabled, but not both.[60] It was later determined that the browser "VPN" operated the same as a proxy rather than other VPN services.[61]
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In July 2018, Opera was listed on the NASDAQ stock exchange [62] in New York City at an initial offering of $12 per share.
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In 2018, a built-in cryptocurrency wallet to the Opera Web Browser was released.[63] Announcing that they would be the first browser with a built-in Crypto Wallet.[64] On 13 December 2018 they released a video showing many decentralized applications like Cryptokitties running on the Android version of the Opera Web Browser.[65]
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In March 2020, Opera updated its Android browser to access .crypto domains, making it the first browser to be able to support a domain name system (DNS) which is not part of the traditional DNS directly without the need of a plugin or add-on.[66] This was through a collaboration with a San Francisco based startup, Unstoppable Domains.[67][68]
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Opera was one of the first browsers to support Cascading Style Sheets (CSS).[69]
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Opera Software uses a release cycle consisting of three "streams" (which correspond to phases of development) that can be downloaded and installed independently of each other: "developer", "beta" and "stable". New features are first introduced in the developer build, then, depending on user feedback, may progress to the beta version and eventually be released.[70]
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The developer stream allows early testing of new features, mainly targeting developers, extension creators, and early adopters. Opera developer is not intended for everyday browsing as it is unstable and is prone to failure or crashing, but it enables advanced users to try out new features that are still under development, without affecting their normal installation of the browser. New versions of the browser are released frequently, generally a few times a week.[71]
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The beta stream, formerly known as "Opera Next", is a feature complete package, allowing stability and quality to mature before the final release. A new version is released every couple of weeks.[72]
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Both streams can be installed alongside the official release without interference. Each has a different icon to help the user distinguish between the variants.
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In 2005, Adobe Systems opted to integrate Opera's rendering engine, Presto, into its Adobe Creative Suite applications. Opera technology was employed in Adobe GoLive, Adobe Photoshop, Adobe Dreamweaver, and other components of the Adobe Creative Suite.[74][75] Opera's layout engine is also found in Virtual Mechanics SiteSpinner Pro.[76] The Internet Channel is a version of the Opera 9 web browser for use on the Nintendo Wii created by Opera Software and Nintendo.[77] Opera Software is also implemented in the Nintendo DS Browser for Nintendo's handheld systems.
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Versions with the Presto layout engine have been positively reviewed,[78][79][80] although they have been criticized for website compatibility issues.[81][82] Because of this issue, Opera 8.01 and higher had included workarounds to help certain popular but problematic web sites display properly.[83][84]
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Versions with the Blink layout engine have been criticized by some users for missing features such as UI customization, and for abandoning Opera Software's own Presto layout engine.[85][86][87][88] Despite that, versions with the Blink layout engine have been praised for being fast and stable, for handling the latest web standards and for having a better website compatibility and a modern-style user interface.[89][90][91]
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Over the years, Opera for personal computers has received several awards. These awards include:[92]
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Related web browsers:
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Web browsers:
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Essential features:
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Operation or Operations may refer to:
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Asia-Pacific
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Mediterranean and Middle East
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Other campaigns
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Coups
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World War II (WWII or WW2), also known as the Second World War, was a global war that lasted from 1939 to 1945. It involved the vast majority of the world's countries—including all the great powers—forming two opposing military alliances: the Allies and the Axis. In a state of total war, directly involving more than 100 million people from more than 30 countries, the major participants threw their entire economic, industrial, and scientific capabilities behind the war effort, blurring the distinction between civilian and military resources. World War II was the deadliest conflict in human history, marked by 70 to 85 million fatalities. Tens of millions of people died due to genocides (including the Holocaust), premeditated death from starvation, massacres, and disease. Aircraft played a major role in the conflict, including in the use of strategic bombing of population centres, and the only uses of nuclear weapons in war.
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World War II is generally considered to have begun on 1 September 1939, with the invasion of Poland by Germany and subsequent declarations of war on Germany by France and the United Kingdom. From late 1939 to early 1941, in a series of campaigns and treaties, Germany conquered or controlled much of continental Europe, and formed the Axis alliance with Italy and Japan. Under the Molotov–Ribbentrop Pact of August 1939, Germany and the Soviet Union partitioned and annexed territories of their European neighbours: Poland, Finland, Romania and the Baltic states. Following the onset of campaigns in North Africa and East Africa, and the Fall of France in mid-1940, the war continued primarily between the European Axis powers and the British Empire, with war in the Balkans, the aerial Battle of Britain, the Blitz, and the Battle of the Atlantic. On 22 June 1941, Germany led the European Axis powers in an invasion of the Soviet Union, opening the largest land theatre of war in history and trapping the Axis, crucially the German Wehrmacht, in a war of attrition.
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Japan, which aimed to dominate Asia and the Pacific, was at war with the Republic of China by 1937. In December 1941, Japan launched a surprise attack on the United States as well as European colonies in East Asia and the Pacific. Following an immediate US declaration of war against Japan, supported by one from the UK, the European Axis powers declared war on the United States in solidarity with their ally. Japan soon captured much of the Western Pacific, but its advances were halted in 1942 after Japan lost the critical Battle of Midway; later, Germany and Italy were defeated in North Africa and at Stalingrad in the Soviet Union. Key setbacks in 1943—which included a series of German defeats on the Eastern Front, the Allied invasions of Sicily and Italy, and Allied offensives in the Pacific—cost the Axis its initiative and forced it into strategic retreat on all fronts. In 1944, the Western Allies invaded German-occupied France, while the Soviet Union regained its territorial losses and turned towards Germany and its allies. During 1944 and 1945, the Japanese suffered reversals in mainland Asia, while the Allies crippled the Japanese Navy and captured key Western Pacific islands.
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The war in Europe concluded with an invasion of Germany by the Western Allies and the Soviet Union, culminating in the capture of Berlin by Soviet troops, the suicide of Adolf Hitler and the German unconditional surrender on 8 May 1945. Following the Potsdam Declaration by the Allies on 26 July 1945 and the refusal of Japan to surrender on its terms, the United States dropped atomic bombs on the Japanese cities of Hiroshima and Nagasaki on 6 and 9 August, respectively. Faced with an imminent invasion of the Japanese archipelago, the possibility of additional atomic bombings, and the Soviet entry into the war against Japan and its invasion of Manchuria on 9 August, Japan announced its intention to surrender on 15 August 1945, cementing total victory in Asia for the Allies. In the wake of the war, Germany and Japan were occupied and war crimes tribunals were conducted against German and Japanese leaders.
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World War II changed the political alignment and social structure of the globe. The United Nations (UN) was established to foster international co-operation and prevent future conflicts, and the victorious great powers—China, France, the Soviet Union, the United Kingdom, and the United States—became the permanent members of its Security Council. The Soviet Union and the United States emerged as rival superpowers, setting the stage for the nearly half-century-long Cold War. In the wake of European devastation, the influence of its great powers waned, triggering the decolonisation of Africa and Asia. Most countries whose industries had been damaged moved towards economic recovery and expansion. Political integration, especially in Europe, began as an effort to forestall future hostilities, end pre-war enmities and forge a sense of common identity.
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The start of the war in Europe is generally held to be 1 September 1939,[1][2] beginning with the German invasion of Poland; the United Kingdom and France declared war on Germany two days later. The dates for the beginning of war in the Pacific include the start of the Second Sino-Japanese War on 7 July 1937,[3][4] or even the Japanese invasion of Manchuria on 19 September 1931.[5][6][7]
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Others follow the British historian A. J. P. Taylor, who held that the Sino-Japanese War and war in Europe and its colonies occurred simultaneously, and the two wars merged in 1941. This article uses the conventional dating. Other starting dates sometimes used for World War II include the Italian invasion of Abyssinia on 3 October 1935.[8] The British historian Antony Beevor views the beginning of World War II as the Battles of Khalkhin Gol fought between Japan and the forces of Mongolia and the Soviet Union from May to September 1939.[9]
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The exact date of the war's end is also not universally agreed upon. It was generally accepted at the time that the war ended with the armistice of 14 August 1945 (V-J Day), rather than the formal surrender of Japan, which was on 2 September 1945 that officially ended the war in Asia. A peace treaty with Japan was signed in 1951.[10] A treaty regarding Germany's future allowed the reunification of East and West Germany to take place in 1990 and resolved most post-World War II issues.[11] No formal peace treaty between Japan and the Soviet Union was ever signed.[12]
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World War I had radically altered the political European map, with the defeat of the Central Powers—including Austria-Hungary, Germany, Bulgaria and the Ottoman Empire—and the 1917 Bolshevik seizure of power in Russia, which eventually led to the founding of the Soviet Union. Meanwhile, the victorious Allies of World War I, such as France, Belgium, Italy, Romania, and Greece, gained territory, and new nation-states were created out of the collapse of Austria-Hungary and the Ottoman and Russian Empires.
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To prevent a future world war, the League of Nations was created during the 1919 Paris Peace Conference. The organisation's primary goals were to prevent armed conflict through collective security, military and naval disarmament, and settling international disputes through peaceful negotiations and arbitration.
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Despite strong pacifist sentiment after World War I,[13] its aftermath still caused irredentist and revanchist nationalism in several European states. These sentiments were especially marked in Germany because of the significant territorial, colonial, and financial losses incurred by the Treaty of Versailles. Under the treaty, Germany lost around 13 percent of its home territory and all its overseas possessions, while German annexation of other states was prohibited, reparations were imposed, and limits were placed on the size and capability of the country's armed forces.[14]
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The German Empire was dissolved in the German Revolution of 1918–1919, and a democratic government, later known as the Weimar Republic, was created. The interwar period saw strife between supporters of the new republic and hardline opponents on both the right and left. Italy, as an Entente ally, had made some post-war territorial gains; however, Italian nationalists were angered that the promises made by the United Kingdom and France to secure Italian entrance into the war were not fulfilled in the peace settlement. From 1922 to 1925, the Fascist movement led by Benito Mussolini seized power in Italy with a nationalist, totalitarian, and class collaborationist agenda that abolished representative democracy, repressed socialist, left-wing and liberal forces, and pursued an aggressive expansionist foreign policy aimed at making Italy a world power, promising the creation of a "New Roman Empire".[15]
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Adolf Hitler, after an unsuccessful attempt to overthrow the German government in 1923, eventually became the Chancellor of Germany in 1933. He abolished democracy, espousing a radical, racially motivated revision of the world order, and soon began a massive rearmament campaign.[16] Meanwhile, France, to secure its alliance, allowed Italy a free hand in Ethiopia, which Italy desired as a colonial possession. The situation was aggravated in early 1935 when the Territory of the Saar Basin was legally reunited with Germany and Hitler repudiated the Treaty of Versailles, accelerated his rearmament programme, and introduced conscription.[17]
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The United Kingdom, France and Italy formed the Stresa Front in April 1935 in order to contain Germany, a key step towards military globalisation; however, that June, the United Kingdom made an independent naval agreement with Germany, easing prior restrictions. The Soviet Union, concerned by Germany's goals of capturing vast areas of Eastern Europe, drafted a treaty of mutual assistance with France. Before taking effect, though, the Franco-Soviet pact was required to go through the bureaucracy of the League of Nations, which rendered it essentially toothless.[18] The United States, concerned with events in Europe and Asia, passed the Neutrality Act in August of the same year.[19]
|
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+
Hitler defied the Versailles and Locarno treaties by remilitarising the Rhineland in March 1936, encountering little opposition due to appeasement.[20] In October 1936, Germany and Italy formed the Rome–Berlin Axis. A month later, Germany and Japan signed the Anti-Comintern Pact, which Italy joined the following year.[21]
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+
The Kuomintang (KMT) party in China launched a unification campaign against regional warlords and nominally unified China in the mid-1920s, but was soon embroiled in a civil war against its former Chinese Communist Party allies[22] and new regional warlords. In 1931, an increasingly militaristic Empire of Japan, which had long sought influence in China[23] as the first step of what its government saw as the country's right to rule Asia, staged the Mukden Incident as a pretext to invade Manchuria and establish the puppet state of Manchukuo.[24]
|
44 |
+
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45 |
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China appealed to the League of Nations to stop the Japanese invasion of Manchuria. Japan withdrew from the League of Nations after being condemned for its incursion into Manchuria. The two nations then fought several battles, in Shanghai, Rehe and Hebei, until the Tanggu Truce was signed in 1933. Thereafter, Chinese volunteer forces continued the resistance to Japanese aggression in Manchuria, and Chahar and Suiyuan.[25] After the 1936 Xi'an Incident, the Kuomintang and communist forces agreed on a ceasefire to present a united front to oppose Japan.[26]
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The Second Italo–Ethiopian War was a brief colonial war that began in October 1935 and ended in May 1936. The war began with the invasion of the Ethiopian Empire (also known as Abyssinia) by the armed forces of the Kingdom of Italy (Regno d'Italia), which was launched from Italian Somaliland and Eritrea.[27] The war resulted in the military occupation of Ethiopia and its annexation into the newly created colony of Italian East Africa (Africa Orientale Italiana, or AOI); in addition it exposed the weakness of the League of Nations as a force to preserve peace. Both Italy and Ethiopia were member nations, but the League did little when the former clearly violated Article X of the League's Covenant.[28] The United Kingdom and France supported imposing sanctions on Italy for the invasion, but they were not fully enforced and failed to end the Italian invasion.[29] Italy subsequently dropped its objections to Germany's goal of absorbing Austria.[30]
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When civil war broke out in Spain, Hitler and Mussolini lent military support to the Nationalist rebels, led by General Francisco Franco. Italy supported the Nationalists to a greater extent than the Nazis did: altogether Mussolini sent to Spain more than 70,000 ground troops and 6,000 aviation personnel, as well as about 720 aircraft.[31] The Soviet Union supported the existing government, the Spanish Republic. More than 30,000 foreign volunteers, known as the International Brigades, also fought against the Nationalists. Both Germany and the Soviet Union used this proxy war as an opportunity to test in combat their most advanced weapons and tactics. The Nationalists won the civil war in April 1939; Franco, now dictator, remained officially neutral during World War II but generally favoured the Axis.[32] His greatest collaboration with Germany was the sending of volunteers to fight on the Eastern Front.[33]
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In July 1937, Japan captured the former Chinese imperial capital of Peking after instigating the Marco Polo Bridge Incident, which culminated in the Japanese campaign to invade all of China.[34] The Soviets quickly signed a non-aggression pact with China to lend materiel support, effectively ending China's prior co-operation with Germany. From September to November, the Japanese attacked Taiyuan, engaged the Kuomintang Army around Xinkou,[35] and fought Communist forces in Pingxingguan.[36][37] Generalissimo Chiang Kai-shek deployed his best army to defend Shanghai, but, after three months of fighting, Shanghai fell. The Japanese continued to push the Chinese forces back, capturing the capital Nanking in December 1937. After the fall of Nanking, tens of thousands if not hundreds of thousands of Chinese civilians and disarmed combatants were murdered by the Japanese.[38][39]
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+
In March 1938, Nationalist Chinese forces won their first major victory at Taierzhuang but then the city of Xuzhou was taken by Japanese in May.[40] In June 1938, Chinese forces stalled the Japanese advance by flooding the Yellow River; this manoeuvre bought time for the Chinese to prepare their defences at Wuhan, but the city was taken by October.[41] Japanese military victories did not bring about the collapse of Chinese resistance that Japan had hoped to achieve; instead, the Chinese government relocated inland to Chongqing and continued the war.[42][43]
|
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+
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55 |
+
In the mid-to-late 1930s, Japanese forces in Manchukuo had sporadic border clashes with the Soviet Union and Mongolia. The Japanese doctrine of Hokushin-ron, which emphasised Japan's expansion northward, was favoured by the Imperial Army during this time. With the Japanese defeat at Khalkin Gol in 1939, the ongoing Second Sino-Japanese War[44] and ally Nazi Germany pursuing neutrality with the Soviets, this policy would prove difficult to maintain. Japan and the Soviet Union eventually signed a Neutrality Pact in April 1941, and Japan adopted the doctrine of Nanshin-ron, promoted by the Navy, which took its focus southward, eventually leading to its war with the United States and the Western Allies.[45][46]
|
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+
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+
In Europe, Germany and Italy were becoming more aggressive. In March 1938, Germany annexed Austria, again provoking little response from other European powers.[47] Encouraged, Hitler began pressing German claims on the Sudetenland, an area of Czechoslovakia with a predominantly ethnic German population. Soon the United Kingdom and France followed the appeasement policy of British Prime Minister Neville Chamberlain and conceded this territory to Germany in the Munich Agreement, which was made against the wishes of the Czechoslovak government, in exchange for a promise of no further territorial demands.[48] Soon afterwards, Germany and Italy forced Czechoslovakia to cede additional territory to Hungary, and Poland annexed Czechoslovakia's Zaolzie region.[49]
|
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+
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+
Although all of Germany's stated demands had been satisfied by the agreement, privately Hitler was furious that British interference had prevented him from seizing all of Czechoslovakia in one operation. In subsequent speeches Hitler attacked British and Jewish "war-mongers" and in January 1939 secretly ordered a major build-up of the German navy to challenge British naval supremacy. In March 1939, Germany invaded the remainder of Czechoslovakia and subsequently split it into the German Protectorate of Bohemia and Moravia and a pro-German client state, the Slovak Republic.[50] Hitler also delivered 20 March 1939 ultimatum to Lithuania, forcing the concession of the Klaipėda Region, formerly the German Memelland.[51]
|
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+
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61 |
+
Greatly alarmed and with Hitler making further demands on the Free City of Danzig, the United Kingdom and France guaranteed their support for Polish independence; when Italy conquered Albania in April 1939, the same guarantee was extended to Romania and Greece.[52] Shortly after the Franco-British pledge to Poland, Germany and Italy formalised their own alliance with the Pact of Steel.[53] Hitler accused the United Kingdom and Poland of trying to "encircle" Germany and renounced the Anglo-German Naval Agreement and the German–Polish Non-Aggression Pact.[54]
|
62 |
+
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63 |
+
The situation reached a general crisis in late August as German troops continued to mobilise against the Polish border. On 23 August, when tripartite negotiations about a military alliance between France, the United Kingdom and Soviet Union stalled,[55] the Soviet Union signed a non-aggression pact with Germany.[56] This pact had a secret protocol that defined German and Soviet "spheres of influence" (western Poland and Lithuania for Germany; eastern Poland, Finland, Estonia, Latvia and Bessarabia for the Soviet Union), and raised the question of continuing Polish independence.[57] The pact neutralised the possibility of Soviet opposition to a campaign against Poland and assured that Germany would not have to face the prospect of a two-front war, as it had in World War I. Immediately after that, Hitler ordered the attack to proceed on 26 August, but upon hearing that the United Kingdom had concluded a formal mutual assistance pact with Poland, and that Italy would maintain neutrality, he decided to delay it.[58]
|
64 |
+
|
65 |
+
In response to British requests for direct negotiations to avoid war, Germany made demands on Poland, which only served as a pretext to worsen relations.[59] On 29 August, Hitler demanded that a Polish plenipotentiary immediately travel to Berlin to negotiate the handover of Danzig, and to allow a plebiscite in the Polish Corridor in which the German minority would vote on secession.[59] The Poles refused to comply with the German demands, and on the night of 30–31 August in a stormy meeting with the British ambassador Neville Henderson, Ribbentrop declared that Germany considered its claims rejected.[60]
|
66 |
+
|
67 |
+
On 1 September 1939, Germany invaded Poland after having staged several false flag border incidents as a pretext to initiate the invasion.[61] The first German attack of the war came against the Polish defenses at Westerplatte.[62] The United Kingdom responded with an ultimatum to Germany to cease military operations, and on 3 September, after the ultimatum was ignored, France and Britain declared war on Germany, followed by Australia, New Zealand, South Africa and Canada. The alliance provided no direct military support to Poland, outside of a cautious French probe into the Saarland.[63] The Western Allies also began a naval blockade of Germany, which aimed to damage the country's economy and the war effort.[64] Germany responded by ordering U-boat warfare against Allied merchant and warships, which would later escalate into the Battle of the Atlantic.[65]
|
68 |
+
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69 |
+
On 8 September, German troops reached the suburbs of Warsaw. The Polish counter offensive to the west halted the German advance for several days, but it was outflanked and encircled by the Wehrmacht. Remnants of the Polish army broke through to besieged Warsaw. On 17 September 1939, after signing a cease-fire with Japan, the Soviets invaded Eastern Poland[66] under a pretext that the Polish state had ostensibly ceased to exist.[67] On 27 September, the Warsaw garrison surrendered to the Germans, and the last large operational unit of the Polish Army surrendered on 6 October. Despite the military defeat, Poland never surrendered; instead it formed the Polish government-in-exile and a clandestine state apparatus remained in occupied Poland.[68] A significant part of Polish military personnel evacuated to Romania and the Baltic countries; many of them would fight against the Axis in other theatres of the war.[69]
|
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+
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71 |
+
Germany annexed the western and occupied the central part of Poland, and the Soviet Union annexed its eastern part; small shares of Polish territory were transferred to Lithuania and Slovakia. On 6 October, Hitler made a public peace overture to the United Kingdom and France but said that the future of Poland was to be determined exclusively by Germany and the Soviet Union. The proposal was rejected,[60] and Hitler ordered an immediate offensive against France,[70] which would be postponed until the spring of 1940 due to bad weather.[71][72][73]
|
72 |
+
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73 |
+
The Soviet Union forced the Baltic countries—Estonia, Latvia and Lithuania, the states that were in the Soviet "sphere of influence" under the Molotov-Ribbentrop pact—to sign "mutual assistance pacts" that stipulated stationing Soviet troops in these countries. Soon after, significant Soviet military contingents were moved there.[74][75][76] Finland refused to sign a similar pact and rejected ceding part of its territory to the Soviet Union. The Soviet Union invaded Finland in November 1939,[77] and the Soviet Union was expelled from the League of Nations.[78] Despite overwhelming numerical superiority, Soviet military success was modest, and the Finno-Soviet war ended in March 1940 with minimal Finnish concessions.[79]
|
74 |
+
|
75 |
+
In June 1940, the Soviet Union forcibly annexed Estonia, Latvia and Lithuania,[75] and the disputed Romanian regions of Bessarabia, northern Bukovina and Hertza. Meanwhile, Nazi-Soviet political rapprochement and economic co-operation[80][81] gradually stalled,[82][83] and both states began preparations for war.[84]
|
76 |
+
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77 |
+
In April 1940, Germany invaded Denmark and Norway to protect shipments of iron ore from Sweden, which the Allies were attempting to cut off.[85] Denmark capitulated after a few hours, and Norway was conquered within two months[86] despite Allied support. British discontent over the Norwegian campaign led to the appointment of Winston Churchill as Prime Minister on 10 May 1940.[87]
|
78 |
+
|
79 |
+
On the same day, Germany launched an offensive against France. To circumvent the strong Maginot Line fortifications on the Franco-German border, Germany directed its attack at the neutral nations of Belgium, the Netherlands, and Luxembourg.[88] The Germans carried out a flanking manoeuvre through the Ardennes region,[89] which was mistakenly perceived by Allies as an impenetrable natural barrier against armoured vehicles.[90][91] By successfully implementing new blitzkrieg tactics, the Wehrmacht rapidly advanced to the Channel and cut off the Allied forces in Belgium, trapping the bulk of the Allied armies in a cauldron on the Franco-Belgian border near Lille. The United Kingdom was able to evacuate a significant number of Allied troops from the continent by early June, although abandoning almost all their equipment.[92]
|
80 |
+
|
81 |
+
On 10 June, Italy invaded France, declaring war on both France and the United Kingdom.[93] The Germans turned south against the weakened French army, and Paris fell to them on 14 June. Eight days later France signed an armistice with Germany; it was divided into German and Italian occupation zones,[94] and an unoccupied rump state under the Vichy Regime, which, though officially neutral, was generally aligned with Germany. France kept its fleet, which the United Kingdom attacked on 3 July in an attempt to prevent its seizure by Germany.[95]
|
82 |
+
|
83 |
+
The Battle of Britain[96] began in early July with Luftwaffe attacks on shipping and harbours.[97] The United Kingdom rejected Hitler's ultimatum,[which?][98] and the German air superiority campaign started in August but failed to defeat RAF Fighter Command, forcing the indefinite postponement of the proposed German invasion of Britain. The German strategic bombing offensive intensified with night attacks on London and other cities in the Blitz, but failed to significantly disrupt the British war effort[97] and largely ended in May 1941.[99]
|
84 |
+
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85 |
+
Using newly captured French ports, the German Navy enjoyed success against an over-extended Royal Navy, using U-boats against British shipping in the Atlantic.[100] The British Home Fleet scored a significant victory on 27 May 1941 by sinking the German battleship Bismarck.[101]
|
86 |
+
|
87 |
+
In November 1939, the United States was taking measures to assist China and the Western Allies, and amended the Neutrality Act to allow "cash and carry" purchases by the Allies.[102] In 1940, following the German capture of Paris, the size of the United States Navy was significantly increased. In September the United States further agreed to a trade of American destroyers for British bases.[103] Still, a large majority of the American public continued to oppose any direct military intervention in the conflict well into 1941.[104] In December 1940 Roosevelt accused Hitler of planning world conquest and ruled out any negotiations as useless, calling for the United States to become an "arsenal of democracy" and promoting Lend-Lease programmes of aid to support the British war effort.[98] The United States started strategic planning to prepare for a full-scale offensive against Germany.[105]
|
88 |
+
|
89 |
+
At the end of September 1940, the Tripartite Pact formally united Japan, Italy, and Germany as the Axis Powers. The Tripartite Pact stipulated that any country, with the exception of the Soviet Union, which attacked any Axis Power would be forced to go to war against all three.[106] The Axis expanded in November 1940 when Hungary, Slovakia and Romania joined.[107] Romania and Hungary would make major contributions to the Axis war against the Soviet Union, in Romania's case partially to recapture territory ceded to the Soviet Union.[108]
|
90 |
+
|
91 |
+
In early June 1940 the Italian Regia Aeronautica attacked and besieged Malta, a British possession. In late summer through early autumn Italy conquered British Somaliland and made an incursion into British-held Egypt. In October Italy attacked Greece, but the attack was repulsed with heavy Italian casualties; the campaign ended within months with minor territorial changes.[109] Germany started preparation for an invasion of the Balkans to assist Italy, to prevent the British from gaining a foothold there, which would be a potential threat for Romanian oil fields, and to strike against the British dominance of the Mediterranean.[110]
|
92 |
+
|
93 |
+
In December 1940, British Empire forces began counter-offensives against Italian forces in Egypt and Italian East Africa.[111] The offensives were highly successful; by early February 1941 Italy had lost control of eastern Libya, and large numbers of Italian troops had been taken, prisoner. The Italian Navy also suffered significant defeats, with the Royal Navy putting three Italian battleships out of commission by a carrier attack at Taranto and neutralising several more warships at the Battle of Cape Matapan.[112]
|
94 |
+
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95 |
+
Italian defeats prompted Germany to deploy an expeditionary force to North Africa, and at the end of March 1941 Rommel's Afrika Korps launched an offensive which drove back the Commonwealth forces.[113] In under a month, Axis forces advanced to western Egypt and besieged the port of Tobruk.[114]
|
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+
|
97 |
+
By late March 1941 Bulgaria and Yugoslavia signed the Tripartite Pact; however, the Yugoslav government was overthrown two days later by pro-British nationalists. Germany responded with simultaneous invasions of both Yugoslavia and Greece, commencing on 6 April 1941; both nations were forced to surrender within the month.[115] The airborne invasion of the Greek island of Crete at the end of May completed the German conquest of the Balkans.[116] Although the Axis victory was swift, bitter and large-scale partisan warfare subsequently broke out against the Axis occupation of Yugoslavia, which continued until the end of the war.[117]
|
98 |
+
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99 |
+
In the Middle East, in May Commonwealth forces quashed an uprising in Iraq which had been supported by German aircraft from bases within Vichy-controlled Syria.[118] Between June and July they invaded and occupied the French possessions Syria and Lebanon, with the assistance of the Free French.[119]
|
100 |
+
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101 |
+
With the situation in Europe and Asia relatively stable, Germany, Japan, and the Soviet Union made preparations. With the Soviets wary of mounting tensions with Germany and the Japanese planning to take advantage of the European War by seizing resource-rich European possessions in Southeast Asia, the two powers signed the Soviet–Japanese Neutrality Pact in April 1941.[120] By contrast, the Germans were steadily making preparations for an attack on the Soviet Union, massing forces on the Soviet border.[121]
|
102 |
+
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103 |
+
Hitler believed that the United Kingdom's refusal to end the war was based on the hope that the United States and the Soviet Union would enter the war against Germany sooner or later.[122] He, therefore, decided to try to strengthen Germany's relations with the Soviets, or failing that to attack and eliminate them as a factor. In November 1940, negotiations took place to determine if the Soviet Union would join the Tripartite Pact. The Soviets showed some interest but asked for concessions from Finland, Bulgaria, Turkey, and Japan that Germany considered unacceptable. On 18 December 1940, Hitler issued the directive to prepare for an invasion of the Soviet Union.[123]
|
104 |
+
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105 |
+
On 22 June 1941, Germany, supported by Italy and Romania, invaded the Soviet Union in Operation Barbarossa, with Germany accusing the Soviets of plotting against them. They were joined shortly by Finland and Hungary.[124] The primary targets of this surprise offensive[125] were the Baltic region, Moscow and Ukraine, with the ultimate goal of ending the 1941 campaign near the Arkhangelsk-Astrakhan line, from the Caspian to the White Seas. Hitler's objectives were to eliminate the Soviet Union as a military power, exterminate Communism, generate Lebensraum ("living space")[126] by dispossessing the native population[127] and guarantee access to the strategic resources needed to defeat Germany's remaining rivals.[128]
|
106 |
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107 |
+
Although the Red Army was preparing for strategic counter-offensives before the war,[129] Barbarossa forced the Soviet supreme command to adopt a strategic defence. During the summer, the Axis made significant gains into Soviet territory, inflicting immense losses in both personnel and materiel. By mid-August, however, the German Army High Command decided to suspend the offensive of a considerably depleted Army Group Centre, and to divert the 2nd Panzer Group to reinforce troops advancing towards central Ukraine and Leningrad.[130] The Kiev offensive was overwhelmingly successful, resulting in encirclement and elimination of four Soviet armies, and made possible further advance into Crimea and industrially developed Eastern Ukraine (the First Battle of Kharkov).[131]
|
108 |
+
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109 |
+
The diversion of three quarters of the Axis troops and the majority of their air forces from France and the central Mediterranean to the Eastern Front[132] prompted the United Kingdom to reconsider its grand strategy.[133] In July, the UK and the Soviet Union formed a military alliance against Germany[134] and in August, the United Kingdom and the United States jointly issued the Atlantic Charter, which outlined British and American goals for the postwar world.[135] In late August the British and Soviets invaded neutral Iran to secure the Persian Corridor, Iran's oil fields, and preempt any Axis advances through Iran toward the Baku oil fields or British India.[136]
|
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+
By October Axis operational objectives in Ukraine and the Baltic region were achieved, with only the sieges of Leningrad[137] and Sevastopol continuing.[138] A major offensive against Moscow was renewed; after two months of fierce battles in increasingly harsh weather, the German army almost reached the outer suburbs of Moscow, where the exhausted troops[139] were forced to suspend their offensive.[140] Large territorial gains were made by Axis forces, but their campaign had failed to achieve its main objectives: two key cities remained in Soviet hands, the Soviet capability to resist was not broken, and the Soviet Union retained a considerable part of its military potential. The blitzkrieg phase of the war in Europe had ended.[141]
|
112 |
+
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113 |
+
By early December, freshly mobilised reserves[142] allowed the Soviets to achieve numerical parity with Axis troops.[143] This, as well as intelligence data which established that a minimal number of Soviet troops in the East would be sufficient to deter any attack by the Japanese Kwantung Army,[144] allowed the Soviets to begin a massive counter-offensive that started on 5 December all along the front and pushed German troops 100–250 kilometres (62–155 mi) west.[145]
|
114 |
+
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115 |
+
Following the Japanese false flag Mukden Incident in 1931, the Japanese shelling of the American gunboat USS Panay in 1937, and the 1937-38 Nanjing Massacre Japanese-American relations deteriorated. In 1939, the United States notified Japan that it would not be extending its trade treaty and American public opinion opposing Japanese expansionism led to a series of economic sanctions, the Export Control Acts, which banned U.S. exports of chemicals, minerals and military parts to Japan and increased economic pressure on the Japanese regime.[98][146][147] During 1939 Japan launched its first attack against Changsha, a strategically important Chinese city, but was repulsed by late September.[148] Despite several offensives by both sides, the war between China and Japan was stalemated by 1940. To increase pressure on China by blocking supply routes, and to better position Japanese forces in the event of a war with the Western powers, Japan invaded and occupied northern Indochina in September, 1940.[149]
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Chinese nationalist forces launched a large-scale counter-offensive in early 1940. In August, Chinese communists launched an offensive in Central China; in retaliation, Japan instituted harsh measures in occupied areas to reduce human and material resources for the communists.[150] Continued antipathy between Chinese communist and nationalist forces culminated in armed clashes in January 1941, effectively ending their co-operation.[151] In March, the Japanese 11th army attacked the headquarters of the Chinese 19th army but was repulsed during Battle of Shanggao.[152] In September, Japan attempted to take the city of Changsha again and clashed with Chinese nationalist forces.[153]
|
118 |
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119 |
+
German successes in Europe encouraged Japan to increase pressure on European governments in Southeast Asia. The Dutch government agreed to provide Japan some oil supplies from the Dutch East Indies, but negotiations for additional access to their resources ended in failure in June 1941.[154] In July 1941 Japan sent troops to southern Indochina, thus threatening British and Dutch possessions in the Far East. The United States, United Kingdom, and other Western governments reacted to this move with a freeze on Japanese assets and a total oil embargo.[155][156] At the same time, Japan was planning an invasion of the Soviet Far East, intending to capitalise off the German invasion in the west, but abandoned the operation after the sanctions.[157]
|
120 |
+
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121 |
+
Since early 1941 the United States and Japan had been engaged in negotiations in an attempt to improve their strained relations and end the war in China. During these negotiations, Japan advanced a number of proposals which were dismissed by the Americans as inadequate.[158] At the same time the United States, the United Kingdom, and the Netherlands engaged in secret discussions for the joint defence of their territories, in the event of a Japanese attack against any of them.[159] Roosevelt reinforced the Philippines (an American protectorate scheduled for independence in 1946) and warned Japan that the United States would react to Japanese attacks against any "neighboring countries".[159]
|
122 |
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123 |
+
Frustrated at the lack of progress and feeling the pinch of the American–British–Dutch sanctions, Japan prepared for war. On 20 November, a new government under Hideki Tojo presented an interim proposal as its final offer. It called for the end of American aid to China and for lifting the embargo on the supply of oil and other resources to Japan. In exchange, Japan promised not to launch any attacks in Southeast Asia and to withdraw its forces from southern Indochina.[158] The American counter-proposal of 26 November required that Japan evacuate all of China without conditions and conclude non-aggression pacts with all Pacific powers.[160] That meant Japan was essentially forced to choose between abandoning its ambitions in China, or seizing the natural resources it needed in the Dutch East Indies by force;[161][162] the Japanese military did not consider the former an option, and many officers considered the oil embargo an unspoken declaration of war.[163]
|
124 |
+
|
125 |
+
Japan planned to rapidly seize European colonies in Asia to create a large defensive perimeter stretching into the Central Pacific. The Japanese would then be free to exploit the resources of Southeast Asia while exhausting the over-stretched Allies by fighting a defensive war.[164][165] To prevent American intervention while securing the perimeter, it was further planned to neutralise the United States Pacific Fleet and the American military presence in the Philippines from the outset.[166] On 7 December 1941 (8 December in Asian time zones), Japan attacked British and American holdings with near-simultaneous offensives against Southeast Asia and the Central Pacific.[167] These included an attack on the American fleets at Pearl Harbor and the Philippines, landings in Malaya,[167], Thailand and the Battle of Hong Kong.[168]
|
126 |
+
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127 |
+
The Japanese invasion of Thailand led to Thailand's decision to ally itself with Japan and the other Japanese attacks led the United States, United Kingdom, China, Australia, and several other states to formally declare war on Japan, whereas the Soviet Union, being heavily involved in large-scale hostilities with European Axis countries, maintained its neutrality agreement with Japan.[169] Germany, followed by the other Axis states, declared war on the United States[170] in solidarity with Japan, citing as justification the American attacks on German war vessels that had been ordered by Roosevelt.[124][171]
|
128 |
+
|
129 |
+
On 1 January 1942, the Allied Big Four[172]—the Soviet Union, China, the United Kingdom and the United States—and 22 smaller or exiled governments issued the Declaration by United Nations, thereby affirming the Atlantic Charter,[173] and agreeing not to sign a separate peace with the Axis powers.[174]
|
130 |
+
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131 |
+
During 1942, Allied officials debated on the appropriate grand strategy to pursue. All agreed that defeating Germany was the primary objective. The Americans favoured a straightforward, large-scale attack on Germany through France. The Soviets were also demanding a second front. The British, on the other hand, argued that military operations should target peripheral areas to wear out German strength, leading to increasing demoralisation, and bolster resistance forces. Germany itself would be subject to a heavy bombing campaign. An offensive against Germany would then be launched primarily by Allied armour without using large-scale armies.[175] Eventually, the British persuaded the Americans that a landing in France was infeasible in 1942 and they should instead focus on driving the Axis out of North Africa.[176]
|
132 |
+
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133 |
+
At the Casablanca Conference in early 1943, the Allies reiterated the statements issued in the 1942 Declaration, and demanded the unconditional surrender of their enemies. The British and Americans agreed to continue to press the initiative in the Mediterranean by invading Sicily to fully secure the Mediterranean supply routes.[177] Although the British argued for further operations in the Balkans to bring Turkey into the war, in May 1943, the Americans extracted a British commitment to limit Allied operations in the Mediterranean to an invasion of the Italian mainland and to invade France in 1944.[178]
|
134 |
+
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135 |
+
By the end of April 1942, Japan and its ally Thailand had almost fully conquered Burma, Malaya, the Dutch East Indies, Singapore, and Rabaul, inflicting severe losses on Allied troops and taking a large number of prisoners.[179] Despite stubborn resistance by Filipino and US forces, the Philippine Commonwealth was eventually captured in May 1942, forcing its government into exile.[180] On 16 April, in Burma, 7,000 British soldiers were encircled by the Japanese 33rd Division during the Battle of Yenangyaung and rescued by the Chinese 38th Division.[181] Japanese forces also achieved naval victories in the South China Sea, Java Sea and Indian Ocean,[182] and bombed the Allied naval base at Darwin, Australia. In January 1942, the only Allied success against Japan was a Chinese victory at Changsha.[183] These easy victories over unprepared US and European opponents left Japan overconfident, as well as overextended.[184]
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136 |
+
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In early May 1942, Japan initiated operations to capture Port Moresby by amphibious assault and thus sever communications and supply lines between the United States and Australia. The planned invasion was thwarted when an Allied task force, centred on two American fleet carriers, fought Japanese naval forces to a draw in the Battle of the Coral Sea.[185] Japan's next plan, motivated by the earlier Doolittle Raid, was to seize Midway Atoll and lure American carriers into battle to be eliminated; as a diversion, Japan would also send forces to occupy the Aleutian Islands in Alaska.[186] In mid-May, Japan started the Zhejiang-Jiangxi Campaign in China, with the goal of inflicting retribution on the Chinese who aided the surviving American airmen in the Doolittle Raid by destroying air bases and fighting against the Chinese 23rd and 32nd Army Groups.[187][188] In early June, Japan put its operations into action, but the Americans, having broken Japanese naval codes in late May, were fully aware of the plans and order of battle, and used this knowledge to achieve a decisive victory at Midway over the Imperial Japanese Navy.[189]
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With its capacity for aggressive action greatly diminished as a result of the Midway battle, Japan chose to focus on a belated attempt to capture Port Moresby by an overland campaign in the Territory of Papua.[190] The Americans planned a counter-attack against Japanese positions in the southern Solomon Islands, primarily Guadalcanal, as a first step towards capturing Rabaul, the main Japanese base in Southeast Asia.[191]
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Both plans started in July, but by mid-September, the Battle for Guadalcanal took priority for the Japanese, and troops in New Guinea were ordered to withdraw from the Port Moresby area to the northern part of the island, where they faced Australian and United States troops in the Battle of Buna-Gona.[192] Guadalcanal soon became a focal point for both sides with heavy commitments of troops and ships in the battle for Guadalcanal. By the start of 1943, the Japanese were defeated on the island and withdrew their troops.[193] In Burma, Commonwealth forces mounted two operations. The first, an offensive into the Arakan region in late 1942, went disastrously, forcing a retreat back to India by May 1943.[194] The second was the insertion of irregular forces behind Japanese front-lines in February which, by the end of April, had achieved mixed results.[195]
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Despite considerable losses, in early 1942 Germany and its allies stopped a major Soviet offensive in central and southern Russia, keeping most territorial gains they had achieved during the previous year.[196] In May the Germans defeated Soviet offensives in the Kerch Peninsula and at Kharkov,[197] and then launched their main summer offensive against southern Russia in June 1942, to seize the oil fields of the Caucasus and occupy Kuban steppe, while maintaining positions on the northern and central areas of the front. The Germans split Army Group South into two groups: Army Group A advanced to the lower Don River and struck south-east to the Caucasus, while Army Group B headed towards the Volga River. The Soviets decided to make their stand at Stalingrad on the Volga.[198]
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By mid-November, the Germans had nearly taken Stalingrad in bitter street fighting. The Soviets began their second winter counter-offensive, starting with an encirclement of German forces at Stalingrad,[199] and an assault on the Rzhev salient near Moscow, though the latter failed disastrously.[200] By early February 1943, the German Army had taken tremendous losses; German troops at Stalingrad had been defeated,[201] and the front-line had been pushed back beyond its position before the summer offensive. In mid-February, after the Soviet push had tapered off, the Germans launched another attack on Kharkov, creating a salient in their front line around the Soviet city of Kursk.[202]
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Exploiting poor American naval command decisions, the German navy ravaged Allied shipping off the American Atlantic coast.[203] By November 1941, Commonwealth forces had launched a counter-offensive, Operation Crusader, in North Africa, and reclaimed all the gains the Germans and Italians had made.[204] In North Africa, the Germans launched an offensive in January, pushing the British back to positions at the Gazala Line by early February,[205] followed by a temporary lull in combat which Germany used to prepare for their upcoming offensives.[206] Concerns the Japanese might use bases in Vichy-held Madagascar caused the British to invade the island in early May 1942.[207] An Axis offensive in Libya forced an Allied retreat deep inside Egypt until Axis forces were stopped at El Alamein.[208] On the Continent, raids of Allied commandos on strategic targets, culminating in the disastrous Dieppe Raid,[209] demonstrated the Western Allies' inability to launch an invasion of continental Europe without much better preparation, equipment, and operational security.[210][page needed]
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In August 1942, the Allies succeeded in repelling a second attack against El Alamein[211] and, at a high cost, managed to deliver desperately needed supplies to the besieged Malta.[212] A few months later, the Allies commenced an attack of their own in Egypt, dislodging the Axis forces and beginning a drive west across Libya.[213] This attack was followed up shortly after by Anglo-American landings in French North Africa, which resulted in the region joining the Allies.[214] Hitler responded to the French colony's defection by ordering the occupation of Vichy France;[214] although Vichy forces did not resist this violation of the armistice, they managed to scuttle their fleet to prevent its capture by German forces.[214][215] The Axis forces in Africa withdrew into Tunisia, which was conquered by the Allies in May 1943.[214][216]
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In June 1943 the British and Americans began a strategic bombing campaign against Germany with a goal to disrupt the war economy, reduce morale, and "de-house" the civilian population.[217] The firebombing of Hamburg was among the first attacks in this campaign, inflicting significant casualties and considerable losses on infrastructure of this important industrial centre.[218]
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After the Guadalcanal Campaign, the Allies initiated several operations against Japan in the Pacific. In May 1943, Canadian and US forces were sent to eliminate Japanese forces from the Aleutians.[219] Soon after, the United States, with support from Australia, New Zealand and Pacific Islander forces, began major ground, sea and air operations to isolate Rabaul by capturing surrounding islands, and breach the Japanese Central Pacific perimeter at the Gilbert and Marshall Islands.[220] By the end of March 1944, the Allies had completed both of these objectives and had also neutralised the major Japanese base at Truk in the Caroline Islands. In April, the Allies launched an operation to retake Western New Guinea.[221]
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In the Soviet Union, both the Germans and the Soviets spent the spring and early summer of 1943 preparing for large offensives in central Russia. On 4 July 1943, Germany attacked Soviet forces around the Kursk Bulge. Within a week, German forces had exhausted themselves against the Soviets' deeply echeloned and well-constructed defences,[222] and for the first time in the war Hitler cancelled the operation before it had achieved tactical or operational success.[223] This decision was partially affected by the Western Allies' invasion of Sicily launched on 9 July, which, combined with previous Italian failures, resulted in the ousting and arrest of Mussolini later that month.[224]
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On 12 July 1943, the Soviets launched their own counter-offensives, thereby dispelling any chance of German victory or even stalemate in the east. The Soviet victory at Kursk marked the end of German superiority,[225] giving the Soviet Union the initiative on the Eastern Front.[226][227] The Germans tried to stabilise their eastern front along the hastily fortified Panther–Wotan line, but the Soviets broke through it at Smolensk and by the Lower Dnieper Offensives.[228]
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On 3 September 1943, the Western Allies invaded the Italian mainland, following Italy's armistice with the Allies.[229] Germany with the help of fascists responded by disarming Italian forces that were in many places without superior orders, seizing military control of Italian areas,[230] and creating a series of defensive lines.[231] German special forces then rescued Mussolini, who then soon established a new client state in German-occupied Italy named the Italian Social Republic,[232] causing an Italian civil war. The Western Allies fought through several lines until reaching the main German defensive line in mid-November.[233]
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German operations in the Atlantic also suffered. By May 1943, as Allied counter-measures became increasingly effective, the resulting sizeable German submarine losses forced a temporary halt of the German Atlantic naval campaign.[234] In November 1943, Franklin D. Roosevelt and Winston Churchill met with Chiang Kai-shek in Cairo and then with Joseph Stalin in Tehran.[235] The former conference determined the post-war return of Japanese territory[236] and the military planning for the Burma Campaign,[237] while the latter included agreement that the Western Allies would invade Europe in 1944 and that the Soviet Union would declare war on Japan within three months of Germany's defeat.[238]
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From November 1943, during the seven-week Battle of Changde, the Chinese forced Japan to fight a costly war of attrition, while awaiting Allied relief.[239][240][241] In January 1944, the Allies launched a series of attacks in Italy against the line at Monte Cassino and tried to outflank it with landings at Anzio.[242]
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On 27 January 1944, Soviet troops launched a major offensive that expelled German forces from the Leningrad region, thereby ending the most lethal siege in history.[243] The following Soviet offensive was halted on the pre-war Estonian border by the German Army Group North aided by Estonians hoping to re-establish national independence. This delay slowed subsequent Soviet operations in the Baltic Sea region.[244] By late May 1944, the Soviets had liberated Crimea, largely expelled Axis forces from Ukraine, and made incursions into Romania, which were repulsed by the Axis troops.[245] The Allied offensives in Italy had succeeded and, at the expense of allowing several German divisions to retreat, on 4 June Rome was captured.[246]
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The Allies had mixed success in mainland Asia. In March 1944, the Japanese launched the first of two invasions, an operation against British positions in Assam, India,[247] and soon besieged Commonwealth positions at Imphal and Kohima.[248] In May 1944, British forces mounted a counter-offensive that drove Japanese troops back to Burma by July,[248] and Chinese forces that had invaded northern Burma in late 1943 besieged Japanese troops in Myitkyina.[249] The second Japanese invasion of China aimed to destroy China's main fighting forces, secure railways between Japanese-held territory and capture Allied airfields.[250] By June, the Japanese had conquered the province of Henan and begun a new attack on Changsha in Hunan province.[251]
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On 6 June 1944 (known as D-Day), after three years of Soviet pressure,[252] the Western Allies invaded northern France. After reassigning several Allied divisions from Italy, they also attacked southern France.[253] These landings were successful, and led to the defeat of the German Army units in France. Paris was liberated on 25 August by the local resistance assisted by the Free French Forces, both led by General Charles de Gaulle,[254] and the Western Allies continued to push back German forces in western Europe during the latter part of the year. An attempt to advance into northern Germany spearheaded by a major airborne operation in the Netherlands failed.[255] After that, the Western Allies slowly pushed into Germany, but failed to cross the Rur river in a large offensive. In Italy, Allied advance also slowed due to the last major German defensive line.[256]
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On 22 June, the Soviets launched a strategic offensive in Belarus ("Operation Bagration") that destroyed the German Army Group Centre almost completely.[257] Soon after that, another Soviet strategic offensive forced German troops from Western Ukraine and Eastern Poland. The Soviets formed the Polish Committee of National Liberation to control territory in Poland and combat the Polish Armia Krajowa; The Soviet Red Army remained in the Praga district on the other side of the Vistula and watched passively as the Germans quelled the Warsaw Uprising initiated by the Armia Krajowa.[258] The national uprising in Slovakia was also quelled by the Germans.[259] The Soviet Red Army's strategic offensive in eastern Romania cut off and destroyed the considerable German troops there and triggered a successful coup d'état in Romania and in Bulgaria, followed by those countries' shift to the Allied side.[260]
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In September 1944, Soviet troops advanced into Yugoslavia and forced the rapid withdrawal of German Army Groups E and F in Greece, Albania and Yugoslavia to rescue them from being cut off.[261] By this point, the Communist-led Partisans under Marshal Josip Broz Tito, who had led an increasingly successful guerrilla campaign against the occupation since 1941, controlled much of the territory of Yugoslavia and engaged in delaying efforts against German forces further south. In northern Serbia, the Soviet Red Army, with limited support from Bulgarian forces, assisted the Partisans in a joint liberation of the capital city of Belgrade on 20 October. A few days later, the Soviets launched a massive assault against German-occupied Hungary that lasted until the fall of Budapest in February 1945.[262] Unlike impressive Soviet victories in the Balkans, bitter Finnish resistance to the Soviet offensive in the Karelian Isthmus denied the Soviets occupation of Finland and led to a Soviet-Finnish armistice on relatively mild conditions,[263] although Finland was forced to fight their former ally Germany.[264][broken footnote]
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By the start of July 1944, Commonwealth forces in Southeast Asia had repelled the Japanese sieges in Assam, pushing the Japanese back to the Chindwin River[265] while the Chinese captured Myitkyina. In September 1944, Chinese forces captured Mount Song and reopened the Burma Road.[266] In China, the Japanese had more successes, having finally captured Changsha in mid-June and the city of Hengyang by early August.[267] Soon after, they invaded the province of Guangxi, winning major engagements against Chinese forces at Guilin and Liuzhou by the end of November[268] and successfully linking up their forces in China and Indochina by mid-December.[269]
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In the Pacific, US forces continued to press back the Japanese perimeter. In mid-June 1944, they began their offensive against the Mariana and Palau islands, and decisively defeated Japanese forces in the Battle of the Philippine Sea. These defeats led to the resignation of the Japanese Prime Minister, Hideki Tojo, and provided the United States with air bases to launch intensive heavy bomber attacks on the Japanese home islands. In late October, American forces invaded the Filipino island of Leyte; soon after, Allied naval forces scored another large victory in the Battle of Leyte Gulf, one of the largest naval battles in history.[270]
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On 16 December 1944, Germany made a last attempt on the Western Front by using most of its remaining reserves to launch a massive counter-offensive in the Ardennes and along with the French-German border to split the Western Allies, encircle large portions of Western Allied troops and capture their primary supply port at Antwerp to prompt a political settlement.[271] By January, the offensive had been repulsed with no strategic objectives fulfilled.[271] In Italy, the Western Allies remained stalemated at the German defensive line. In mid-January 1945, the Soviets and Poles attacked in Poland, pushing from the Vistula to the Oder river in Germany, and overran East Prussia.[272] On 4 February Soviet, British, and US leaders met for the Yalta Conference. They agreed on the occupation of post-war Germany, and on when the Soviet Union would join the war against Japan.[273]
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In February, the Soviets entered Silesia and Pomerania, while Western Allies entered western Germany and closed to the Rhine river. By March, the Western Allies crossed the Rhine north and south of the Ruhr, encircling the German Army Group B.[274] In early March, in an attempt to protect its last oil reserves in Hungary and to retake Budapest, Germany launched its last major offensive against Soviet troops near Lake Balaton. In two weeks, the offensive had been repulsed, the Soviets advanced to Vienna, and captured the city. In early April, Soviet troops captured Königsberg, while the Western Allies finally pushed forward in Italy and swept across western Germany capturing Hamburg and Nuremberg. American and Soviet forces met at the Elbe river on 25 April, leaving several unoccupied pockets in southern Germany and around Berlin.
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Soviet and Polish forces stormed and captured Berlin in late April. In Italy, German forces surrendered on 29 April. On 30 April, the Reichstag was captured, signalling the military defeat of Nazi Germany,[275] Berlin garrison surrendered on 2 May.
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Several changes in leadership occurred during this period. On 12 April, President Roosevelt died and was succeeded by Harry S. Truman. Benito Mussolini was killed by Italian partisans on 28 April.[276] Two days later, Hitler committed suicide in besieged Berlin, and he was succeeded by Grand Admiral Karl Dönitz.[277]
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Total and unconditional surrender in Europe was signed on 7 and 8 May, to be effective by the end of 8 May.[278] German Army Group Centre resisted in Prague until 11 May.[279]
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In the Pacific theatre, American forces accompanied by the forces of the Philippine Commonwealth advanced in the Philippines, clearing Leyte by the end of April 1945. They landed on Luzon in January 1945 and recaptured Manila in March. Fighting continued on Luzon, Mindanao, and other islands of the Philippines until the end of the war.[280] Meanwhile, the United States Army Air Forces launched a massive firebombing campaign of strategic cities in Japan in an effort to destroy Japanese war industry and civilian morale. A devastating bombing raid on Tokyo of 9–10 March was the deadliest conventional bombing raid in history.[281]
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In May 1945, Australian troops landed in Borneo, over-running the oilfields there. British, American, and Chinese forces defeated the Japanese in northern Burma in March, and the British pushed on to reach Rangoon by 3 May.[282] Chinese forces started a counterattack in the Battle of West Hunan that occurred between 6 April and 7 June 1945. American naval and amphibious forces also moved towards Japan, taking Iwo Jima by March, and Okinawa by the end of June.[283] At the same time, American submarines cut off Japanese imports, drastically reducing Japan's ability to supply its overseas forces.[284]
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On 11 July, Allied leaders met in Potsdam, Germany. They confirmed earlier agreements about Germany,[285] and the American, British and Chinese governments reiterated the demand for unconditional surrender of Japan, specifically stating that "the alternative for Japan is prompt and utter destruction".[286] During this conference, the United Kingdom held its general election, and Clement Attlee replaced Churchill as Prime Minister.[287]
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The call for unconditional surrender was rejected by the Japanese government, which believed it would be capable of negotiating for more favourable surrender terms.[288] In early August, the United States dropped atomic bombs on the Japanese cities of Hiroshima and Nagasaki. Between the two bombings, the Soviets, pursuant to the Yalta agreement, invaded Japanese-held Manchuria and quickly defeated the Kwantung Army, which was the largest Japanese fighting force.[289] These two events persuaded previously adamant Imperial Army leaders to accept surrender terms.[290] The Red Army also captured the southern part of Sakhalin Island and the Kuril Islands. On 15 August 1945, Japan surrendered, with the surrender documents finally signed at Tokyo Bay on the deck of the American battleship USS Missouri on 2 September 1945, ending the war.[291]
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The Allies established occupation administrations in Austria and Germany. The former became a neutral state, non-aligned with any political bloc. The latter was divided into western and eastern occupation zones controlled by the Western Allies and the Soviet Union. A denazification programme in Germany led to the prosecution of Nazi war criminals in the Nuremberg trials and the removal of ex-Nazis from power, although this policy moved towards amnesty and re-integration of ex-Nazis into West German society.[292]
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Germany lost a quarter of its pre-war (1937) territory. Among the eastern territories, Silesia, Neumark and most of Pomerania were taken over by Poland,[293] and East Prussia was divided between Poland and the Soviet Union, followed by the expulsion to Germany of the nine million Germans from these provinces,[294][295] as well as three million Germans from the Sudetenland in Czechoslovakia. By the 1950s, one-fifth of West Germans were refugees from the east. The Soviet Union also took over the Polish provinces east of the Curzon line,[296] from which 2 million Poles were expelled;[295][297] north-east Romania,[298][299] parts of eastern Finland,[300] and the three Baltic states were incorporated into the Soviet Union.[301][302]
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In an effort to maintain world peace,[303] the Allies formed the United Nations, which officially came into existence on 24 October 1945,[304] and adopted the Universal Declaration of Human Rights in 1948 as a common standard for all member nations.[305] The great powers that were the victors of the war—France, China, the United Kingdom, the Soviet Union and the United States—became the permanent members of the UN's Security Council.[306] The five permanent members remain so to the present, although there have been two seat changes, between the Republic of China and the People's Republic of China in 1971, and between the Soviet Union and its successor state, the Russian Federation, following the dissolution of the Soviet Union in 1991. The alliance between the Western Allies and the Soviet Union had begun to deteriorate even before the war was over.[307]
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Germany had been de facto divided, and two independent states, the Federal Republic of Germany (West Germany) and the German Democratic Republic (East Germany),[308] were created within the borders of Allied and Soviet occupation zones. The rest of Europe was also divided into Western and Soviet spheres of influence.[309] Most eastern and central European countries fell into the Soviet sphere, which led to establishment of Communist-led regimes, with full or partial support of the Soviet occupation authorities. As a result, East Germany,[310] Poland, Hungary, Romania, Czechoslovakia, and Albania[311] became Soviet satellite states. Communist Yugoslavia conducted a fully independent policy, causing tension with the Soviet Union.[312]
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Post-war division of the world was formalised by two international military alliances, the United States-led NATO and the Soviet-led Warsaw Pact.[313] The long period of political tensions and military competition between them, the Cold War, would be accompanied by an unprecedented arms race and proxy wars.[314]
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In Asia, the United States led the occupation of Japan and administrated Japan's former islands in the Western Pacific, while the Soviets annexed South Sakhalin and the Kuril Islands.[315] Korea, formerly under Japanese rule, was divided and occupied by the Soviet Union in the North and the United States in the South between 1945 and 1948. Separate republics emerged on both sides of the 38th parallel in 1948, each claiming to be the legitimate government for all of Korea, which led ultimately to the Korean War.[316]
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In China, nationalist and communist forces resumed the civil war in June 1946. Communist forces were victorious and established the People's Republic of China on the mainland, while nationalist forces retreated to Taiwan in 1949.[317] In the Middle East, the Arab rejection of the United Nations Partition Plan for Palestine and the creation of Israel marked the escalation of the Arab–Israeli conflict. While European powers attempted to retain some or all of their colonial empires, their losses of prestige and resources during the war rendered this unsuccessful, leading to decolonisation.[318][319]
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The global economy suffered heavily from the war, although participating nations were affected differently. The United States emerged much richer than any other nation, leading to a baby boom, and by 1950 its gross domestic product per person was much higher than that of any of the other powers, and it dominated the world economy.[320] The UK and US pursued a policy of industrial disarmament in Western Germany in the years 1945–1948.[321] Because of international trade interdependencies this led to European economic stagnation and delayed European recovery for several years.[322][323]
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Recovery began with the mid-1948 currency reform in Western Germany, and was sped up by the liberalisation of European economic policy that the Marshall Plan (1948–1951) both directly and indirectly caused.[324][325] The post-1948 West German recovery has been called the German economic miracle.[326] Italy also experienced an economic boom[327] and the French economy rebounded.[328] By contrast, the United Kingdom was in a state of economic ruin,[329] and although receiving a quarter of the total Marshall Plan assistance, more than any other European country,[330] it continued in relative economic decline for decades.[331]
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The Soviet Union, despite enormous human and material losses, also experienced rapid increase in production in the immediate post-war era.[332] Japan recovered much later.[333] China returned to its pre-war industrial production by 1952.[334]
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Estimates for the total number of casualties in the war vary, because many deaths went unrecorded.[335] Most suggest that some 60 million people died in the war, including about 20 million military personnel and 40 million civilians.[336][337][338]
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Many of the civilians died because of deliberate genocide, massacres, mass bombings, disease, and starvation.
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The Soviet Union alone lost around 27 million people during the war,[339] including 8.7 million military and 19 million civilian deaths.[340] A quarter of the people in the Soviet Union were wounded or killed.[341] Germany sustained 5.3 million military losses, mostly on the Eastern Front and during the final battles in Germany.[342]
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An estimated 11[343] to 17 million[344] civilians died as a direct or as an indirect result of Nazi racist policies, including mass killing of around 6 million Jews, along with Roma, homosexuals, at least 1.9 million ethnic Poles[345][346] and millions of other Slavs (including Russians, Ukrainians and Belarusians), and other ethnic and minority groups.[347][344] Between 1941 and 1945, more than 200,000 ethnic Serbs, along with gypsies and Jews, were persecuted and murdered by the Axis-aligned Croatian Ustaše in Yugoslavia.[348] Also, more than 100,000 Poles were massacred by the Ukrainian Insurgent Army in the Volhynia massacres, between 1943 and 1945.[349] At the same time about 10,000–15,000 Ukrainians were killed by the Polish Home Army and other Polish units, in reprisal attacks.[350]
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In Asia and the Pacific, between 3 million and more than 10 million civilians, mostly Chinese (estimated at 7.5 million[351]), were killed by the Japanese occupation forces.[352] The most infamous Japanese atrocity was the Nanking Massacre, in which fifty to three hundred thousand Chinese civilians were raped and murdered.[353] Mitsuyoshi Himeta reported that 2.7 million casualties occurred during the Sankō Sakusen. General Yasuji Okamura implemented the policy in Heipei and Shantung.[354]
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Axis forces employed biological and chemical weapons. The Imperial Japanese Army used a variety of such weapons during its invasion and occupation of China (see Unit 731)[355][356] and in early conflicts against the Soviets.[357] Both the Germans and the Japanese tested such weapons against civilians,[358] and sometimes on prisoners of war.[359]
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The Soviet Union was responsible for the Katyn massacre of 22,000 Polish officers,[360] and the imprisonment or execution of thousands of political prisoners by the NKVD, along with mass civilian deportations to Siberia, in the Baltic states and eastern Poland annexed by the Red Army.[361]
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The mass bombing of cities in Europe and Asia has often been called a war crime, although no positive or specific customary international humanitarian law with respect to aerial warfare existed before or during World War II.[362] The USAAF firebombed a total of 67 Japanese cities, killing 393,000 civilians and destroying 65% of built-up areas.[363]
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Nazi Germany was responsible for the Holocaust (which killed approximately 6 million Jews) as well as for killing 2.7 million ethnic Poles[364] and 4 million others who were deemed "unworthy of life" (including the disabled and mentally ill, Soviet prisoners of war, Romani, homosexuals, Freemasons, and Jehovah's Witnesses) as part of a programme of deliberate extermination, in effect becoming a "genocidal state".[365] Soviet POWs were kept in especially unbearable conditions, and 3.6 million Soviet POWs out of 5.7 died in Nazi camps during the war.[366][367] In addition to concentration camps, death camps were created in Nazi Germany to exterminate people on an industrial scale. Nazi Germany extensively used forced labourers; about 12 million Europeans from German occupied countries were abducted and used as a slave work force in German industry, agriculture and war economy.[368]
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The Soviet Gulag became a de facto system of deadly camps during 1942–43, when wartime privation and hunger caused numerous deaths of inmates,[369] including foreign citizens of Poland and other countries occupied in 1939–40 by the Soviet Union, as well as Axis POWs.[370] By the end of the war, most Soviet POWs liberated from Nazi camps and many repatriated civilians were detained in special filtration camps where they were subjected to NKVD evaluation, and 226,127 were sent to the Gulag as real or perceived Nazi collaborators.[371]
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Japanese prisoner-of-war camps, many of which were used as labour camps, also had high death rates. The International Military Tribunal for the Far East found the death rate of Western prisoners was 27 per cent (for American POWs, 37 per cent),[372] seven times that of POWs under the Germans and Italians.[373] While 37,583 prisoners from the UK, 28,500 from the Netherlands, and 14,473 from the United States were released after the surrender of Japan, the number of Chinese released was only 56.[374]
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At least five million Chinese civilians from northern China and Manchukuo were enslaved between 1935 and 1941 by the East Asia Development Board, or Kōain, for work in mines and war industries. After 1942, the number reached 10 million.[375] In Java, between 4 and 10 million rōmusha (Japanese: "manual labourers"), were forced to work by the Japanese military. About 270,000 of these Javanese labourers were sent to other Japanese-held areas in South East Asia, and only 52,000 were repatriated to Java.[376]
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In Europe, occupation came under two forms. In Western, Northern, and Central Europe (France, Norway, Denmark, the Low Countries, and the annexed portions of Czechoslovakia) Germany established economic policies through which it collected roughly 69.5 billion reichmarks (27.8 billion US dollars) by the end of the war; this figure does not include the sizeable plunder of industrial products, military equipment, raw materials and other goods.[377] Thus, the income from occupied nations was over 40 per cent of the income Germany collected from taxation, a figure which increased to nearly 40 per cent of total German income as the war went on.[378]
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In the East, the intended gains of Lebensraum were never attained as fluctuating front-lines and Soviet scorched earth policies denied resources to the German invaders.[379] Unlike in the West, the Nazi racial policy encouraged extreme brutality against what it considered to be the "inferior people" of Slavic descent; most German advances were thus followed by mass executions.[380] Although resistance groups formed in most occupied territories, they did not significantly hamper German operations in either the East[381] or the West[382] until late 1943.
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In Asia, Japan termed nations under its occupation as being part of the Greater East Asia Co-Prosperity Sphere, essentially a Japanese hegemony which it claimed was for purposes of liberating colonised peoples.[383] Although Japanese forces were sometimes welcomed as liberators from European domination, Japanese war crimes frequently turned local public opinion against them.[384] During Japan's initial conquest it captured 4,000,000 barrels (640,000 m3) of oil (~5.5×105 tonnes) left behind by retreating Allied forces, and by 1943 was able to get production in the Dutch East Indies up to 50 million barrels (~6.8×10^6 t), 76 per cent of its 1940 output rate.[384]
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In Europe, before the outbreak of the war, the Allies had significant advantages in both population and economics. In 1938, the Western Allies (United Kingdom, France, Poland and the British Dominions) had a 30 per cent larger population and a 30 per cent higher gross domestic product than the European Axis powers (Germany and Italy); if colonies are included, the Allies had more than a 5:1 advantage in population and a nearly 2:1 advantage in GDP.[385] In Asia at the same time, China had roughly six times the population of Japan but only an 89 per cent higher GDP; this is reduced to three times the population and only a 38 per cent higher GDP if Japanese colonies are included.[385]
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The United States produced about two-thirds of all the munitions used by the Allies in WWII, including warships, transports, warplanes, artillery, tanks, trucks, and ammunition.[386]
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Though the Allies' economic and population advantages were largely mitigated during the initial rapid blitzkrieg attacks of Germany and Japan, they became the decisive factor by 1942, after the United States and Soviet Union joined the Allies, as the war largely settled into one of attrition.[387] While the Allies' ability to out-produce the Axis is often attributed[by whom?] to the Allies having more access to natural resources, other factors, such as Germany and Japan's reluctance to employ women in the labour force,[388] Allied strategic bombing,[389] and Germany's late shift to a war economy[390] contributed significantly. Additionally, neither Germany nor Japan planned to fight a protracted war, and had not equipped themselves to do so.[391] To improve their production, Germany and Japan used millions of slave labourers;[392] Germany used about 12 million people, mostly from Eastern Europe,[368] while Japan used more than 18 million people in Far East Asia.[375][376]
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Aircraft were used for reconnaissance, as fighters, bombers, and ground-support, and each role was advanced considerably. Innovation included airlift (the capability to quickly move limited high-priority supplies, equipment, and personnel);[393] and of strategic bombing (the bombing of enemy industrial and population centres to destroy the enemy's ability to wage war).[394] Anti-aircraft weaponry also advanced, including defences such as radar and surface-to-air artillery. The use of the jet aircraft was pioneered and, though late introduction meant it had little impact, it led to jets becoming standard in air forces worldwide.[395] Although guided missiles were being developed, they were not advanced enough to reliably target aircraft until some years after the war.
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Advances were made in nearly every aspect of naval warfare, most notably with aircraft carriers and submarines. Although aeronautical warfare had relatively little success at the start of the war, actions at Taranto, Pearl Harbor, and the Coral Sea established the carrier as the dominant capital ship in place of the battleship.[396][397][398] In the Atlantic, escort carriers proved to be a vital part of Allied convoys, increasing the effective protection radius and helping to close the Mid-Atlantic gap.[399] Carriers were also more economical than battleships because of the relatively low cost of aircraft[400] and their not requiring to be as heavily armoured.[401] Submarines, which had proved to be an effective weapon during the First World War,[402] were anticipated by all sides to be important in the second. The British focused development on anti-submarine weaponry and tactics, such as sonar and convoys, while Germany focused on improving its offensive capability, with designs such as the Type VII submarine and wolfpack tactics.[403][better source needed] Gradually, improving Allied technologies such as the Leigh light, hedgehog, squid, and homing torpedoes proved victorious over the German submarines.[citation needed]
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Land warfare changed from the static front lines of trench warfare of World War I, which had relied on improved artillery that outmatched the speed of both infantry and cavalry, to increased mobility and combined arms. The tank, which had been used predominantly for infantry support in the First World War, had evolved into the primary weapon.[404] In the late 1930s, tank design was considerably more advanced than it had been during World War I,[405] and advances continued throughout the war with increases in speed, armour and firepower.[citation needed] At the start of the war, most commanders thought enemy tanks should be met by tanks with superior specifications.[406] This idea was challenged by the poor performance of the relatively light early tank guns against armour, and German doctrine of avoiding tank-versus-tank combat. This, along with Germany's use of combined arms, were among the key elements of their highly successful blitzkrieg tactics across Poland and France.[404] Many means of destroying tanks, including indirect artillery, anti-tank guns (both towed and self-propelled), mines, short-ranged infantry antitank weapons, and other tanks were used.[406] Even with large-scale mechanisation, infantry remained the backbone of all forces,[407] and throughout the war, most infantry were equipped similarly to World War I.[408] The portable machine gun spread, a notable example being the German MG34, and various submachine guns which were suited to close combat in urban and jungle settings.[408] The assault rifle, a late war development incorporating many features of the rifle and submachine gun, became the standard postwar infantry weapon for most armed forces.[409]
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Most major belligerents attempted to solve the problems of complexity and security involved in using large codebooks for cryptography by designing ciphering machines, the most well known being the German Enigma machine.[410] Development of SIGINT (signals intelligence) and cryptanalysis enabled the countering process of decryption. Notable examples were the Allied decryption of Japanese naval codes[411] and British Ultra, a pioneering method for decoding Enigma benefiting from information given to the United Kingdom by the Polish Cipher Bureau, which had been decoding early versions of Enigma before the war.[412] Another aspect of military intelligence was the use of deception, which the Allies used to great effect, such as in operations Mincemeat and Bodyguard.[411][413]
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Other technological and engineering feats achieved during, or as a result of, the war include the world's first programmable computers (Z3, Colossus, and ENIAC), guided missiles and modern rockets, the Manhattan Project's development of nuclear weapons, operations research and the development of artificial harbours and oil pipelines under the English Channel.[citation needed] Penicillin was first mass-produced and used during the war (see Stabilization and mass production of penicillin).[414]
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1 |
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Snakes are elongated, legless, carnivorous reptiles of the suborder Serpentes.[2] Like all other squamates, snakes are ectothermic, amniote vertebrates covered in overlapping scales. Many species of snakes have skulls with several more joints than their lizard ancestors, enabling them to swallow prey much larger than their heads with their highly mobile jaws. To accommodate their narrow bodies, snakes' paired organs (such as kidneys) appear one in front of the other instead of side by side, and most have only one functional lung. Some species retain a pelvic girdle with a pair of vestigial claws on either side of the cloaca. Lizards have evolved elongate bodies without limbs or with greatly reduced limbs about twenty-five times independently via convergent evolution, leading to many lineages of legless lizards.[3] Legless lizards resemble snakes, but several common groups of legless lizards have eyelids and external ears, which snakes lack, although this rule is not universal (see Amphisbaenia, Dibamidae, and Pygopodidae).
|
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Living snakes are found on every continent except Antarctica, and on most smaller land masses; exceptions include some large islands, such as Ireland, Iceland, Greenland, the Hawaiian archipelago, and the islands of New Zealand, and many small islands of the Atlantic and central Pacific oceans.[4] Additionally, sea snakes are widespread throughout the Indian and Pacific Oceans. More than 20 families are currently recognized, comprising about 520 genera and about 3,600 species.[5][6] They range in size from the tiny, 10.4 cm (4.1 in)-long Barbados thread snake[7] to the reticulated python of 6.95 meters (22.8 ft) in length.[8] The fossil species Titanoboa cerrejonensis was 12.8 meters (42 ft) long.[9] Snakes are thought to have evolved from either burrowing or aquatic lizards, perhaps during the Jurassic period, with the earliest known fossils dating to between 143 and 167 Ma ago.[10] The diversity of modern snakes appeared during the Paleocene epoch (c 66 to 56 Ma ago, after the Cretaceous–Paleogene extinction event). The oldest preserved descriptions of snakes can be found in the Brooklyn Papyrus.
|
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Most species are nonvenomous and those that have venom use it primarily to kill and subdue prey rather than for self-defense. Some possess venom potent enough to cause painful injury or death to humans. Nonvenomous snakes either swallow prey alive or kill by constriction.
|
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|
9 |
+
The English word snake comes from Old English snaca, itself from Proto-Germanic *snak-an- (cf. Germanic Schnake "ring snake", Swedish snok "grass snake"), from Proto-Indo-European root *(s)nēg-o- "to crawl", "to creep", which also gave sneak as well as Sanskrit nāgá "snake".[11] The word ousted adder, as adder went on to narrow in meaning, though in Old English næddre was the general word for snake.[12] The other term, serpent, is from French, ultimately from Indo-European *serp- (to creep),[13] which also gave Ancient Greek hérpō (ἕρπω) "I crawl".
|
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|
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Leptotyphlopidae
|
12 |
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|
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Anomalepididae
|
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|
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Typhlopidae
|
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|
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Anilius
|
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Tropidophiidae
|
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|
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+
Uropeltidae
|
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+
Anomochilus
|
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+
|
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+
Cylindrophis
|
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|
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+
Pythonidae
|
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+
|
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+
Xenopeltis
|
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+
|
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+
Loxocemus
|
32 |
+
|
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+
Acrochordidae
|
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+
|
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+
Xenodermidae
|
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+
|
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Pareidae
|
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|
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+
Viperidae
|
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|
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+
Homalopsidae
|
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+
|
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+
Lamprophiidae
|
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+
|
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Elapidae
|
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|
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+
Colubridae
|
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|
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+
Boidae
|
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|
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Erycinae
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|
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+
Calabaria
|
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+
|
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+
Ungaliophiinae
|
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+
|
57 |
+
Sanzinia
|
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+
|
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+
Candoia
|
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+
|
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+
The fossil record of snakes is relatively poor because snake skeletons are typically small and fragile making fossilization uncommon. Fossils readily identifiable as snakes (though often retaining hind limbs) first appear in the fossil record during the Cretaceous period.[15] The earliest known true snake fossils (members of the crown group Serpentes) come from the marine simoliophiids, the oldest of which is the Late Cretaceous (Cenomanian age) Haasiophis terrasanctus,[1] dated to between 112 and 94 million years old.[16]
|
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Based on comparative anatomy, there is consensus that snakes descended from lizards.[17]:11[18] Pythons and boas—primitive groups among modern snakes—have vestigial hind limbs: tiny, clawed digits known as anal spurs, which are used to grasp during mating.[17]:11[19] The families Leptotyphlopidae and Typhlopidae also possess remnants of the pelvic girdle, appearing as horny projections when visible.
|
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|
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+
Front limbs are nonexistent in all known snakes. This is caused by the evolution of their Hox genes, controlling limb morphogenesis. The axial skeleton of the snakes’ common ancestor, like most other tetrapods, had regional specializations consisting of cervical (neck), thoracic (chest), lumbar (lower back), sacral (pelvic), and caudal (tail) vertebrae. Early in snake evolution, the Hox gene expression in the axial skeleton responsible for the development of the thorax became dominant. As a result, the vertebrae anterior to the hindlimb buds (when present) all have the same thoracic-like identity (except from the atlas, axis, and 1–3 neck vertebrae). In other words, most of a snake's skeleton is an extremely extended thorax. Ribs are found exclusively on the thoracic vertebrae. Neck, lumbar and pelvic vertebrae are very reduced in number (only 2–10 lumbar and pelvic vertebrae are present), while only a short tail remains of the caudal vertebrae. However, the tail is still long enough to be of important use in many species, and is modified in some aquatic and tree-dwelling species.
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|
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Many modern snake groups originated during the Paleocene, alongside the adaptive radiation of mammals following the extinction of (non-avian) dinosaurs. The expansion of grasslands in North America also led to an explosive radiation among snakes.[20] Previously, snakes were a minor component of the North American fauna, but during the Miocene, the number of species and their prevalence increased dramatically with the first appearances of vipers and elapids in North America and the significant diversification of Colubridae (including the origin of many modern genera such as Nerodia, Lampropeltis, Pituophis, and Pantherophis).[20]
|
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|
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There is fossil evidence to suggest that snakes may have evolved from burrowing lizards, such as the varanids (or a similar group) during the Cretaceous Period.[21] An early fossil snake relative, Najash rionegrina, was a two-legged burrowing animal with a sacrum, and was fully terrestrial.[22] One extant analog of these putative ancestors is the earless monitor Lanthanotus of Borneo (though it also is semiaquatic).[23] Subterranean species evolved bodies streamlined for burrowing, and eventually lost their limbs.[23] According to this hypothesis, features such as the transparent, fused eyelids (brille) and loss of external ears evolved to cope with fossorial difficulties, such as scratched corneas and dirt in the ears.[21][23] Some primitive snakes are known to have possessed hindlimbs, but their pelvic bones lacked a direct connection to the vertebrae. These include fossil species like Haasiophis, Pachyrhachis and Eupodophis, which are slightly older than Najash.[19]
|
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|
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This hypothesis was strengthened in 2015 by the discovery of a 113m year-old fossil of a four-legged snake in Brazil that has been named Tetrapodophis amplectus. It has many snake-like features, is adapted for burrowing and its stomach indicates that it was preying on other animals.[24] It is currently uncertain if Tetrapodophis is a snake or another species, in the squamate order, as a snake-like body has independently evolved at least 26 times. Tetrapodophis does not have distinctive snake features in its spine and skull.[25][26]
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An alternative hypothesis, based on morphology, suggests the ancestors of snakes were related to mosasaurs—extinct aquatic reptiles from the Cretaceous—which in turn are thought to have derived from varanid lizards.[18] According to this hypothesis, the fused, transparent eyelids of snakes are thought to have evolved to combat marine conditions (corneal water loss through osmosis), and the external ears were lost through disuse in an aquatic environment. This ultimately led to an animal similar to today's sea snakes. In the Late Cretaceous, snakes recolonized land, and continued to diversify into today's snakes. Fossilized snake remains are known from early Late Cretaceous marine sediments, which is consistent with this hypothesis; particularly so, as they are older than the terrestrial Najash rionegrina. Similar skull structure, reduced or absent limbs, and other anatomical features found in both mosasaurs and snakes lead to a positive cladistical correlation, although some of these features are shared with varanids.[citation needed]
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Tetrapodophis
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Pachyrhachis
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Eupodophis descouensi
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Eupodophis descouensi hind leg
|
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Genetic studies in recent years have indicated snakes are not as closely related to monitor lizards as was once believed—and therefore not to mosasaurs, the proposed ancestor in the aquatic scenario of their evolution. However, more evidence links mosasaurs to snakes than to varanids. Fragmented remains found from the Jurassic and Early Cretaceous indicate deeper fossil records for these groups, which may potentially refute either hypothesis.[27][28]
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In 2016 two studies reported that limb loss in snakes is associated with DNA mutations in the Zone of Polarizing Activity Regulatory Sequence (ZRS), a regulatory region of the sonic hedgehog gene which is critically required for limb development. More advanced snakes have no remnants of limbs, but basal snakes such as pythons and boas do have traces of highly reduced, vestigial hind limbs. Python embryos even have fully developed hind limb buds, but their later development is stopped by the DNA mutations in the ZRS.[29][30][31][32]
|
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There are over 2,900 species of snakes ranging as far northward as the Arctic Circle in Scandinavia and southward through Australia.[18] Snakes can be found on every continent except Antarctica, in the sea, and as high as 16,000 feet (4,900 m) in the Himalayan Mountains of Asia.[18][33]:143 There are numerous islands from which snakes are absent, such as Ireland, Iceland, and New Zealand[4][33] (although New Zealand's waters are infrequently visited by the yellow-bellied sea snake and the banded sea krait).[34]
|
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All modern snakes are grouped within the suborder Serpentes in Linnean taxonomy, part of the order Squamata, though their precise placement within squamates remains controversial.[5]
|
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The two infraorders of Serpentes are: Alethinophidia and Scolecophidia.[5] This separation is based on morphological characteristics and mitochondrial DNA sequence similarity. Alethinophidia is sometimes split into Henophidia and Caenophidia, with the latter consisting of "colubroid" snakes (colubrids, vipers, elapids, hydrophiids, and atractaspids) and acrochordids, while the other alethinophidian families comprise Henophidia.[35] While not extant today, the Madtsoiidae, a family of giant, primitive, python-like snakes, was around until 50,000 years ago in Australia, represented by genera such as Wonambi.
|
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|
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There are numerous debates in the systematics within the group. For instance, many sources classify Boidae and Pythonidae as one family, while some keep the Elapidae and Hydrophiidae (sea snakes) separate for practical reasons despite their extremely close relation.
|
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|
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Recent molecular studies support the monophyly of the clades of modern snakes, scolecophidians, typhlopids + anomalepidids, alethinophidians, core alethinophidians, uropeltids (Cylindrophis, Anomochilus, uropeltines), macrostomatans, booids, boids, pythonids and caenophidians.[14]
|
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|
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|
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|
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While snakes are limbless reptiles, which evolved from (and are grouped with) lizards, there are many other species of lizards which have lost their limbs independently and superficially look similar to snakes. These include the slowworm and glass snake.
|
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+
|
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Other serpentine tetrapods unrelated to snakes include caecilians (amphibians), amphisbaenians (near-lizard squamates), and the extinct aistopods (amphibians).
|
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+
|
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+
The now extinct Titanoboa cerrejonensis snakes were 12.8 m (42 ft) in length.[9] By comparison, the largest extant snakes are the reticulated python, which measures about 6.95 m (22.8 ft) long,[8] and the green anaconda, which measures about 5.21 m (17.1 ft) long and is considered the heaviest snake on Earth at 97.5 kg (215 lb).[39]
|
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|
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At the other end of the scale, the smallest extant snake is Leptotyphlops carlae, with a length of about 10.4 cm (4.1 in).[7] Most snakes are fairly small animals, approximately 1 m (3.3 ft) in length.[40]
|
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|
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+
Pit vipers, pythons, and some boas have infrared-sensitive receptors in deep grooves on the snout, which allow them to "see" the radiated heat of warm-blooded prey. In pit vipers, the grooves are located between the nostril and the eye in a large "pit" on each side of the head. Other infrared-sensitive snakes have multiple, smaller labial pits lining the upper lip, just below the nostrils.[41]
|
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|
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+
Snakes use smell to track their prey. They smell by using their forked tongues to collect airborne particles, then passing them to the vomeronasal organ or Jacobson's organ in the mouth for examination.[41] The fork in the tongue gives snakes a sort of directional sense of smell and taste simultaneously.[41] They keep their tongues constantly in motion, sampling particles from the air, ground, and water, analyzing the chemicals found, and determining the presence of prey or predators in the local environment. In water-dwelling snakes, such as the anaconda, the tongue functions efficiently underwater.[41]
|
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|
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+
The underside is very sensitive to vibration. This allows snakes to be able to sense approaching animals by detecting faint vibrations in the ground.[41]
|
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|
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Snake vision varies greatly, from only being able to distinguish light from dark to keen eyesight, but the main trend is that their vision is adequate although not sharp, and allows them to track movements.[42] Generally, vision is best in arboreal snakes and weakest in burrowing snakes. Some snakes, such as the Asian vine snake (genus Ahaetulla), have binocular vision, with both eyes capable of focusing on the same point. Most snakes focus by moving the lens back and forth in relation to the retina, while in the other amniote groups, the lens is stretched. Many nocturnal snakes have slit pupils while diurnal snakes have round pupils. Most species possess three visual pigments and are probably able to see two primary colors in daylight. It's concluded that the last common ancestors of all snakes had UV-sensitive vision, but that most snakes that depends on their eyesight to hunt in daylight have evolved lenses that act as sunglasses which filters out UV-light, and probably also sharpens their vision by improving the contrasts.[43]
|
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|
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+
The skin of a snake is covered in scales. Contrary to the popular notion of snakes being slimy because of possible confusion of snakes with worms, snakeskin has a smooth, dry texture. Most snakes use specialized belly scales to travel, gripping surfaces. The body scales may be smooth, keeled, or granular. The eyelids of a snake are transparent "spectacle" scales, which remain permanently closed, also known as brille.
|
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|
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The shedding of scales is called ecdysis (or in normal usage, molting or sloughing). In the case of snakes, the complete outer layer of skin is shed in one layer.[44] Snake scales are not discrete, but extensions of the epidermis—hence they are not shed separately but as a complete outer layer during each molt, akin to a sock being turned inside out.[45]
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Snakes have a wide diversity of skin coloration patterns. These patterns are often related to behavior, such as a tendency to have to flee from predators. Snakes that are plain or have longitudinal stripes often have to escape from predators, with the pattern (or lack thereof) not providing reference points to predators, thus allowing the snake to escape without being notice. Plain snakes usually adopt active hunting strategies, as their pattern allows them to send little information to prey about motion. Blotched snakes, on the other hand, usually use ambush-based strategies, likely because it helps them blend into an environment with irregularly shaped objects, like sticks or rocks. Spotted patterning can similarly help snakes to blend into their environment.[46]
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The shape and number of scales on the head, back, and belly are often characteristic and used for taxonomic purposes. Scales are named mainly according to their positions on the body. In "advanced" (Caenophidian) snakes, the broad belly scales and rows of dorsal scales correspond to the vertebrae, allowing scientists to count the vertebrae without dissection.
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Molting, or ecdysis, serves a number of functions. Firstly, the old and worn skin is replaced; secondly, it helps get rid of parasites such as mites and ticks. Renewal of the skin by molting is supposed to allow growth in some animals such as insects; however, this has been disputed in the case of snakes.[45][47]
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Molting occurs periodically throughout the snake's life. Before a molt, the snake stops eating and often hides or moves to a safe place. Just before shedding, the skin becomes dull and dry looking and the eyes become cloudy or blue-colored. The inner surface of the old skin liquefies. This causes the old skin to separate from the new skin beneath it. After a few days, the eyes clear and the snake "crawls" out of its old skin. The old skin breaks near the mouth and the snake wriggles out, aided by rubbing against rough surfaces. In many cases, the cast skin peels backward over the body from head to tail in one piece, like pulling a sock off inside-out. A new, larger, brighter layer of skin has formed underneath.[45][48]
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An older snake may shed its skin only once or twice a year. But a younger snake, still growing, may shed up to four times a year.[48] The discarded skin gives a perfect imprint of the scale pattern, and it is usually possible to identify the snake if the discarded skin is reasonably intact.[45] This periodic renewal has led to the snake being a symbol of healing and medicine, as pictured in the Rod of Asclepius.[49]
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Scale counts can sometimes be used to tell the sex of a snake when the species is not distinctly sexually dimorphic. A probe is inserted into the cloaca until it can go no further. The probe is marked at the point where it stops, removed, and compared to the subcaudal depth by laying it alongside the scales.[50] The scalation count determines whether the snake is a male or female as hemipenes of a male will probe to a different depth (usually longer) than the cloaca of a female.[50][clarification needed]
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The skeleton of most snakes consists solely of the skull, hyoid, vertebral column, and ribs, though henophidian snakes retain vestiges of the pelvis and rear limbs.
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The skull of the snake consists of a solid and complete neurocranium, to which many of the other bones are only loosely attached, particularly the highly mobile jaw bones, which facilitate manipulation and ingestion of large prey items. The left and right sides of the lower jaw are joined only by a flexible ligament at the anterior tips, allowing them to separate widely, while the posterior end of the lower jaw bones articulate with a quadrate bone, allowing further mobility. The bones of the mandible and quadrate bones can also pick up ground borne vibrations.[51] Because the sides of the jaw can move independently of one another, snakes resting their jaws on a surface have sensitive stereo hearing which can detect the position of prey. The jaw-quadrate-stapes pathway is capable of detecting vibrations on the angstrom scale, despite the absence of an outer ear and the ossicle mechanism of impedance matching used in other vertebrates to receive vibrations from the air.[52][53]
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The hyoid is a small bone located posterior and ventral to the skull, in the 'neck' region, which serves as an attachment for muscles of the snake's tongue, as it does in all other tetrapods.
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The vertebral column consists of anywhere between 200 and 400 (or more) vertebrae. Tail vertebrae are comparatively few in number (often less than 20% of the total) and lack ribs, while body vertebrae each have two ribs articulating with them. The vertebrae have projections that allow for strong muscle attachment enabling locomotion without limbs.
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Autotomy of the tail, a feature found in some lizards is absent in most snakes.[54] Caudal autotomy in snakes is rare and is intervertebral, unlike that in lizards, which is intravertebral—that is, the break happens along a predefined fracture plane present on a vertebra.[55][56]
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In some snakes, most notably boas and pythons, there are vestiges of the hindlimbs in the form of a pair of pelvic spurs. These small, claw-like protrusions on each side of the cloaca are the external portion of the vestigial hindlimb skeleton, which includes the remains of an ilium and femur.
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Snakes are polyphyodonts with teeth that are continuously replaced.[57]
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Snake's and other reptiles have a three-chambered heart that controls the circulatory system via the left and right atrium, and one ventricle.[58] Internally, the ventricle is divided into three interconnected cavities which include the cavum arteriosum, the cavum pulmonale, and the cavum venosum.[59] The cavum venosum receives deoxygenated blood from the right atrium while the cavum arteriosum receives oxygenated blood directly from the left atrium. Located beneath the cavum venosum is the cavum pulmonale, which pumps blood to the pulmonary trunk.[60]
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The snake's heart is encased in a sac, called the pericardium, located at the bifurcation of the bronchi. The heart is able to move around, however, owing to the lack of a diaphragm. This adjustment protects the heart from potential damage when large ingested prey is passed through the esophagus. The spleen is attached to the gall bladder and pancreas and filters the blood. The thymus is located in fatty tissue above the heart and is responsible for the generation of immune cells in the blood. The cardiovascular system of snakes is also unique for the presence of a renal portal system in which the blood from the snake's tail passes through the kidneys before returning to the heart.[61]
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The vestigial left lung is often small or sometimes even absent, as snakes' tubular bodies require all of their organs to be long and thin.[61] In the majority of species, only one lung is functional. This lung contains a vascularized anterior portion and a posterior portion that does not function in gas exchange.[61] This 'saccular lung' is used for hydrostatic purposes to adjust buoyancy in some aquatic snakes and its function remains unknown in terrestrial species.[61] Many organs that are paired, such as kidneys or reproductive organs, are staggered within the body, with one located ahead of the other.[61]
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Snakes have no lymph nodes.[61]
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Cobras, vipers, and closely related species use venom to immobilize, injure or kill their prey. The venom is modified saliva, delivered through fangs.[17]:243 The fangs of 'advanced' venomous snakes like viperids and elapids are hollow to inject venom more effectively, while the fangs of rear-fanged snakes such as the boomslang merely have a groove on the posterior edge to channel venom into the wound. Snake venoms are often prey specific—their role in self-defense is secondary.[17]:243
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Venom, like all salivary secretions, is a predigestant that initiates the breakdown of food into soluble compounds, facilitating proper digestion. Even nonvenomous snake bites (like any animal bite) will cause tissue damage.[17]:209
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Certain birds, mammals, and other snakes (such as kingsnakes) that prey on venomous snakes have developed resistance and even immunity to certain venoms.[17]:243 Venomous snakes include three families of snakes, and do not constitute a formal classification group used in taxonomy.
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The colloquial term "poisonous snake" is generally an incorrect label for snakes. A poison is inhaled or ingested, whereas venom produced by snakes is injected into its victim via fangs.[62] There are, however, two exceptions: Rhabdophis sequesters toxins from the toads it eats, then secretes them from nuchal glands to ward off predators, and a small unusual population of garter snakes in the U.S. state of Oregon retains enough toxins in their livers from the newts they eat to be effectively poisonous to small local predators (such as crows and foxes).[63]
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Snake venoms are complex mixtures of proteins, and are stored in venom glands at the back of the head.[63] In all venomous snakes, these glands open through ducts into grooved or hollow teeth in the upper jaw.[17]:243[62] These proteins can potentially be a mix of neurotoxins (which attack the nervous system), hemotoxins (which attack the circulatory system), cytotoxins, bungarotoxins and many other toxins that affect the body in different ways.[62] Almost all snake venom contains hyaluronidase, an enzyme that ensures rapid diffusion of the venom.[17]:243
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Venomous snakes that use hemotoxins usually have fangs in the front of their mouths, making it easier for them to inject the venom into their victims.[62] Some snakes that use neurotoxins (such as the mangrove snake) have fangs in the back of their mouths, with the fangs curled backwards.[64] This makes it difficult both for the snake to use its venom and for scientists to milk them.[62] Elapids, however, such as cobras and kraits are proteroglyphous—they possess hollow fangs that cannot be erected toward the front of their mouths, and cannot "stab" like a viper. They must actually bite the victim.[17]:242
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It has recently been suggested that all snakes may be venomous to a certain degree, with harmless snakes having weak venom and no fangs.[65] Most snakes currently labelled "nonvenomous" would still be considered harmless according to this theory, as they either lack a venom delivery method or are incapable of delivering enough to endanger a human. This theory postulates that snakes may have evolved from a common lizard ancestor that was venomous—and that venomous lizards like the gila monster, beaded lizard, monitor lizards, and the now-extinct mosasaurs may also have derived from it. They share this venom clade with various other saurian species.
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Venomous snakes are classified in two taxonomic families:
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There is a third family containing the opistoglyphous (rear-fanged) snakes (as well as the majority of other snake species):
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Although a wide range of reproductive modes are used by snakes, all snakes employ internal fertilization. This is accomplished by means of paired, forked hemipenes, which are stored, inverted, in the male's tail.[66] The hemipenes are often grooved, hooked, or spined in order to grip the walls of the female's cloaca.[67][66]
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Most species of snakes lay eggs which they abandon shortly after laying. However, a few species (such as the king cobra) actually construct nests and stay in the vicinity of the hatchlings after incubation.[66] Most pythons coil around their egg-clutches and remain with them until they hatch.[68] A female python will not leave the eggs, except to occasionally bask in the sun or drink water. She will even "shiver" to generate heat to incubate the eggs.[68]
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Some species of snake are ovoviviparous and retain the eggs within their bodies until they are almost ready to hatch.[69][70] Recently, it has been confirmed that several species of snake are fully viviparous, such as the boa constrictor and green anaconda, nourishing their young through a placenta as well as a yolk sac, which is highly unusual among reptiles, or anything else outside of requiem sharks or placental mammals.[69][70] Retention of eggs and live birth are most often associated with colder environments.[66][70]
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Sexual selection in snakes is demonstrated by the 3,000 species that each use different tactics in acquiring mates.[71] Ritual combat between males for the females they want to mate with includes topping, a behavior exhibited by most viperids in which one male will twist around the vertically elevated fore body of its opponent and forcing it downward. It is common for neck biting to occur while the snakes are entwined.[72]
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Parthenogenesis is a natural form of reproduction in which growth and development of embryos occur without fertilization. Agkistrodon contortrix (copperhead) and Agkistrodon piscivorus (cotton mouth) can reproduce by facultative parthenogenesis. That is, they are capable of switching from a sexual mode of reproduction to an asexual mode.[73] The type of parthenogenesis that likely occurs is automixis with terminal fusion, a process in which two terminal products from the same meiosis fuse to form a diploid zygote. This process leads to genome wide homozygosity, expression of deleterious recessive alleles and often to developmental abnormalities. Both captive-born and wild-born A. contortrix and A. piscivorus appear to be capable of this form of parthenogenesis.[73]
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Reproduction in squamate reptiles is almost exclusively sexual. Males ordinarily have a ZZ pair of sex determining chromosomes, and females a ZW pair. However, the Colombian Rainbow boa (Epicrates maurus) can also reproduce by facultative parthenogenesis resulting in production of WW female progeny.[74] The WW females are likely produced by terminal automixis.
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In regions where winters are colder than snakes can tolerate while remaining active, local species will brumate. Unlike hibernation, in which mammals are actually asleep, brumating reptiles are awake but inactive. Individual snakes may brumate in burrows, under rock piles, or inside fallen trees, or snakes may aggregate in large numbers at hibernacula.
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All snakes are strictly carnivorous, eating small animals including lizards, frogs, other snakes, small mammals, birds, eggs, fish, snails, worms or insects.[17][3][18][75] Because snakes cannot bite or tear their food to pieces, they must swallow prey whole. The body size of a snake has a major influence on its eating habits. Smaller snakes eat smaller prey. Juvenile pythons might start out feeding on lizards or mice and graduate to small deer or antelope as an adult, for example.
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The snake's jaw is a complex structure. Contrary to the popular belief that snakes can dislocate their jaws, snakes have a very flexible lower jaw, the two halves of which are not rigidly attached, and numerous other joints in their skull (see snake skull), allowing them to open their mouths wide enough to swallow their prey whole, even if it is larger in diameter than the snake itself.[75] For example, the African egg-eating snake has flexible jaws adapted for eating eggs much larger than the diameter of its head.[17]:81 This snake has no teeth, but does have bony protrusions on the inside edge of its spine, which it uses to break shells when it eats eggs.[17]:81
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While the majority of snakes eat a variety of prey animals, there is some specialization by some species. King cobras and the Australian bandy-bandy consume other snakes. Snakes of the family Pareidae have more teeth on the right side of their mouths than on the left, as the shells of their prey usually spiral clockwise.[17]:184[76][77]
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Some snakes have a venomous bite, which they use to kill their prey before eating it.[75][78] Other snakes kill their prey by constriction.[75] Still others swallow their prey whole and alive.[17]:81[75]
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After eating, snakes become dormant while the process of digestion takes place.[50] Digestion is an intense activity, especially after consumption of large prey. In species that feed only sporadically, the entire intestine enters a reduced state between meals to conserve energy. The digestive system is then 'up-regulated' to full capacity within 48 hours of prey consumption. Being ectothermic ("cold-blooded"), the surrounding temperature plays a large role in snake digestion. The ideal temperature for snakes to digest is 30 °C (86 °F). So much metabolic energy is involved in a snake's digestion that in the South American rattlesnake (Crotalus durissus), surface body temperature increases by as much as 1.2 °C (2.2 °F) during the digestive process.[79] Because of this, a snake disturbed after having eaten recently will often regurgitate its prey to be able to escape the perceived threat. When undisturbed, the digestive process is highly efficient, with the snake's digestive enzymes dissolving and absorbing everything but the prey's hair (or feathers) and claws, which are excreted along with waste.
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The lack of limbs does not impede the movement of snakes. They have developed several different modes of locomotion to deal with particular environments. Unlike the gaits of limbed animals, which form a continuum, each mode of snake locomotion is discrete and distinct from the others; transitions between modes are abrupt.[80][81]
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Lateral undulation is the sole mode of aquatic locomotion, and the most common mode of terrestrial locomotion.[81] In this mode, the body of the snake alternately flexes to the left and right, resulting in a series of rearward-moving "waves".[80] While this movement appears rapid, snakes have rarely been documented moving faster than two body-lengths per second, often much less.[82] This mode of movement has the same net cost of transport (calories burned per meter moved) as running in lizards of the same mass.[83]
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Terrestrial lateral undulation is the most common mode of terrestrial locomotion for most snake species.[80] In this mode, the posteriorly moving waves push against contact points in the environment, such as rocks, twigs, irregularities in the soil, etc.[80] Each of these environmental objects, in turn, generates a reaction force directed forward and towards the midline of the snake, resulting in forward thrust while the lateral components cancel out.[84] The speed of this movement depends upon the density of push-points in the environment, with a medium density of about 8[clarification needed] along the snake's length being ideal.[82] The wave speed is precisely the same as the snake speed, and as a result, every point on the snake's body follows the path of the point ahead of it, allowing snakes to move through very dense vegetation and small openings.[84]
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When swimming, the waves become larger as they move down the snake's body, and the wave travels backwards faster than the snake moves forwards.[85] Thrust is generated by pushing their body against the water, resulting in the observed slip. In spite of overall similarities, studies show that the pattern of muscle activation is different in aquatic versus terrestrial lateral undulation, which justifies calling them separate modes.[86] All snakes can laterally undulate forward (with backward-moving waves), but only sea snakes have been observed reversing the motion (moving backwards with forward-moving waves).[80]
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Most often employed by colubroid snakes (colubrids, elapids, and vipers) when the snake must move in an environment that lacks irregularities to push against (rendering lateral undulation impossible), such as a slick mud flat, or a sand dune, sidewinding is a modified form of lateral undulation in which all of the body segments oriented in one direction remain in contact with the ground, while the other segments are lifted up, resulting in a peculiar "rolling" motion.[87][88] This mode of locomotion overcomes the slippery nature of sand or mud by pushing off with only static portions on the body, thereby minimizing slipping.[87] The static nature of the contact points can be shown from the tracks of a sidewinding snake, which show each belly scale imprint, without any smearing. This mode of locomotion has very low caloric cost, less than ⅓ of the cost for a lizard to move the same distance.[83] Contrary to popular belief, there is no evidence that sidewinding is associated with the sand being hot.[87]
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When push-points are absent, but there is not enough space to use sidewinding because of lateral constraints, such as in tunnels, snakes rely on concertina locomotion.[80][88] In this mode, the snake braces the posterior portion of its body against the tunnel wall while the front of the snake extends and straightens.[87] The front portion then flexes and forms an anchor point, and the posterior is straightened and pulled forwards. This mode of locomotion is slow and very demanding, up to seven times the cost of laterally undulating over the same distance.[83] This high cost is due to the repeated stops and starts of portions of the body as well as the necessity of using active muscular effort to brace against the tunnel walls.
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The movement of snakes in arboreal habitats has only recently been studied.[89] While on tree branches, snakes use several modes of locomotion depending on species and bark texture.[89] In general, snakes will use a modified form of concertina locomotion on smooth branches, but will laterally undulate if contact points are available.[89] Snakes move faster on small branches and when contact points are present, in contrast to limbed animals, which do better on large branches with little 'clutter'.[89]
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Gliding snakes (Chrysopelea) of Southeast Asia launch themselves from branch tips, spreading their ribs and laterally undulating as they glide between trees.[87][90][91] These snakes can perform a controlled glide for hundreds of feet depending upon launch altitude and can even turn in midair.[87][90]
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The slowest mode of snake locomotion is rectilinear locomotion, which is also the only one where the snake does not need to bend its body laterally, though it may do so when turning.[92] In this mode, the belly scales are lifted and pulled forward before being placed down and the body pulled over them. Waves of movement and stasis pass posteriorly, resulting in a series of ripples in the skin.[92] The ribs of the snake do not move in this mode of locomotion and this method is most often used by large pythons, boas, and vipers when stalking prey across open ground as the snake's movements are subtle and harder to detect by their prey in this manner.[87]
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Snakes do not ordinarily prey on humans. Unless startled or injured, most snakes prefer to avoid contact and will not attack humans. With the exception of large constrictors, nonvenomous snakes are not a threat to humans. The bite of a nonvenomous snake is usually harmless; their teeth are not adapted for tearing or inflicting a deep puncture wound, but rather grabbing and holding. Although the possibility of infection and tissue damage is present in the bite of a nonvenomous snake, venomous snakes present far greater hazard to humans.[17]:209 The World Health Organisation (WHO) lists snakebite under the "other neglected conditions" category.[95]
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Documented deaths resulting from snake bites are uncommon. Nonfatal bites from venomous snakes may result in the need for amputation of a limb or part thereof. Of the roughly 725 species of venomous snakes worldwide, only 250 are able to kill a human with one bite. Australia averages only one fatal snake bite per year. In India, 250,000 snakebites are recorded in a single year, with as many as 50,000 recorded initial deaths.[96] The WHO estimates that on the order of 100 000 people die each year as a result of snake bites, and around three times as many amputations and other permanent disabilities are caused by snakebites annually.[97]
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The treatment for a snakebite is as variable as the bite itself. The most common and effective method is through antivenom (or antivenin), a serum made from the venom of the snake. Some antivenom is species-specific (monovalent) while some is made for use with multiple species in mind (polyvalent). In the United States for example, all species of venomous snakes are pit vipers, with the exception of the coral snake. To produce antivenom, a mixture of the venoms of the different species of rattlesnakes, copperheads, and cottonmouths is injected into the body of a horse in ever-increasing dosages until the horse is immunized. Blood is then extracted from the immunized horse. The serum is separated and further purified and freeze-dried. It is reconstituted with sterile water and becomes antivenom. For this reason, people who are allergic to horses are more likely to suffer an allergic reaction to antivenom.[98] Antivenom for the more dangerous species (such as mambas, taipans, and cobras) is made in a similar manner in India, South Africa, and Australia, although these antivenoms are species-specific.
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In some parts of the world, especially in India, snake charming is a roadside show performed by a charmer. In such a show, the snake charmer carries a basket that contains a snake that he seemingly charms by playing tunes from his flutelike musical instrument, to which the snake responds.[99] Snakes lack external ears, though they do have internal ears, and respond to the movement of the flute, not the actual noise.[99]
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The Wildlife Protection Act of 1972 in India technically proscribes snake charming on grounds of reducing animal cruelty. Other snake charmers also have a snake and mongoose show, where both the animals have a mock fight; however, this is not very common, as the snakes, as well as the mongooses, may be seriously injured or killed. Snake charming as a profession is dying out in India because of competition from modern forms of entertainment and environment laws proscribing the practice. Many Indians have never seen snake charming and it is becoming a folktale of the past.[99][100][101][102]
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The Irulas tribe of Andhra Pradesh and Tamil Nadu in India have been hunter-gatherers in the hot, dry plains forests, and have practiced the art of snake catching for generations. They have a vast knowledge of snakes in the field. They generally catch the snakes with the help of a simple stick. Earlier, the Irulas caught thousands of snakes for the snake-skin industry. After the complete ban of the snake-skin industry in India and protection of all snakes under the Indian Wildlife (Protection) Act 1972, they formed the Irula Snake Catcher's Cooperative and switched to catching snakes for removal of venom, releasing them in the wild after four extractions. The venom so collected is used for producing life-saving antivenom, biomedical research and for other medicinal products.[103] The Irulas are also known to eat some of the snakes they catch and are very useful in rat extermination in the villages.
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Despite the existence of snake charmers, there have also been professional snake catchers or wranglers. Modern-day snake trapping involves a herpetologist using a long stick with a V- shaped end. Some television show hosts, like Bill Haast, Austin Stevens, Steve Irwin, and Jeff Corwin, prefer to catch them using bare hands.
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While not commonly thought of as food in most cultures, in others the consumption of snakes is acceptable, or even considered a delicacy. Snake soup of Cantonese cuisine is consumed by locals in autumn, to warm up their body. Western cultures document the consumption of snakes under extreme circumstances of hunger.[104] Cooked rattlesnake meat is an exception, which is commonly consumed in Texas[105] and parts of the Midwestern United States. In Asian countries such as China, Taiwan, Thailand, Indonesia, Vietnam and Cambodia, drinking the blood of snakes—particularly the cobra—is believed to increase sexual virility.[106] The blood is drained while the cobra is still alive when possible, and is usually mixed with some form of liquor to improve the taste.[106]
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In some Asian countries, the use of snakes in alcohol is also accepted. In such cases, the body of a snake or several snakes is left to steep in a jar or container of liquor. It is claimed that this makes the liquor stronger (as well as more expensive). One example of this is the Habu snake sometimes placed in the Okinawan liquor Habushu (ブ酒,) also known as "Habu Sake".[107]
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Snake wine (蛇酒) is an alcoholic beverage produced by infusing whole snakes in rice wine or grain alcohol. The drink was first recorded to have been consumed in China during the Western Zhou dynasty and considered an important curative and believed to reinvigorate a person according to traditional Chinese medicine.[108]
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In the Western world, some snakes (especially docile species such as the ball python and corn snake) are kept as pets. To meet this demand a captive breeding industry has developed. Snakes bred in captivity tend to make better pets and are considered preferable to wild caught specimens.[109] Snakes can be very low maintenance pets, especially compared to more traditional species. They require minimal space, as most common species do not exceed 5 feet (1.5 m) in length. Pet snakes can be fed relatively infrequently, usually once every 5 to 14 days. Certain snakes have a lifespan of more than 40 years if given proper care.
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In ancient Mesopotamia, Nirah, the messenger god of Ištaran, was represented as a serpent on kudurrus, or boundary stones.[110] Representations of two intertwined serpents are common in Sumerian art and Neo-Sumerian artwork[110] and still appear sporadically on cylinder seals and amulets until as late as the thirteenth century BC.[110] The horned viper (Cerastes cerastes) appears in Kassite and Neo-Assyrian kudurrus[110] and is invoked in Assyrian texts as a magical protective entity.[110] A dragon-like creature with horns, the body and neck of a snake, the forelegs of a lion, and the hind-legs of a bird appears in Mesopotamian art from the Akkadian Period until the Hellenistic Period (323 BC–31 BC).[110] This creature, known in Akkadian as the mušḫuššu, meaning "furious serpent", was used as a symbol for particular deities and also as a general protective emblem.[110] It seems to have originally been the attendant of the Underworld god Ninazu,[110] but later became the attendant to the Hurrian storm-god Tishpak, as well as, later, Ninazu's son Ningishzida, the Babylonian national god Marduk, the scribal god Nabu, and the Assyrian national god Ashur.[110]
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In Egyptian history, the snake occupies a primary role with the Nile cobra adorning the crown of the pharaoh in ancient times. It was worshipped as one of the gods and was also used for sinister purposes: murder of an adversary and ritual suicide (Cleopatra).[citation needed] The ouroboros was a well-known ancient Egyptian symbol of a serpent swallowing its own tail.[111] The precursor to the ouroboros was the "Many-Faced",[111] a serpent with five heads, who, according to the Amduat, the oldest surviving Book of the Afterlife, was said to coil around the corpse of the sun god Ra protectively.[111] The earliest surviving depiction of a "true" ouroboros comes from the gilded shrines in the tomb of Tutankhamun.[111] In the early centuries AD, the ouroboros was adopted as a symbol by Gnostic Christians[111] and chapter 136 of the Pistis Sophia, an early Gnostic text, describes "a great dragon whose tail is in its mouth".[111] In medieval alchemy, the ouroboros became a typical western dragon with wings, legs, and a tail.[111]
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In the Bible, King Nahash of Ammon, whose name means "Snake", is depicted very negatively, as a particularly cruel and despicable enemy of the ancient Hebrews.
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The ancient Greeks used the Gorgoneion, a depiction of a hideous face with serpents for hair, as an apotropaic symbol to ward off evil.[112] In a Greek myth described by Pseudo-Apollodorus in his Bibliotheca, Medusa was a Gorgon with serpents for hair whose gaze turned all those who looked at her to stone and was slain by the hero Perseus.[113][114][115] In the Roman poet Ovid's Metamorphoses, Medusa is said to have once been a beautiful priestess of Athena, whom Athena turned into a serpent-haired monster after she was raped by the god Poseidon in Athena's temple.[116] In another myth referenced by the Boeotian poet Hesiod and described in detail by Pseudo-Apollodorus, the hero Heracles is said to have slain the Lernaean Hydra,[117][118] a multiple-headed serpent which dwelt in the swamps of Lerna.[117][118]
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The legendary account of the foundation of Thebes mentioned a monster snake guarding the spring from which the new settlement was to draw its water. In fighting and killing the snake, the companions of the founder Cadmus all perished – leading to the term "Cadmean victory" (i.e. a victory involving one's own ruin).[citation needed]
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Three medical symbols involving snakes that are still used today are Bowl of Hygieia, symbolizing pharmacy, and the Caduceus and Rod of Asclepius, which are symbols denoting medicine in general.[49]
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One of the etymologies proposed for the common female first name Linda is that it might derive from Old German Lindi or Linda, meaning a serpent.
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India is often called the land of snakes and is steeped in tradition regarding snakes.[119] Snakes are worshipped as gods even today with many women pouring milk on snake pits (despite snakes' aversion for milk).[119] The cobra is seen on the neck of Shiva and Vishnu is depicted often as sleeping on a seven-headed snake or within the coils of a serpent.[120] There are also several temples in India solely for cobras sometimes called Nagraj (King of Snakes) and it is believed that snakes are symbols of fertility. There is a Hindu festival called Nag Panchami each year on which day snakes are venerated and prayed to. See also Nāga.[citation needed]
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In India there is another mythology about snakes. Commonly known in Hindi as "Ichchhadhari" snakes. Such snakes can take the form of any living creature, but prefer human form. These mythical snakes possess a valuable gem called "Mani", which is more brilliant than diamond. There are many stories in India about greedy people trying to possess this gem and ending up getting killed.[citation needed]
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The snake is one of the 12 celestial animals of Chinese zodiac, in the Chinese calendar.[121]
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Many ancient Peruvian cultures worshipped nature.[122] They emphasized animals and often depicted snakes in their art.[123]
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Snakes are a part of Hindu worship. A festival, Nag Panchami, in which participants worship either images of or live Nāgas (cobras) is celebrated every year. Most images of Lord Shiva depict snake around his neck. Puranas have various stories associated with snakes. In the Puranas, Shesha is said to hold all the planets of the Universe on his hoods and to constantly sing the glories of Vishnu from all his mouths. He is sometimes referred to as "Ananta-Shesha", which means "Endless Shesha". Other notable snakes in Hinduism are Ananta, Vasuki, Taxak, Karkotaka and Pingala. The term Nāga is used to refer to entities that take the form of large snakes in Hinduism and Buddhism.
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Snakes have also been widely revered, such as in ancient Greece, where the serpent was seen as a healer. Asclepius carried a serpent wound around his wand, a symbol seen today on many ambulances.
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In religious terms, the snake and jaguar are arguably the most important animals in ancient Mesoamerica. "In states of ecstasy, lords dance a serpent dance; great descending snakes adorn and support buildings from Chichen Itza to Tenochtitlan, and the Nahuatl word coatl meaning serpent or twin, forms part of primary deities such as Mixcoatl, Quetzalcoatl, and Coatlicue."[124] In both Maya and Aztec calendars, the fifth day of the week was known as Snake Day.
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In Judaism, the snake of brass is also a symbol of healing, of one's life being saved from imminent death.[125]
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In some parts of Christianity, Christ's redemptive work is compared to saving one's life through beholding the Nehushtan (serpent of brass).[126] Snake handlers use snakes as an integral part of church worship in order to exhibit their faith in divine protection. However, more commonly in Christianity, the serpent has been seen as a representative of evil and sly plotting, which can be seen in the description in Genesis chapter 3 of a snake in the Garden of Eden tempting Eve.[127] Saint Patrick is reputed to have expelled all snakes from Ireland while converting the country to Christianity in the 5th century, thus explaining the absence of snakes there.
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In Christianity and Judaism, the snake makes its infamous appearance in the first book of the Bible when a serpent appears before the first couple Adam and Eve and tempts them with the forbidden fruit from the Tree of Knowledge.[127] The snake returns in Exodus when Moses, as a sign of God's power, turns his staff into a snake and when Moses made the Nehushtan, a bronze snake on a pole that when looked at cured the people of bites from the snakes that plagued them in the desert. The serpent makes its final appearance symbolizing Satan in the Book of Revelation: "And he laid hold on the dragon the old serpent, which is the devil and Satan, and bound him for a thousand years."[128]
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In Neo-Paganism and Wicca, the snake is seen as a symbol of wisdom and knowledge.
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Several compounds from snake venoms are being researched as potential treatments or preventatives for pain, cancers, arthritis, stroke, heart disease, hemophilia, and hypertension, and to control bleeding (e.g. during surgery).[130][131][132]
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Caldwell MW, Nydam RL, Palci A, Apesteguía S (January 2015). "The oldest known snakes from the Middle Jurassic-Lower Cretaceous provide insights on snake evolution". Nature Communications. 6 (5996): 5996. Bibcode:2015NatCo...6.5996C. doi:10.1038/ncomms6996. PMID 25625704.
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1 |
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Snakes are elongated, legless, carnivorous reptiles of the suborder Serpentes.[2] Like all other squamates, snakes are ectothermic, amniote vertebrates covered in overlapping scales. Many species of snakes have skulls with several more joints than their lizard ancestors, enabling them to swallow prey much larger than their heads with their highly mobile jaws. To accommodate their narrow bodies, snakes' paired organs (such as kidneys) appear one in front of the other instead of side by side, and most have only one functional lung. Some species retain a pelvic girdle with a pair of vestigial claws on either side of the cloaca. Lizards have evolved elongate bodies without limbs or with greatly reduced limbs about twenty-five times independently via convergent evolution, leading to many lineages of legless lizards.[3] Legless lizards resemble snakes, but several common groups of legless lizards have eyelids and external ears, which snakes lack, although this rule is not universal (see Amphisbaenia, Dibamidae, and Pygopodidae).
|
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|
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Living snakes are found on every continent except Antarctica, and on most smaller land masses; exceptions include some large islands, such as Ireland, Iceland, Greenland, the Hawaiian archipelago, and the islands of New Zealand, and many small islands of the Atlantic and central Pacific oceans.[4] Additionally, sea snakes are widespread throughout the Indian and Pacific Oceans. More than 20 families are currently recognized, comprising about 520 genera and about 3,600 species.[5][6] They range in size from the tiny, 10.4 cm (4.1 in)-long Barbados thread snake[7] to the reticulated python of 6.95 meters (22.8 ft) in length.[8] The fossil species Titanoboa cerrejonensis was 12.8 meters (42 ft) long.[9] Snakes are thought to have evolved from either burrowing or aquatic lizards, perhaps during the Jurassic period, with the earliest known fossils dating to between 143 and 167 Ma ago.[10] The diversity of modern snakes appeared during the Paleocene epoch (c 66 to 56 Ma ago, after the Cretaceous–Paleogene extinction event). The oldest preserved descriptions of snakes can be found in the Brooklyn Papyrus.
|
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|
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Most species are nonvenomous and those that have venom use it primarily to kill and subdue prey rather than for self-defense. Some possess venom potent enough to cause painful injury or death to humans. Nonvenomous snakes either swallow prey alive or kill by constriction.
|
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|
9 |
+
The English word snake comes from Old English snaca, itself from Proto-Germanic *snak-an- (cf. Germanic Schnake "ring snake", Swedish snok "grass snake"), from Proto-Indo-European root *(s)nēg-o- "to crawl", "to creep", which also gave sneak as well as Sanskrit nāgá "snake".[11] The word ousted adder, as adder went on to narrow in meaning, though in Old English næddre was the general word for snake.[12] The other term, serpent, is from French, ultimately from Indo-European *serp- (to creep),[13] which also gave Ancient Greek hérpō (ἕρπω) "I crawl".
|
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|
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+
Leptotyphlopidae
|
12 |
+
|
13 |
+
Anomalepididae
|
14 |
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|
15 |
+
Typhlopidae
|
16 |
+
|
17 |
+
Anilius
|
18 |
+
|
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+
Tropidophiidae
|
20 |
+
|
21 |
+
Uropeltidae
|
22 |
+
|
23 |
+
Anomochilus
|
24 |
+
|
25 |
+
Cylindrophis
|
26 |
+
|
27 |
+
Pythonidae
|
28 |
+
|
29 |
+
Xenopeltis
|
30 |
+
|
31 |
+
Loxocemus
|
32 |
+
|
33 |
+
Acrochordidae
|
34 |
+
|
35 |
+
Xenodermidae
|
36 |
+
|
37 |
+
Pareidae
|
38 |
+
|
39 |
+
Viperidae
|
40 |
+
|
41 |
+
Homalopsidae
|
42 |
+
|
43 |
+
Lamprophiidae
|
44 |
+
|
45 |
+
Elapidae
|
46 |
+
|
47 |
+
Colubridae
|
48 |
+
|
49 |
+
Boidae
|
50 |
+
|
51 |
+
Erycinae
|
52 |
+
|
53 |
+
Calabaria
|
54 |
+
|
55 |
+
Ungaliophiinae
|
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+
|
57 |
+
Sanzinia
|
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+
|
59 |
+
Candoia
|
60 |
+
|
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+
The fossil record of snakes is relatively poor because snake skeletons are typically small and fragile making fossilization uncommon. Fossils readily identifiable as snakes (though often retaining hind limbs) first appear in the fossil record during the Cretaceous period.[15] The earliest known true snake fossils (members of the crown group Serpentes) come from the marine simoliophiids, the oldest of which is the Late Cretaceous (Cenomanian age) Haasiophis terrasanctus,[1] dated to between 112 and 94 million years old.[16]
|
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+
|
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+
Based on comparative anatomy, there is consensus that snakes descended from lizards.[17]:11[18] Pythons and boas—primitive groups among modern snakes—have vestigial hind limbs: tiny, clawed digits known as anal spurs, which are used to grasp during mating.[17]:11[19] The families Leptotyphlopidae and Typhlopidae also possess remnants of the pelvic girdle, appearing as horny projections when visible.
|
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|
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+
Front limbs are nonexistent in all known snakes. This is caused by the evolution of their Hox genes, controlling limb morphogenesis. The axial skeleton of the snakes’ common ancestor, like most other tetrapods, had regional specializations consisting of cervical (neck), thoracic (chest), lumbar (lower back), sacral (pelvic), and caudal (tail) vertebrae. Early in snake evolution, the Hox gene expression in the axial skeleton responsible for the development of the thorax became dominant. As a result, the vertebrae anterior to the hindlimb buds (when present) all have the same thoracic-like identity (except from the atlas, axis, and 1–3 neck vertebrae). In other words, most of a snake's skeleton is an extremely extended thorax. Ribs are found exclusively on the thoracic vertebrae. Neck, lumbar and pelvic vertebrae are very reduced in number (only 2–10 lumbar and pelvic vertebrae are present), while only a short tail remains of the caudal vertebrae. However, the tail is still long enough to be of important use in many species, and is modified in some aquatic and tree-dwelling species.
|
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+
|
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+
Many modern snake groups originated during the Paleocene, alongside the adaptive radiation of mammals following the extinction of (non-avian) dinosaurs. The expansion of grasslands in North America also led to an explosive radiation among snakes.[20] Previously, snakes were a minor component of the North American fauna, but during the Miocene, the number of species and their prevalence increased dramatically with the first appearances of vipers and elapids in North America and the significant diversification of Colubridae (including the origin of many modern genera such as Nerodia, Lampropeltis, Pituophis, and Pantherophis).[20]
|
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|
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+
There is fossil evidence to suggest that snakes may have evolved from burrowing lizards, such as the varanids (or a similar group) during the Cretaceous Period.[21] An early fossil snake relative, Najash rionegrina, was a two-legged burrowing animal with a sacrum, and was fully terrestrial.[22] One extant analog of these putative ancestors is the earless monitor Lanthanotus of Borneo (though it also is semiaquatic).[23] Subterranean species evolved bodies streamlined for burrowing, and eventually lost their limbs.[23] According to this hypothesis, features such as the transparent, fused eyelids (brille) and loss of external ears evolved to cope with fossorial difficulties, such as scratched corneas and dirt in the ears.[21][23] Some primitive snakes are known to have possessed hindlimbs, but their pelvic bones lacked a direct connection to the vertebrae. These include fossil species like Haasiophis, Pachyrhachis and Eupodophis, which are slightly older than Najash.[19]
|
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|
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This hypothesis was strengthened in 2015 by the discovery of a 113m year-old fossil of a four-legged snake in Brazil that has been named Tetrapodophis amplectus. It has many snake-like features, is adapted for burrowing and its stomach indicates that it was preying on other animals.[24] It is currently uncertain if Tetrapodophis is a snake or another species, in the squamate order, as a snake-like body has independently evolved at least 26 times. Tetrapodophis does not have distinctive snake features in its spine and skull.[25][26]
|
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|
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An alternative hypothesis, based on morphology, suggests the ancestors of snakes were related to mosasaurs—extinct aquatic reptiles from the Cretaceous—which in turn are thought to have derived from varanid lizards.[18] According to this hypothesis, the fused, transparent eyelids of snakes are thought to have evolved to combat marine conditions (corneal water loss through osmosis), and the external ears were lost through disuse in an aquatic environment. This ultimately led to an animal similar to today's sea snakes. In the Late Cretaceous, snakes recolonized land, and continued to diversify into today's snakes. Fossilized snake remains are known from early Late Cretaceous marine sediments, which is consistent with this hypothesis; particularly so, as they are older than the terrestrial Najash rionegrina. Similar skull structure, reduced or absent limbs, and other anatomical features found in both mosasaurs and snakes lead to a positive cladistical correlation, although some of these features are shared with varanids.[citation needed]
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|
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Tetrapodophis
|
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|
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Pachyrhachis
|
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|
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Eupodophis descouensi
|
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|
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+
Eupodophis descouensi hind leg
|
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Genetic studies in recent years have indicated snakes are not as closely related to monitor lizards as was once believed—and therefore not to mosasaurs, the proposed ancestor in the aquatic scenario of their evolution. However, more evidence links mosasaurs to snakes than to varanids. Fragmented remains found from the Jurassic and Early Cretaceous indicate deeper fossil records for these groups, which may potentially refute either hypothesis.[27][28]
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In 2016 two studies reported that limb loss in snakes is associated with DNA mutations in the Zone of Polarizing Activity Regulatory Sequence (ZRS), a regulatory region of the sonic hedgehog gene which is critically required for limb development. More advanced snakes have no remnants of limbs, but basal snakes such as pythons and boas do have traces of highly reduced, vestigial hind limbs. Python embryos even have fully developed hind limb buds, but their later development is stopped by the DNA mutations in the ZRS.[29][30][31][32]
|
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|
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There are over 2,900 species of snakes ranging as far northward as the Arctic Circle in Scandinavia and southward through Australia.[18] Snakes can be found on every continent except Antarctica, in the sea, and as high as 16,000 feet (4,900 m) in the Himalayan Mountains of Asia.[18][33]:143 There are numerous islands from which snakes are absent, such as Ireland, Iceland, and New Zealand[4][33] (although New Zealand's waters are infrequently visited by the yellow-bellied sea snake and the banded sea krait).[34]
|
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|
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+
All modern snakes are grouped within the suborder Serpentes in Linnean taxonomy, part of the order Squamata, though their precise placement within squamates remains controversial.[5]
|
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|
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+
The two infraorders of Serpentes are: Alethinophidia and Scolecophidia.[5] This separation is based on morphological characteristics and mitochondrial DNA sequence similarity. Alethinophidia is sometimes split into Henophidia and Caenophidia, with the latter consisting of "colubroid" snakes (colubrids, vipers, elapids, hydrophiids, and atractaspids) and acrochordids, while the other alethinophidian families comprise Henophidia.[35] While not extant today, the Madtsoiidae, a family of giant, primitive, python-like snakes, was around until 50,000 years ago in Australia, represented by genera such as Wonambi.
|
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+
|
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+
There are numerous debates in the systematics within the group. For instance, many sources classify Boidae and Pythonidae as one family, while some keep the Elapidae and Hydrophiidae (sea snakes) separate for practical reasons despite their extremely close relation.
|
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+
|
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+
Recent molecular studies support the monophyly of the clades of modern snakes, scolecophidians, typhlopids + anomalepidids, alethinophidians, core alethinophidians, uropeltids (Cylindrophis, Anomochilus, uropeltines), macrostomatans, booids, boids, pythonids and caenophidians.[14]
|
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|
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|
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+
|
99 |
+
While snakes are limbless reptiles, which evolved from (and are grouped with) lizards, there are many other species of lizards which have lost their limbs independently and superficially look similar to snakes. These include the slowworm and glass snake.
|
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+
|
101 |
+
Other serpentine tetrapods unrelated to snakes include caecilians (amphibians), amphisbaenians (near-lizard squamates), and the extinct aistopods (amphibians).
|
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+
|
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+
The now extinct Titanoboa cerrejonensis snakes were 12.8 m (42 ft) in length.[9] By comparison, the largest extant snakes are the reticulated python, which measures about 6.95 m (22.8 ft) long,[8] and the green anaconda, which measures about 5.21 m (17.1 ft) long and is considered the heaviest snake on Earth at 97.5 kg (215 lb).[39]
|
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+
|
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+
At the other end of the scale, the smallest extant snake is Leptotyphlops carlae, with a length of about 10.4 cm (4.1 in).[7] Most snakes are fairly small animals, approximately 1 m (3.3 ft) in length.[40]
|
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+
|
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+
Pit vipers, pythons, and some boas have infrared-sensitive receptors in deep grooves on the snout, which allow them to "see" the radiated heat of warm-blooded prey. In pit vipers, the grooves are located between the nostril and the eye in a large "pit" on each side of the head. Other infrared-sensitive snakes have multiple, smaller labial pits lining the upper lip, just below the nostrils.[41]
|
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+
|
109 |
+
Snakes use smell to track their prey. They smell by using their forked tongues to collect airborne particles, then passing them to the vomeronasal organ or Jacobson's organ in the mouth for examination.[41] The fork in the tongue gives snakes a sort of directional sense of smell and taste simultaneously.[41] They keep their tongues constantly in motion, sampling particles from the air, ground, and water, analyzing the chemicals found, and determining the presence of prey or predators in the local environment. In water-dwelling snakes, such as the anaconda, the tongue functions efficiently underwater.[41]
|
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+
|
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The underside is very sensitive to vibration. This allows snakes to be able to sense approaching animals by detecting faint vibrations in the ground.[41]
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Snake vision varies greatly, from only being able to distinguish light from dark to keen eyesight, but the main trend is that their vision is adequate although not sharp, and allows them to track movements.[42] Generally, vision is best in arboreal snakes and weakest in burrowing snakes. Some snakes, such as the Asian vine snake (genus Ahaetulla), have binocular vision, with both eyes capable of focusing on the same point. Most snakes focus by moving the lens back and forth in relation to the retina, while in the other amniote groups, the lens is stretched. Many nocturnal snakes have slit pupils while diurnal snakes have round pupils. Most species possess three visual pigments and are probably able to see two primary colors in daylight. It's concluded that the last common ancestors of all snakes had UV-sensitive vision, but that most snakes that depends on their eyesight to hunt in daylight have evolved lenses that act as sunglasses which filters out UV-light, and probably also sharpens their vision by improving the contrasts.[43]
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The skin of a snake is covered in scales. Contrary to the popular notion of snakes being slimy because of possible confusion of snakes with worms, snakeskin has a smooth, dry texture. Most snakes use specialized belly scales to travel, gripping surfaces. The body scales may be smooth, keeled, or granular. The eyelids of a snake are transparent "spectacle" scales, which remain permanently closed, also known as brille.
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The shedding of scales is called ecdysis (or in normal usage, molting or sloughing). In the case of snakes, the complete outer layer of skin is shed in one layer.[44] Snake scales are not discrete, but extensions of the epidermis—hence they are not shed separately but as a complete outer layer during each molt, akin to a sock being turned inside out.[45]
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Snakes have a wide diversity of skin coloration patterns. These patterns are often related to behavior, such as a tendency to have to flee from predators. Snakes that are plain or have longitudinal stripes often have to escape from predators, with the pattern (or lack thereof) not providing reference points to predators, thus allowing the snake to escape without being notice. Plain snakes usually adopt active hunting strategies, as their pattern allows them to send little information to prey about motion. Blotched snakes, on the other hand, usually use ambush-based strategies, likely because it helps them blend into an environment with irregularly shaped objects, like sticks or rocks. Spotted patterning can similarly help snakes to blend into their environment.[46]
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The shape and number of scales on the head, back, and belly are often characteristic and used for taxonomic purposes. Scales are named mainly according to their positions on the body. In "advanced" (Caenophidian) snakes, the broad belly scales and rows of dorsal scales correspond to the vertebrae, allowing scientists to count the vertebrae without dissection.
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Molting, or ecdysis, serves a number of functions. Firstly, the old and worn skin is replaced; secondly, it helps get rid of parasites such as mites and ticks. Renewal of the skin by molting is supposed to allow growth in some animals such as insects; however, this has been disputed in the case of snakes.[45][47]
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Molting occurs periodically throughout the snake's life. Before a molt, the snake stops eating and often hides or moves to a safe place. Just before shedding, the skin becomes dull and dry looking and the eyes become cloudy or blue-colored. The inner surface of the old skin liquefies. This causes the old skin to separate from the new skin beneath it. After a few days, the eyes clear and the snake "crawls" out of its old skin. The old skin breaks near the mouth and the snake wriggles out, aided by rubbing against rough surfaces. In many cases, the cast skin peels backward over the body from head to tail in one piece, like pulling a sock off inside-out. A new, larger, brighter layer of skin has formed underneath.[45][48]
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An older snake may shed its skin only once or twice a year. But a younger snake, still growing, may shed up to four times a year.[48] The discarded skin gives a perfect imprint of the scale pattern, and it is usually possible to identify the snake if the discarded skin is reasonably intact.[45] This periodic renewal has led to the snake being a symbol of healing and medicine, as pictured in the Rod of Asclepius.[49]
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Scale counts can sometimes be used to tell the sex of a snake when the species is not distinctly sexually dimorphic. A probe is inserted into the cloaca until it can go no further. The probe is marked at the point where it stops, removed, and compared to the subcaudal depth by laying it alongside the scales.[50] The scalation count determines whether the snake is a male or female as hemipenes of a male will probe to a different depth (usually longer) than the cloaca of a female.[50][clarification needed]
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The skeleton of most snakes consists solely of the skull, hyoid, vertebral column, and ribs, though henophidian snakes retain vestiges of the pelvis and rear limbs.
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The skull of the snake consists of a solid and complete neurocranium, to which many of the other bones are only loosely attached, particularly the highly mobile jaw bones, which facilitate manipulation and ingestion of large prey items. The left and right sides of the lower jaw are joined only by a flexible ligament at the anterior tips, allowing them to separate widely, while the posterior end of the lower jaw bones articulate with a quadrate bone, allowing further mobility. The bones of the mandible and quadrate bones can also pick up ground borne vibrations.[51] Because the sides of the jaw can move independently of one another, snakes resting their jaws on a surface have sensitive stereo hearing which can detect the position of prey. The jaw-quadrate-stapes pathway is capable of detecting vibrations on the angstrom scale, despite the absence of an outer ear and the ossicle mechanism of impedance matching used in other vertebrates to receive vibrations from the air.[52][53]
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The hyoid is a small bone located posterior and ventral to the skull, in the 'neck' region, which serves as an attachment for muscles of the snake's tongue, as it does in all other tetrapods.
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The vertebral column consists of anywhere between 200 and 400 (or more) vertebrae. Tail vertebrae are comparatively few in number (often less than 20% of the total) and lack ribs, while body vertebrae each have two ribs articulating with them. The vertebrae have projections that allow for strong muscle attachment enabling locomotion without limbs.
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Autotomy of the tail, a feature found in some lizards is absent in most snakes.[54] Caudal autotomy in snakes is rare and is intervertebral, unlike that in lizards, which is intravertebral—that is, the break happens along a predefined fracture plane present on a vertebra.[55][56]
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In some snakes, most notably boas and pythons, there are vestiges of the hindlimbs in the form of a pair of pelvic spurs. These small, claw-like protrusions on each side of the cloaca are the external portion of the vestigial hindlimb skeleton, which includes the remains of an ilium and femur.
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Snakes are polyphyodonts with teeth that are continuously replaced.[57]
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Snake's and other reptiles have a three-chambered heart that controls the circulatory system via the left and right atrium, and one ventricle.[58] Internally, the ventricle is divided into three interconnected cavities which include the cavum arteriosum, the cavum pulmonale, and the cavum venosum.[59] The cavum venosum receives deoxygenated blood from the right atrium while the cavum arteriosum receives oxygenated blood directly from the left atrium. Located beneath the cavum venosum is the cavum pulmonale, which pumps blood to the pulmonary trunk.[60]
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The snake's heart is encased in a sac, called the pericardium, located at the bifurcation of the bronchi. The heart is able to move around, however, owing to the lack of a diaphragm. This adjustment protects the heart from potential damage when large ingested prey is passed through the esophagus. The spleen is attached to the gall bladder and pancreas and filters the blood. The thymus is located in fatty tissue above the heart and is responsible for the generation of immune cells in the blood. The cardiovascular system of snakes is also unique for the presence of a renal portal system in which the blood from the snake's tail passes through the kidneys before returning to the heart.[61]
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The vestigial left lung is often small or sometimes even absent, as snakes' tubular bodies require all of their organs to be long and thin.[61] In the majority of species, only one lung is functional. This lung contains a vascularized anterior portion and a posterior portion that does not function in gas exchange.[61] This 'saccular lung' is used for hydrostatic purposes to adjust buoyancy in some aquatic snakes and its function remains unknown in terrestrial species.[61] Many organs that are paired, such as kidneys or reproductive organs, are staggered within the body, with one located ahead of the other.[61]
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Snakes have no lymph nodes.[61]
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Cobras, vipers, and closely related species use venom to immobilize, injure or kill their prey. The venom is modified saliva, delivered through fangs.[17]:243 The fangs of 'advanced' venomous snakes like viperids and elapids are hollow to inject venom more effectively, while the fangs of rear-fanged snakes such as the boomslang merely have a groove on the posterior edge to channel venom into the wound. Snake venoms are often prey specific—their role in self-defense is secondary.[17]:243
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Venom, like all salivary secretions, is a predigestant that initiates the breakdown of food into soluble compounds, facilitating proper digestion. Even nonvenomous snake bites (like any animal bite) will cause tissue damage.[17]:209
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Certain birds, mammals, and other snakes (such as kingsnakes) that prey on venomous snakes have developed resistance and even immunity to certain venoms.[17]:243 Venomous snakes include three families of snakes, and do not constitute a formal classification group used in taxonomy.
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The colloquial term "poisonous snake" is generally an incorrect label for snakes. A poison is inhaled or ingested, whereas venom produced by snakes is injected into its victim via fangs.[62] There are, however, two exceptions: Rhabdophis sequesters toxins from the toads it eats, then secretes them from nuchal glands to ward off predators, and a small unusual population of garter snakes in the U.S. state of Oregon retains enough toxins in their livers from the newts they eat to be effectively poisonous to small local predators (such as crows and foxes).[63]
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Snake venoms are complex mixtures of proteins, and are stored in venom glands at the back of the head.[63] In all venomous snakes, these glands open through ducts into grooved or hollow teeth in the upper jaw.[17]:243[62] These proteins can potentially be a mix of neurotoxins (which attack the nervous system), hemotoxins (which attack the circulatory system), cytotoxins, bungarotoxins and many other toxins that affect the body in different ways.[62] Almost all snake venom contains hyaluronidase, an enzyme that ensures rapid diffusion of the venom.[17]:243
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Venomous snakes that use hemotoxins usually have fangs in the front of their mouths, making it easier for them to inject the venom into their victims.[62] Some snakes that use neurotoxins (such as the mangrove snake) have fangs in the back of their mouths, with the fangs curled backwards.[64] This makes it difficult both for the snake to use its venom and for scientists to milk them.[62] Elapids, however, such as cobras and kraits are proteroglyphous—they possess hollow fangs that cannot be erected toward the front of their mouths, and cannot "stab" like a viper. They must actually bite the victim.[17]:242
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It has recently been suggested that all snakes may be venomous to a certain degree, with harmless snakes having weak venom and no fangs.[65] Most snakes currently labelled "nonvenomous" would still be considered harmless according to this theory, as they either lack a venom delivery method or are incapable of delivering enough to endanger a human. This theory postulates that snakes may have evolved from a common lizard ancestor that was venomous—and that venomous lizards like the gila monster, beaded lizard, monitor lizards, and the now-extinct mosasaurs may also have derived from it. They share this venom clade with various other saurian species.
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Venomous snakes are classified in two taxonomic families:
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There is a third family containing the opistoglyphous (rear-fanged) snakes (as well as the majority of other snake species):
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Although a wide range of reproductive modes are used by snakes, all snakes employ internal fertilization. This is accomplished by means of paired, forked hemipenes, which are stored, inverted, in the male's tail.[66] The hemipenes are often grooved, hooked, or spined in order to grip the walls of the female's cloaca.[67][66]
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Most species of snakes lay eggs which they abandon shortly after laying. However, a few species (such as the king cobra) actually construct nests and stay in the vicinity of the hatchlings after incubation.[66] Most pythons coil around their egg-clutches and remain with them until they hatch.[68] A female python will not leave the eggs, except to occasionally bask in the sun or drink water. She will even "shiver" to generate heat to incubate the eggs.[68]
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Some species of snake are ovoviviparous and retain the eggs within their bodies until they are almost ready to hatch.[69][70] Recently, it has been confirmed that several species of snake are fully viviparous, such as the boa constrictor and green anaconda, nourishing their young through a placenta as well as a yolk sac, which is highly unusual among reptiles, or anything else outside of requiem sharks or placental mammals.[69][70] Retention of eggs and live birth are most often associated with colder environments.[66][70]
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Sexual selection in snakes is demonstrated by the 3,000 species that each use different tactics in acquiring mates.[71] Ritual combat between males for the females they want to mate with includes topping, a behavior exhibited by most viperids in which one male will twist around the vertically elevated fore body of its opponent and forcing it downward. It is common for neck biting to occur while the snakes are entwined.[72]
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Parthenogenesis is a natural form of reproduction in which growth and development of embryos occur without fertilization. Agkistrodon contortrix (copperhead) and Agkistrodon piscivorus (cotton mouth) can reproduce by facultative parthenogenesis. That is, they are capable of switching from a sexual mode of reproduction to an asexual mode.[73] The type of parthenogenesis that likely occurs is automixis with terminal fusion, a process in which two terminal products from the same meiosis fuse to form a diploid zygote. This process leads to genome wide homozygosity, expression of deleterious recessive alleles and often to developmental abnormalities. Both captive-born and wild-born A. contortrix and A. piscivorus appear to be capable of this form of parthenogenesis.[73]
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Reproduction in squamate reptiles is almost exclusively sexual. Males ordinarily have a ZZ pair of sex determining chromosomes, and females a ZW pair. However, the Colombian Rainbow boa (Epicrates maurus) can also reproduce by facultative parthenogenesis resulting in production of WW female progeny.[74] The WW females are likely produced by terminal automixis.
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In regions where winters are colder than snakes can tolerate while remaining active, local species will brumate. Unlike hibernation, in which mammals are actually asleep, brumating reptiles are awake but inactive. Individual snakes may brumate in burrows, under rock piles, or inside fallen trees, or snakes may aggregate in large numbers at hibernacula.
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All snakes are strictly carnivorous, eating small animals including lizards, frogs, other snakes, small mammals, birds, eggs, fish, snails, worms or insects.[17][3][18][75] Because snakes cannot bite or tear their food to pieces, they must swallow prey whole. The body size of a snake has a major influence on its eating habits. Smaller snakes eat smaller prey. Juvenile pythons might start out feeding on lizards or mice and graduate to small deer or antelope as an adult, for example.
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The snake's jaw is a complex structure. Contrary to the popular belief that snakes can dislocate their jaws, snakes have a very flexible lower jaw, the two halves of which are not rigidly attached, and numerous other joints in their skull (see snake skull), allowing them to open their mouths wide enough to swallow their prey whole, even if it is larger in diameter than the snake itself.[75] For example, the African egg-eating snake has flexible jaws adapted for eating eggs much larger than the diameter of its head.[17]:81 This snake has no teeth, but does have bony protrusions on the inside edge of its spine, which it uses to break shells when it eats eggs.[17]:81
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While the majority of snakes eat a variety of prey animals, there is some specialization by some species. King cobras and the Australian bandy-bandy consume other snakes. Snakes of the family Pareidae have more teeth on the right side of their mouths than on the left, as the shells of their prey usually spiral clockwise.[17]:184[76][77]
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Some snakes have a venomous bite, which they use to kill their prey before eating it.[75][78] Other snakes kill their prey by constriction.[75] Still others swallow their prey whole and alive.[17]:81[75]
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After eating, snakes become dormant while the process of digestion takes place.[50] Digestion is an intense activity, especially after consumption of large prey. In species that feed only sporadically, the entire intestine enters a reduced state between meals to conserve energy. The digestive system is then 'up-regulated' to full capacity within 48 hours of prey consumption. Being ectothermic ("cold-blooded"), the surrounding temperature plays a large role in snake digestion. The ideal temperature for snakes to digest is 30 °C (86 °F). So much metabolic energy is involved in a snake's digestion that in the South American rattlesnake (Crotalus durissus), surface body temperature increases by as much as 1.2 °C (2.2 °F) during the digestive process.[79] Because of this, a snake disturbed after having eaten recently will often regurgitate its prey to be able to escape the perceived threat. When undisturbed, the digestive process is highly efficient, with the snake's digestive enzymes dissolving and absorbing everything but the prey's hair (or feathers) and claws, which are excreted along with waste.
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The lack of limbs does not impede the movement of snakes. They have developed several different modes of locomotion to deal with particular environments. Unlike the gaits of limbed animals, which form a continuum, each mode of snake locomotion is discrete and distinct from the others; transitions between modes are abrupt.[80][81]
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Lateral undulation is the sole mode of aquatic locomotion, and the most common mode of terrestrial locomotion.[81] In this mode, the body of the snake alternately flexes to the left and right, resulting in a series of rearward-moving "waves".[80] While this movement appears rapid, snakes have rarely been documented moving faster than two body-lengths per second, often much less.[82] This mode of movement has the same net cost of transport (calories burned per meter moved) as running in lizards of the same mass.[83]
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Terrestrial lateral undulation is the most common mode of terrestrial locomotion for most snake species.[80] In this mode, the posteriorly moving waves push against contact points in the environment, such as rocks, twigs, irregularities in the soil, etc.[80] Each of these environmental objects, in turn, generates a reaction force directed forward and towards the midline of the snake, resulting in forward thrust while the lateral components cancel out.[84] The speed of this movement depends upon the density of push-points in the environment, with a medium density of about 8[clarification needed] along the snake's length being ideal.[82] The wave speed is precisely the same as the snake speed, and as a result, every point on the snake's body follows the path of the point ahead of it, allowing snakes to move through very dense vegetation and small openings.[84]
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When swimming, the waves become larger as they move down the snake's body, and the wave travels backwards faster than the snake moves forwards.[85] Thrust is generated by pushing their body against the water, resulting in the observed slip. In spite of overall similarities, studies show that the pattern of muscle activation is different in aquatic versus terrestrial lateral undulation, which justifies calling them separate modes.[86] All snakes can laterally undulate forward (with backward-moving waves), but only sea snakes have been observed reversing the motion (moving backwards with forward-moving waves).[80]
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Most often employed by colubroid snakes (colubrids, elapids, and vipers) when the snake must move in an environment that lacks irregularities to push against (rendering lateral undulation impossible), such as a slick mud flat, or a sand dune, sidewinding is a modified form of lateral undulation in which all of the body segments oriented in one direction remain in contact with the ground, while the other segments are lifted up, resulting in a peculiar "rolling" motion.[87][88] This mode of locomotion overcomes the slippery nature of sand or mud by pushing off with only static portions on the body, thereby minimizing slipping.[87] The static nature of the contact points can be shown from the tracks of a sidewinding snake, which show each belly scale imprint, without any smearing. This mode of locomotion has very low caloric cost, less than ⅓ of the cost for a lizard to move the same distance.[83] Contrary to popular belief, there is no evidence that sidewinding is associated with the sand being hot.[87]
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When push-points are absent, but there is not enough space to use sidewinding because of lateral constraints, such as in tunnels, snakes rely on concertina locomotion.[80][88] In this mode, the snake braces the posterior portion of its body against the tunnel wall while the front of the snake extends and straightens.[87] The front portion then flexes and forms an anchor point, and the posterior is straightened and pulled forwards. This mode of locomotion is slow and very demanding, up to seven times the cost of laterally undulating over the same distance.[83] This high cost is due to the repeated stops and starts of portions of the body as well as the necessity of using active muscular effort to brace against the tunnel walls.
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The movement of snakes in arboreal habitats has only recently been studied.[89] While on tree branches, snakes use several modes of locomotion depending on species and bark texture.[89] In general, snakes will use a modified form of concertina locomotion on smooth branches, but will laterally undulate if contact points are available.[89] Snakes move faster on small branches and when contact points are present, in contrast to limbed animals, which do better on large branches with little 'clutter'.[89]
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Gliding snakes (Chrysopelea) of Southeast Asia launch themselves from branch tips, spreading their ribs and laterally undulating as they glide between trees.[87][90][91] These snakes can perform a controlled glide for hundreds of feet depending upon launch altitude and can even turn in midair.[87][90]
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The slowest mode of snake locomotion is rectilinear locomotion, which is also the only one where the snake does not need to bend its body laterally, though it may do so when turning.[92] In this mode, the belly scales are lifted and pulled forward before being placed down and the body pulled over them. Waves of movement and stasis pass posteriorly, resulting in a series of ripples in the skin.[92] The ribs of the snake do not move in this mode of locomotion and this method is most often used by large pythons, boas, and vipers when stalking prey across open ground as the snake's movements are subtle and harder to detect by their prey in this manner.[87]
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Snakes do not ordinarily prey on humans. Unless startled or injured, most snakes prefer to avoid contact and will not attack humans. With the exception of large constrictors, nonvenomous snakes are not a threat to humans. The bite of a nonvenomous snake is usually harmless; their teeth are not adapted for tearing or inflicting a deep puncture wound, but rather grabbing and holding. Although the possibility of infection and tissue damage is present in the bite of a nonvenomous snake, venomous snakes present far greater hazard to humans.[17]:209 The World Health Organisation (WHO) lists snakebite under the "other neglected conditions" category.[95]
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Documented deaths resulting from snake bites are uncommon. Nonfatal bites from venomous snakes may result in the need for amputation of a limb or part thereof. Of the roughly 725 species of venomous snakes worldwide, only 250 are able to kill a human with one bite. Australia averages only one fatal snake bite per year. In India, 250,000 snakebites are recorded in a single year, with as many as 50,000 recorded initial deaths.[96] The WHO estimates that on the order of 100 000 people die each year as a result of snake bites, and around three times as many amputations and other permanent disabilities are caused by snakebites annually.[97]
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The treatment for a snakebite is as variable as the bite itself. The most common and effective method is through antivenom (or antivenin), a serum made from the venom of the snake. Some antivenom is species-specific (monovalent) while some is made for use with multiple species in mind (polyvalent). In the United States for example, all species of venomous snakes are pit vipers, with the exception of the coral snake. To produce antivenom, a mixture of the venoms of the different species of rattlesnakes, copperheads, and cottonmouths is injected into the body of a horse in ever-increasing dosages until the horse is immunized. Blood is then extracted from the immunized horse. The serum is separated and further purified and freeze-dried. It is reconstituted with sterile water and becomes antivenom. For this reason, people who are allergic to horses are more likely to suffer an allergic reaction to antivenom.[98] Antivenom for the more dangerous species (such as mambas, taipans, and cobras) is made in a similar manner in India, South Africa, and Australia, although these antivenoms are species-specific.
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In some parts of the world, especially in India, snake charming is a roadside show performed by a charmer. In such a show, the snake charmer carries a basket that contains a snake that he seemingly charms by playing tunes from his flutelike musical instrument, to which the snake responds.[99] Snakes lack external ears, though they do have internal ears, and respond to the movement of the flute, not the actual noise.[99]
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The Wildlife Protection Act of 1972 in India technically proscribes snake charming on grounds of reducing animal cruelty. Other snake charmers also have a snake and mongoose show, where both the animals have a mock fight; however, this is not very common, as the snakes, as well as the mongooses, may be seriously injured or killed. Snake charming as a profession is dying out in India because of competition from modern forms of entertainment and environment laws proscribing the practice. Many Indians have never seen snake charming and it is becoming a folktale of the past.[99][100][101][102]
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The Irulas tribe of Andhra Pradesh and Tamil Nadu in India have been hunter-gatherers in the hot, dry plains forests, and have practiced the art of snake catching for generations. They have a vast knowledge of snakes in the field. They generally catch the snakes with the help of a simple stick. Earlier, the Irulas caught thousands of snakes for the snake-skin industry. After the complete ban of the snake-skin industry in India and protection of all snakes under the Indian Wildlife (Protection) Act 1972, they formed the Irula Snake Catcher's Cooperative and switched to catching snakes for removal of venom, releasing them in the wild after four extractions. The venom so collected is used for producing life-saving antivenom, biomedical research and for other medicinal products.[103] The Irulas are also known to eat some of the snakes they catch and are very useful in rat extermination in the villages.
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Despite the existence of snake charmers, there have also been professional snake catchers or wranglers. Modern-day snake trapping involves a herpetologist using a long stick with a V- shaped end. Some television show hosts, like Bill Haast, Austin Stevens, Steve Irwin, and Jeff Corwin, prefer to catch them using bare hands.
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While not commonly thought of as food in most cultures, in others the consumption of snakes is acceptable, or even considered a delicacy. Snake soup of Cantonese cuisine is consumed by locals in autumn, to warm up their body. Western cultures document the consumption of snakes under extreme circumstances of hunger.[104] Cooked rattlesnake meat is an exception, which is commonly consumed in Texas[105] and parts of the Midwestern United States. In Asian countries such as China, Taiwan, Thailand, Indonesia, Vietnam and Cambodia, drinking the blood of snakes—particularly the cobra—is believed to increase sexual virility.[106] The blood is drained while the cobra is still alive when possible, and is usually mixed with some form of liquor to improve the taste.[106]
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In some Asian countries, the use of snakes in alcohol is also accepted. In such cases, the body of a snake or several snakes is left to steep in a jar or container of liquor. It is claimed that this makes the liquor stronger (as well as more expensive). One example of this is the Habu snake sometimes placed in the Okinawan liquor Habushu (ブ酒,) also known as "Habu Sake".[107]
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Snake wine (蛇酒) is an alcoholic beverage produced by infusing whole snakes in rice wine or grain alcohol. The drink was first recorded to have been consumed in China during the Western Zhou dynasty and considered an important curative and believed to reinvigorate a person according to traditional Chinese medicine.[108]
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In the Western world, some snakes (especially docile species such as the ball python and corn snake) are kept as pets. To meet this demand a captive breeding industry has developed. Snakes bred in captivity tend to make better pets and are considered preferable to wild caught specimens.[109] Snakes can be very low maintenance pets, especially compared to more traditional species. They require minimal space, as most common species do not exceed 5 feet (1.5 m) in length. Pet snakes can be fed relatively infrequently, usually once every 5 to 14 days. Certain snakes have a lifespan of more than 40 years if given proper care.
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In ancient Mesopotamia, Nirah, the messenger god of Ištaran, was represented as a serpent on kudurrus, or boundary stones.[110] Representations of two intertwined serpents are common in Sumerian art and Neo-Sumerian artwork[110] and still appear sporadically on cylinder seals and amulets until as late as the thirteenth century BC.[110] The horned viper (Cerastes cerastes) appears in Kassite and Neo-Assyrian kudurrus[110] and is invoked in Assyrian texts as a magical protective entity.[110] A dragon-like creature with horns, the body and neck of a snake, the forelegs of a lion, and the hind-legs of a bird appears in Mesopotamian art from the Akkadian Period until the Hellenistic Period (323 BC–31 BC).[110] This creature, known in Akkadian as the mušḫuššu, meaning "furious serpent", was used as a symbol for particular deities and also as a general protective emblem.[110] It seems to have originally been the attendant of the Underworld god Ninazu,[110] but later became the attendant to the Hurrian storm-god Tishpak, as well as, later, Ninazu's son Ningishzida, the Babylonian national god Marduk, the scribal god Nabu, and the Assyrian national god Ashur.[110]
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In Egyptian history, the snake occupies a primary role with the Nile cobra adorning the crown of the pharaoh in ancient times. It was worshipped as one of the gods and was also used for sinister purposes: murder of an adversary and ritual suicide (Cleopatra).[citation needed] The ouroboros was a well-known ancient Egyptian symbol of a serpent swallowing its own tail.[111] The precursor to the ouroboros was the "Many-Faced",[111] a serpent with five heads, who, according to the Amduat, the oldest surviving Book of the Afterlife, was said to coil around the corpse of the sun god Ra protectively.[111] The earliest surviving depiction of a "true" ouroboros comes from the gilded shrines in the tomb of Tutankhamun.[111] In the early centuries AD, the ouroboros was adopted as a symbol by Gnostic Christians[111] and chapter 136 of the Pistis Sophia, an early Gnostic text, describes "a great dragon whose tail is in its mouth".[111] In medieval alchemy, the ouroboros became a typical western dragon with wings, legs, and a tail.[111]
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In the Bible, King Nahash of Ammon, whose name means "Snake", is depicted very negatively, as a particularly cruel and despicable enemy of the ancient Hebrews.
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The ancient Greeks used the Gorgoneion, a depiction of a hideous face with serpents for hair, as an apotropaic symbol to ward off evil.[112] In a Greek myth described by Pseudo-Apollodorus in his Bibliotheca, Medusa was a Gorgon with serpents for hair whose gaze turned all those who looked at her to stone and was slain by the hero Perseus.[113][114][115] In the Roman poet Ovid's Metamorphoses, Medusa is said to have once been a beautiful priestess of Athena, whom Athena turned into a serpent-haired monster after she was raped by the god Poseidon in Athena's temple.[116] In another myth referenced by the Boeotian poet Hesiod and described in detail by Pseudo-Apollodorus, the hero Heracles is said to have slain the Lernaean Hydra,[117][118] a multiple-headed serpent which dwelt in the swamps of Lerna.[117][118]
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The legendary account of the foundation of Thebes mentioned a monster snake guarding the spring from which the new settlement was to draw its water. In fighting and killing the snake, the companions of the founder Cadmus all perished – leading to the term "Cadmean victory" (i.e. a victory involving one's own ruin).[citation needed]
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Three medical symbols involving snakes that are still used today are Bowl of Hygieia, symbolizing pharmacy, and the Caduceus and Rod of Asclepius, which are symbols denoting medicine in general.[49]
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One of the etymologies proposed for the common female first name Linda is that it might derive from Old German Lindi or Linda, meaning a serpent.
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India is often called the land of snakes and is steeped in tradition regarding snakes.[119] Snakes are worshipped as gods even today with many women pouring milk on snake pits (despite snakes' aversion for milk).[119] The cobra is seen on the neck of Shiva and Vishnu is depicted often as sleeping on a seven-headed snake or within the coils of a serpent.[120] There are also several temples in India solely for cobras sometimes called Nagraj (King of Snakes) and it is believed that snakes are symbols of fertility. There is a Hindu festival called Nag Panchami each year on which day snakes are venerated and prayed to. See also Nāga.[citation needed]
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In India there is another mythology about snakes. Commonly known in Hindi as "Ichchhadhari" snakes. Such snakes can take the form of any living creature, but prefer human form. These mythical snakes possess a valuable gem called "Mani", which is more brilliant than diamond. There are many stories in India about greedy people trying to possess this gem and ending up getting killed.[citation needed]
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The snake is one of the 12 celestial animals of Chinese zodiac, in the Chinese calendar.[121]
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Many ancient Peruvian cultures worshipped nature.[122] They emphasized animals and often depicted snakes in their art.[123]
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Snakes are a part of Hindu worship. A festival, Nag Panchami, in which participants worship either images of or live Nāgas (cobras) is celebrated every year. Most images of Lord Shiva depict snake around his neck. Puranas have various stories associated with snakes. In the Puranas, Shesha is said to hold all the planets of the Universe on his hoods and to constantly sing the glories of Vishnu from all his mouths. He is sometimes referred to as "Ananta-Shesha", which means "Endless Shesha". Other notable snakes in Hinduism are Ananta, Vasuki, Taxak, Karkotaka and Pingala. The term Nāga is used to refer to entities that take the form of large snakes in Hinduism and Buddhism.
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Snakes have also been widely revered, such as in ancient Greece, where the serpent was seen as a healer. Asclepius carried a serpent wound around his wand, a symbol seen today on many ambulances.
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In religious terms, the snake and jaguar are arguably the most important animals in ancient Mesoamerica. "In states of ecstasy, lords dance a serpent dance; great descending snakes adorn and support buildings from Chichen Itza to Tenochtitlan, and the Nahuatl word coatl meaning serpent or twin, forms part of primary deities such as Mixcoatl, Quetzalcoatl, and Coatlicue."[124] In both Maya and Aztec calendars, the fifth day of the week was known as Snake Day.
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In Judaism, the snake of brass is also a symbol of healing, of one's life being saved from imminent death.[125]
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In some parts of Christianity, Christ's redemptive work is compared to saving one's life through beholding the Nehushtan (serpent of brass).[126] Snake handlers use snakes as an integral part of church worship in order to exhibit their faith in divine protection. However, more commonly in Christianity, the serpent has been seen as a representative of evil and sly plotting, which can be seen in the description in Genesis chapter 3 of a snake in the Garden of Eden tempting Eve.[127] Saint Patrick is reputed to have expelled all snakes from Ireland while converting the country to Christianity in the 5th century, thus explaining the absence of snakes there.
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In Christianity and Judaism, the snake makes its infamous appearance in the first book of the Bible when a serpent appears before the first couple Adam and Eve and tempts them with the forbidden fruit from the Tree of Knowledge.[127] The snake returns in Exodus when Moses, as a sign of God's power, turns his staff into a snake and when Moses made the Nehushtan, a bronze snake on a pole that when looked at cured the people of bites from the snakes that plagued them in the desert. The serpent makes its final appearance symbolizing Satan in the Book of Revelation: "And he laid hold on the dragon the old serpent, which is the devil and Satan, and bound him for a thousand years."[128]
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In Neo-Paganism and Wicca, the snake is seen as a symbol of wisdom and knowledge.
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Several compounds from snake venoms are being researched as potential treatments or preventatives for pain, cancers, arthritis, stroke, heart disease, hemophilia, and hypertension, and to control bleeding (e.g. during surgery).[130][131][132]
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Caldwell MW, Nydam RL, Palci A, Apesteguía S (January 2015). "The oldest known snakes from the Middle Jurassic-Lower Cretaceous provide insights on snake evolution". Nature Communications. 6 (5996): 5996. Bibcode:2015NatCo...6.5996C. doi:10.1038/ncomms6996. PMID 25625704.
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Orange usually refers to:
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Orange may also refer to:
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Athens (/ˈæθɪnz/ ATH-inz;[3] Greek: Αθήνα, romanized: Athína [aˈθina] (listen); Ancient Greek: Ἀθῆναι, romanized: Athênai (pl.) [atʰɛ̂ːnai̯]), is the capital and largest city of Greece. Athens dominates the Attica region and is one of the world's oldest cities, with its recorded history spanning over 3,400 years[citation needed] and its earliest human presence started somewhere between the 11th and 7th millennium BC.[4]
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Classical Athens was a powerful city-state that emerged in conjunction with the seagoing development of the port of Piraeus. A center for the arts, learning and philosophy, home of Plato's Academy and Aristotle's Lyceum,[5][6] it is widely referred to as the cradle of Western civilization and the birthplace of democracy,[7][8] largely because of its cultural and political impact on the European continent, and in particular the Romans.[9] In modern times, Athens is a large cosmopolitan metropolis and central to economic, financial, industrial, maritime, political and cultural life in Greece.
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Athens is a global city and one of the biggest economic centers in southeastern Europe. It has a large financial sector, and its port Piraeus is both the largest passenger port in Europe,[10][11] and the second largest in the world.[12]
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The Municipality of Athens (also City of Athens), which actually constitutes a small administrative unit of the entire city, had a population of 664,046 (in 2011)[2] within its official limits, and a land area of 38.96 km2 (15.04 sq mi).[13][14] The Athens Urban Area (Greater Athens and Greater Piraeus) extends beyond its administrative municipal city limits, with a population of 3,090,508 (in 2011)[15] over an area of 412 km2 (159 sq mi).[14] According to Eurostat[16] in 2011, the functional urban area (FUA) of Athens was the 9th most populous FUA in the European Union (the 6th most populous capital city of the EU), with a population of 3.8 million people. Athens is also the southernmost capital on the European mainland and the warmest major city in Europe.
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The heritage of the classical era is still evident in the city, represented by ancient monuments and works of art, the most famous of all being the Parthenon, considered a key landmark of early Western civilization. The city also retains Roman and Byzantine monuments, as well as a smaller number of Ottoman monuments. Athens is home to two UNESCO World Heritage Sites, the Acropolis of Athens and the medieval Daphni Monastery. Landmarks of the modern era, dating back to the establishment of Athens as the capital of the independent Greek state in 1834, include the Hellenic Parliament and the so-called "architectural trilogy of Athens", consisting of the National Library of Greece, the National and Kapodistrian University of Athens and the Academy of Athens. Athens is also home to several museums and cultural institutions, such as the National Archeological Museum, featuring the world's largest collection of ancient Greek antiquities, the Acropolis Museum, the Museum of Cycladic Art, the Benaki Museum and the Byzantine and Christian Museum. Athens was the host city of the first modern-day Olympic Games in 1896, and 108 years later it hosted the 2004 Summer Olympics, making it one of the few cities to have hosted the Olympics more than once.[17]
|
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|
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In Ancient Greek, the name of the city was Ἀθῆναι (Athênai, pronounced [atʰɛ̂ːnai̯] in Classical Attic) a plural. In earlier Greek, such as Homeric Greek, the name had been current in the singular form though, as Ἀθήνη (Athḗnē).[18] It was possibly rendered in the plural later on, like those of Θῆβαι (Thêbai) and Μυκῆναι (Μukênai). The root of the word is probably not of Greek or Indo-European origin,[19] and is possibly a remnant of the Pre-Greek substrate of Attica.[19] In antiquity, it was debated whether Athens took its name from its patron goddess Athena (Attic Ἀθηνᾶ, Athēnâ, Ionic Ἀθήνη, Athḗnē, and Doric Ἀθάνα, Athā́nā) or Athena took her name from the city.[20] Modern scholars now generally agree that the goddess takes her name from the city,[20] because the ending -ene is common in names of locations, but rare for personal names.[20] During the medieval period, the name of the city was rendered once again in the singular as Ἀθήνα. However, after the establishment of the modern Greek state, and partly due to the conservatism of the written language, Ἀθῆναι [aˈθine] became again the official name of the city and remained so until the abandonment of Katharevousa in the 1970s, when Ἀθήνα, Athína, became the official name.
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According to the ancient Athenian founding myth, Athena, the goddess of wisdom, competed against Poseidon, the God of the Seas, for patronage of the yet-unnamed city;[21] they agreed that whoever gave the Athenians the better gift would become their patron[21] and appointed Cecrops, the king of Athens, as the judge.[21] According to the account given by Pseudo-Apollodorus, Poseidon struck the ground with his trident and a salt water spring welled up.[21] In an alternative version of the myth from Vergil's Georgics, Poseidon instead gave the Athenians the first horse.[21] In both versions, Athena offered the Athenians the first domesticated olive tree.[21][22] Cecrops accepted this gift[21] and declared Athena the patron goddess of Athens.[21][22] Eight different etymologies, now commonly rejected, were proposed during the 17th century.[citation needed] Christian Lobeck proposed as the root of the name the word ἄθος (áthos) or ἄνθος (ánthos) meaning "flower", to denote Athens as the "flowering city". Ludwig von Döderlein proposed the stem of the verb θάω, stem θη- (tháō, thē-, "to suck") to denote Athens as having fertile soil.[23]
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|
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+
In classical literature, the city was sometimes referred to as the City of the Violet Crown, first documented in Pindar's ἰοστέφανοι Ἀθᾶναι (iostéphanoi Athânai), or as τὸ κλεινὸν ἄστυ (tò kleinòn ásty, "the glorious city"). In medieval texts, variant names include Setines, Satine, and Astines, all derivations involving false splitting of prepositional phrases.[24] Today the caption η πρωτεύουσα (ī protévousa), "the capital", has become somewhat common. King Alphonse X of Castile in his 'General Storia' (Global History) cites etymology of Athens as: 'Athinai', 'the one without death', death here meaning ignorance.
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|
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+
Athenians were called cicada-wearers (Ancient Greek: Τεττιγοφόροι) because they used to wear pins of golden cicadas. A symbol of being autochthon (earth-born), because the legendary founder of Athens, Erechtheus was an autochthon or of being musicians, because the cicada is a "musician" insect.[25]
|
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The oldest known human presence in Athens is the Cave of Schist, which has been dated to between the 11th and 7th millennia BC.[4] Athens has been continuously inhabited for at least 5,000 years.[26][27] By 1400 BC the settlement had become an important centre of the Mycenaean civilization and the Acropolis was the site of a major Mycenaean fortress, whose remains can be recognised from sections of the characteristic Cyclopean walls.[28] Unlike other Mycenaean centers, such as Mycenae and Pylos, it is not known whether Athens suffered destruction in about 1200 BC, an event often attributed to a Dorian invasion, and the Athenians always maintained that they were pure Ionians with no Dorian element. However, Athens, like many other Bronze Age settlements, went into economic decline for around 150 years afterwards.
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|
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Iron Age burials, in the Kerameikos and other locations, are often richly provided for and demonstrate that from 900 BC onwards Athens was one of the leading centres of trade and prosperity in the region.[29] The leading position of Athens may well have resulted from its central location in the Greek world, its secure stronghold on the Acropolis and its access to the sea, which gave it a natural advantage over inland rivals such as Thebes and Sparta.
|
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By the 6th century BC, widespread social unrest led to the reforms of Solon. These would pave the way for the eventual introduction of democracy by Cleisthenes in 508 BC. Athens had by this time become a significant naval power with a large fleet, and helped the rebellion of the Ionian cities against Persian rule. In the ensuing Greco-Persian Wars Athens, together with Sparta, led the coalition of Greek states that would eventually repel the Persians, defeating them decisively at Marathon in 490 BC, and crucially at Salamis in 480 BC. However, this did not prevent Athens from being captured and sacked twice by the Persians within one year, after a heroic but ultimately failed resistance at Thermopylae by Spartans and other Greeks led by King Leonidas,[30] after both Boeotia and Attica fell to the Persians.
|
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The decades that followed became known as the Golden Age of Athenian democracy, during which time Athens became the leading city of Ancient Greece, with its cultural achievements laying the foundations for Western civilization. The playwrights Aeschylus, Sophocles and Euripides flourished in Athens during this time, as did the historians Herodotus and Thucydides, the physician Hippocrates, and the philosopher Socrates. Guided by Pericles, who promoted the arts and fostered democracy, Athens embarked on an ambitious building program that saw the construction of the Acropolis of Athens (including the Parthenon), as well as empire-building via the Delian League. Originally intended as an association of Greek city-states to continue the fight against the Persians, the league soon turned into a vehicle for Athens's own imperial ambitions. The resulting tensions brought about the Peloponnesian War (431–404 BC), in which Athens was defeated by its rival Sparta.
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By the mid-4th century BC, the northern Greek kingdom of Macedon was becoming dominant in Athenian affairs. In 338 BC the armies of Philip II defeated an alliance of some of the Greek city-states including Athens and Thebes at the Battle of Chaeronea, effectively ending Athenian independence. Later, under Rome, Athens was given the status of a free city because of its widely admired schools. The Roman emperor Hadrian, in the 2nd century CE, ordered the construction of a library, a gymnasium, an aqueduct which is still in use, several temples and sanctuaries, a bridge and financed the completion of the Temple of Olympian Zeus.
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By the end of Late Antiquity, Athens had shrunk due to sacks by the Herulians, Visigoths, and Early Slavs which caused massive destruction in the city. In this era, the first Christian churches were built in Athens, and the Parthenon and other temples were converted into churches. Athens expanded its settlement in the second half of the Middle Byzantine Period, in the 9th to 10th centuries CE, and was relatively prosperous during the Crusades, benefiting from Italian trade. After the Fourth Crusade the Duchy of Athens was established. In 1458 it was conquered by the Ottoman Empire and entered a long period of decline.
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Following the Greek War of Independence and the establishment of the Greek Kingdom, Athens was chosen as the capital of the newly independent Greek state in 1834, largely because of historical and sentimental reasons. At the time, it was reduced to a town of about 4,000 people in a loose swarm of houses along the foot of the Acropolis. The first King of Greece, Otto of Bavaria, commissioned the architects Stamatios Kleanthis and Eduard Schaubert to design a modern city plan fit for the capital of a state.
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The first modern city plan consisted of a triangle defined by the Acropolis, the ancient cemetery of Kerameikos and the new palace of the Bavarian king (now housing the Greek Parliament), so as to highlight the continuity between modern and ancient Athens. Neoclassicism, the international style of this epoch, was the architectural style through which Bavarian, French and Greek architects such as Hansen, Klenze, Boulanger or Kaftantzoglou designed the first important public buildings of the new capital. In 1896, Athens hosted the first modern Olympic Games. During the 1920s a number of Greek refugees, expelled from Asia Minor after the Greco-Turkish War, swelled Athens's population; nevertheless it was most particularly following World War II, and from the 1950s and 1960s, that the population of the city exploded, and Athens experienced a gradual expansion.
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In the 1980s it became evident that smog from factories and an ever-increasing fleet of automobiles, as well as a lack of adequate free space due to congestion, had evolved into the city's most important challenge. A series of anti-pollution measures taken by the city's authorities in the 1990s, combined with a substantial improvement of the city's infrastructure (including the Attiki Odos motorway, the expansion of the Athens Metro, and the new Athens International Airport), considerably alleviated pollution and transformed Athens into a much more functional city. In 2004 Athens hosted the 2004 Summer Olympics.
|
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Tondo of the Aison Cup, showing the victory of Theseus over the Minotaur in the presence of Athena. Theseus was responsible, according to the myth, for the synoikismos ("dwelling together")—the political unification of Attica under Athens.
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The earliest coinage of Athens, c. 545–525/15 BC
|
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Coat of Arms of the Duchy of Athens during the rule of the de la Roche family (13th century)
|
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The Roman Agora and the Gate of Athena in Plaka district.
|
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The Temple of Olympian Zeus with river Ilisos by Edward Dodwell, 1821
|
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The Entry of King Otto in Athens, Peter von Hess, 1839.
|
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The Stadiou Street in Central Athens in 1908.
|
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Athens sprawls across the central plain of Attica that is often referred to as the Athens Basin or the Attica Basin (Greek: Λεκανοπέδιο Αθηνών/Αττικής). The basin is bounded by four large mountains: Mount Aigaleo to the west, Mount Parnitha to the north, Mount Pentelicus to the northeast and Mount Hymettus to the east.[31] Beyond Mount Aegaleo lies the Thriasian plain, which forms an extension of the central plain to the west. The Saronic Gulf lies to the southwest. Mount Parnitha is the tallest of the four mountains (1,413 m (4,636 ft)),[32] and has been declared a national park.
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Athens is built around a number of hills. Lycabettus is one of the tallest hills of the city proper and provides a view of the entire Attica Basin. The meteorology of Athens is deemed to be one of the most complex in the world because its mountains cause a temperature inversion phenomenon which, along with the Greek Government's difficulties controlling industrial pollution, was responsible for the air pollution problems the city has faced.[27] This issue is not unique to Athens; for instance, Los Angeles and Mexico City also suffer from similar atmospheric inversion problems.[27]
|
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The Cephissus river, the Ilisos and the Eridanos stream are the historical rivers of Athens.
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By the late 1970s, the pollution of Athens had become so destructive that according to the then Greek Minister of Culture, Constantine Trypanis, "...the carved details on the five the caryatids of the Erechtheum had seriously degenerated, while the face of the horseman on the Parthenon's west side was all but obliterated."[33] A series of measures taken by the authorities of the city throughout the 1990s resulted in the improvement of air quality; the appearance of smog (or nefos as the Athenians used to call it) has become less common.
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Measures taken by the Greek authorities throughout the 1990s have improved the quality of air over the Attica Basin. Nevertheless, air pollution still remains an issue for Athens, particularly during the hottest summer days. In late June 2007,[34] the Attica region experienced a number of brush fires,[34] including a blaze that burned a significant portion of a large forested national park in Mount Parnitha,[35] considered critical to maintaining a better air quality in Athens all year round.[34] Damage to the park has led to worries over a stalling in the improvement of air quality in the city.[34]
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The major waste management efforts undertaken in the last decade (particularly the plant built on the small island of Psytalia) have greatly improved water quality in the Saronic Gulf, and the coastal waters of Athens are now accessible again to swimmers. In January 2007, Athens faced a waste management problem when its landfill near Ano Liosia, an Athenian suburb, reached capacity.[36] The crisis eased by mid-January when authorities began taking the garbage to a temporary landfill.[36]
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Athens ranks in the lowest percentage for the risk on frequency and severity of terrorist attacks according to the EU Global Terrorism Database (EIU 2007-2016 calculations). The city also ranked 35th in Digital Security, 21st on Health Security, 29th on Infrastructure Security and 41st on Personal Security globally in a 2017 The Economist Intelligence Unit report.[37] It also ranks as a very safe city (39th globally out of 162 cities overall) on the ranking of the safest and most dangerous countries.[38] A 2019 crime index from Numbeo places Athens at 130th position, rating safer than Tampa, Florida or Dublin, Ireland.[39] According to a Mercer 2019 Quality of Living Survey, Athens ranks 89th on the Mercer Quality of Living Survey ranking.[40]
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Athens has a hot-summer Mediterranean climate (Köppen climate classification: Csa). The dominant feature of Athens' climate is alternation between prolonged hot and dry summers and mild winters with moderate rainfall.[41] With an average of 455 millimetres (17.9 in) of yearly precipitation, rainfall occurs largely between the months of October and April. July and August are the driest months when thunderstorms occur sparsely. Furthermore, some coastal areas of Athens, known as the Athens Riviera have a hot semi-arid climate (BSh) according to the climate atlas published by the Hellenic National Meteorological Service (H.N.M.S).[42]
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Owing to the rain shadow of the Pindus Mountains, annual precipitation of Athens is lower than most other parts of Greece, especially western Greece. As an example, Ioannina receives around 1,300 mm (51 in) per year, and Agrinio around 800 mm (31 in) per year. Daily average highs for July have been measured around 34 °C or 93 °F in downtown Athens, but some parts of the city may be even hotter for the higher density of buildings, such as the center,[43] in particular, western areas due to a combination of industrialization and a number of natural factors, knowledge of which has existed since the mid-19th century.[44][45][46]
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Athens is affected by the urban heat island effect in some areas which is caused by human activity,[47][48] altering its temperatures compared to the surrounding rural areas,[49][50][51][52] and leaving detrimental effects on energy usage, expenditure for cooling,[53][54] and health.[48] The urban heat island of the city has also been found to be partially responsible for alterations of the climatological temperature time-series of specific Athens meteorological stations, because of its impact on the temperatures and the temperature trends recorded by some meteorological stations.[55][56][57][58][59] On the other hand, specific meteorological stations, such as the National Garden station and Thiseio meteorological station, are less affected or do not experience the urban heat island.[49][60]
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Athens holds the World Meteorological Organization record for the highest temperature ever recorded in Europe, at 48 °C (118.4 °F), which was recorded in the Elefsina and Tatoi suburbs of Athens on 10 July 1977.[61]
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The Municipality of Athens, the City Centre of the Athens Urban Area, is divided into several districts: Omonoia, Syntagma, Exarcheia, Agios Nikolaos, Neapolis, Lykavittos, Lofos Strefi, Lofos Finopoulou, Lofos Filopappou, Pedion Areos, Metaxourgeio, Aghios Kostantinos, Larissa Station, Kerameikos, Psiri, Monastiraki, Gazi, Thission, Kapnikarea, Aghia Irini, Aerides, Anafiotika, Plaka, Acropolis, Pnyka, Makrygianni, Lofos Ardittou, Zappeion, Aghios Spyridon, Pangrati, Kolonaki, Dexameni, Evaggelismos, Gouva, Aghios Ioannis, Neos Kosmos, Koukaki, Kynosargous, Fix, Ano Petralona, Kato Petralona, Rouf, Votanikos, Profitis Daniil, Akadimia Platonos, Kolonos, Kolokynthou, Attikis Square, Lofos Skouze, Sepolia, Kypseli, Aghios Meletios, Nea Kypseli, Gyzi, Polygono, Ampelokipoi, Panormou-Gerokomeio, Pentagono, Ellinorosson, Nea Filothei, Ano Kypseli, Tourkovounia-Lofos Patatsou, Lofos Elikonos, Koliatsou, Thymarakia, Kato Patisia, Treis Gefyres, Aghios Eleftherios, Ano Patisia, Kypriadou, Menidi, Prompona, Aghios Panteleimonas, Pangrati, Goudi, Vyronas and Ilisia.
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Parnitha National Park is punctuated by well-marked paths, gorges, springs, torrents and caves dotting the protected area. Hiking and mountain-biking in all four mountains are popular outdoor activities for residents of the city. The National Garden of Athens was completed in 1840 and is a green refuge of 15.5 hectares in the centre of the Greek capital. It is to be found between the Parliament and Zappeion buildings, the latter of which maintains its own garden of seven hectares.
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Parts of the City Centre have been redeveloped under a masterplan called the Unification of Archeological Sites of Athens, which has also gathered funding from the EU to help enhance the project.[68][69] The landmark Dionysiou Areopagitou Street has been pedestrianised, forming a scenic route. The route starts from the Temple of Olympian Zeus at Vasilissis Olgas Avenue, continues under the southern slopes of the Acropolis near Plaka, and finishes just beyond the Temple of Hephaestus in Thiseio. The route in its entirety provides visitors with views of the Parthenon and the Agora (the meeting point of ancient Athenians), away from the busy City Centre.
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The hills of Athens also provide green space. Lycabettus, Philopappos hill and the area around it, including Pnyx and Ardettos hill, are planted with pines and other trees, with the character of a small forest rather than typical metropolitan parkland. Also to be found is the Pedion tou Areos (Field of Mars) of 27.7 hectares, near the National Archaeological Museum.
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Athens' largest zoo is the Attica Zoological Park, a 20-hectare (49-acre) private zoo located in the suburb of Spata. The zoo is home to around 2000 animals representing 400 species, and is open 365 days a year. Smaller zoos exist within public gardens or parks, such as the zoo within the National Garden of Athens.
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The Athens Metropolitan Area consists of 58[15] densely populated municipalities, sprawling around the Municipality of Athens (the City Centre) in virtually all directions. For the Athenians, all the urban municipalities surrounding the City Centre are called suburbs. According to their geographic location in relation to the City of Athens, the suburbs are divided into four zones; the northern suburbs (including Agios Stefanos, Dionysos, Ekali, Nea Erythraia, Kifissia, Kryoneri, Maroussi, Pefki, Lykovrysi, Metamorfosi, Nea Ionia, Nea Filadelfeia, Irakleio, Vrilissia, Melissia, Penteli, Chalandri, Agia Paraskevi, Gerakas, Pallini, Galatsi, Psychiko and Filothei); the southern suburbs (including Alimos, Nea Smyrni, Moschato, Tavros, Agios Ioannis Rentis, Kallithea, Piraeus, Agios Dimitrios, Palaio Faliro, Elliniko, Glyfada, Lagonisi, Saronida, Argyroupoli, Ilioupoli, Varkiza, Voula, Vari and Vouliagmeni); the eastern suburbs (including Zografou, Dafni, Vyronas, Kaisariani, Cholargos and Papagou); and the western suburbs (including Peristeri, Ilion, Egaleo, Koridallos, Agia Varvara, Keratsini, Perama, Nikaia, Drapetsona, Chaidari, Petroupoli, Agioi Anargyroi, Ano Liosia, Aspropyrgos, Eleusina, Acharnes and Kamatero).
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The Athens city coastline, extending from the major commercial port of Piraeus to the southernmost suburb of Varkiza for some 25 km (20 mi),[70] is also connected to the City Centre by tram.
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In the northern suburb of Maroussi, the upgraded main Olympic Complex (known by its Greek acronym OAKA) dominates the skyline. The area has been redeveloped according to a design by the Spanish architect Santiago Calatrava, with steel arches, landscaped gardens, fountains, futuristic glass, and a landmark new blue glass roof which was added to the main stadium. A second Olympic complex, next to the sea at the beach of Palaio Faliro, also features modern stadia, shops and an elevated esplanade. Work is underway to transform the grounds of the old Athens Airport – named Elliniko – in the southern suburbs, into one of the largest landscaped parks in Europe, to be named the Hellenikon Metropolitan Park.[71]
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Many of the southern suburbs (such as Alimos, Palaio Faliro, Elliniko, Glyfada, Voula, Vouliagmeni and Varkiza) known as the Athens Riviera, host a number of sandy beaches, most of which are operated by the Greek National Tourism Organisation and require an entrance fee. Casinos operate on both Mount Parnitha, some 25 km (16 mi)[72] from downtown Athens (accessible by car or cable car), and the nearby town of Loutraki (accessible by car via the Athens – Corinth National Highway, or the suburban rail service Proastiakos).
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The large City Centre (Greek: Κέντρο της Αθήνας) of the Greek capital falls directly within the Municipality of Athens or Athens Municipality (Greek: Δήμος Αθηναίων)—also City of Athens. Athens Municipality is the largest in population size in Greece. Piraeus also forms a significant city centre on its own,[73] within the Athens Urban Area and it is the second largest in population size within it, with Peristeri, Kallithea and Kypseli following.
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The Athens Urban Area (Greek: Πολεοδομικό Συγκρότημα Αθηνών), also known as Urban Area of the Capital (Greek: Πολεοδομικό Συγκρότημα Πρωτεύουσας), today consists of 40 municipalities, 35 of which make up what is referred to as the Greater Athens municipalities, located within 4 regional units (North Athens, West Athens, Central Athens, South Athens); and a further 5, which make up the Greater Piraeus municipalities, located within the regional unit of Piraeus as mentioned above. The densely built up urban area of the Greek capital sprawls across 412 km2 (159 sq mi)[14] throughout the Attica Basin and has a total population of 3,074,160 (in 2011).
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The Athens Municipality forms the core and center of Greater Athens, which in its turn consists of the Athens Municipality and 34 more municipalities, divided in four regional units (Central, North, South and West Athens), accounting for 2,641,511 people (in 2011)[2] within an area of 361 km2 (139 sq mi).[14] Until 2010, these four regional units made up the abolished Athens Prefecture. The municipality of Piraeus, the historic Athenian port, with its 4 suburban municipalities make up the regional unit of Piraeus, which in turn forms Greater Piraeus.
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Greater Athens and Greater Piraeus with part of East[74] and West Attica[75] regional units combined make up the continuous Athens Urban Area,[75][76][77] also called the "Urban Area of the Capital" or simply "Athens" (the most common use of the term), spanning over 412 km2 (159 sq mi),[78] with a population of 3,090,508 people as of 2011. The Athens Urban Area is considered to form the city of Athens as a whole, despite its administrative divisions, which is the largest in Greece and one of the most populated urban areas in Europe.
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The Athens Metropolitan Area (Greek: Μητροπολιτική Περιοχή της Αθήνας), also commonly known in English as Greater Athens,[79] spans 2,928.717 km2 (1,131 sq mi) within the Attica region and includes a total of 58 municipalities, which are organized in seven regional units (those outlined above, along with East Attica and West Attica), having reached a population of 3,737,550 based on the preliminary results of the 2011 census. Athens and Piraeus municipalities serve as the two metropolitan centres of the Athens Metropolitan Area.[80] There are also some inter-municipal centres serving specific areas. For example, Kifissia and Glyfada serve as inter-municipal centres for northern and southern suburbs respectively.
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The Municipality of Athens has an official population of 664,046 people.[2] The four regional units that make up what is referred to as Greater Athens have a combined population of 2,640,701. They together with the regional unit of Piraeus (Greater Piraeus) make up the dense Athens Urban Area which reaches a total population of 3,090,508 inhabitants (in 2011).[15] As Eurostat the FUA of Athens had in 2013 3,828,434 inhabitants, being apparently decreasing compared with the pre-economic crisis date of 2009 (4,164,175)[16]
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The municipality (Center) of Athens is the most populous in Greece, with a population of 664,046 people (in 2011)[2] and an area of 38.96 km2 (15.04 sq mi),[13] forming the core of the Athens Urban Area within the Attica Basin. The incumbent Mayor of Athens is Kostas Bakoyannis of New Democracy. The municipality is divided into seven municipal districts which are mainly used for administrative purposes.
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As of the 2011 census, the population for each of the seven municipal districts of Athens is as follows:[81]
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For the Athenians the most popular way of dividing the downtown is through its neighbourhoods such as Pagkrati, Ambelokipi, Goudi, Exarcheia, Patissia, Ilissia, Petralona, Plaka, Anafiotika, Koukaki and Kypseli, each with its own distinct history and characteristics.
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The Athens Metropolitan Area, with an area of 2,928.717 km2 (1,131 sq mi) and inhabited by 3,753,783 people in 2011,[2] consists of the Athens Urban Area with the addition of the towns and villages of East and West Attica, which surround the dense urban area of the Greek capital. It actually sprawls over the whole peninsula of Attica, which is the best part of the region of Attica, excluding the islands.
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Mycenean Athens in 1600–1100 BC could have reached the size of Tiryns; that would put the population at the range of 10,000–15,000.[82] During the Greek Dark Ages the population of Athens was around 4,000 people. In 700 BC the population grew to 10,000. In 500 BC the area probably contained 200,000 people. During the classical period the city's population is estimated from 150,000–350,000 and up to 610,000 according to Thucydides. When Demetrius of Phalerum conducted a population census in 317 BC the population was 21,000 free citizens, plus 10,000 resident aliens and 400,000 slaves. This suggests a total population of 431,000.[83][84][85][86][87][88][89][90][91] This figure is highly suspect because of the lopsided number of slaves and does not include free women and children and resident foreigners: an estimated based on Thucydides is: 40,000 male citizens, 100,000 family members, 70,000 metics (resident foreigners) and 150,000-400,000 slaves. However the numbers would include all of Athenian territory which consisted almost entirely of modern Attica region expect the territory of the city-state of Megaris and the island section.
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The ancient site of the main city is centred on the rocky hill of the acropolis. In the whole of Athenian territory they existed many towns. Acharnae, Afidnes, Cytherus, Colonus, Corydallus, Cropia, Decelea, Euonymos, Vravron among others was important towns in Athenian countryside. The new port of Piraeus was a prototype harbour with the infrastructure and housing located in the site between modern passenger section of the port (named Kantharos in ancient times) and the Pasalimani harbour (named Zea in ancient times). The old one Phaliro was in the site of modern Palaio Faliro and gradually declined after the construction of the new prototype port but remained as a minor port and important settlement with historic significance in late Classical times. The rapid expansion of the modern city, which continues to this day, was initiated in the 1950s and 1960s, because of Greece's transition from an agricultural to an industrial nation.[92] The expansion is now particularly toward the East and North East (a tendency greatly related to the new Eleftherios Venizelos International Airport and the Attiki Odos, the freeway that cuts across Attica). By this process Athens has engulfed many former suburbs and villages in Attica, and continues to do so. The table below shows the historical population of Athens in recent times.
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Athens became the capital of Greece in 1834, following Nafplion, which was the provisional capital from 1829. The municipality (City) of Athens is also the capital of the Attica region. The term Athens can refer either to the Municipality of Athens, to Greater Athens, or to the entire Athens Urban Area.
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The Hellenic Parliament
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The Presidential Mansion, formerly the Crown Prince Palace, in Herodou Attikou Street.
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The Maximos Mansion, official office of the Prime Minister of the Hellenic Republic, in Herodou Attikou Street.
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The Athens City Hall in Kotzia Square was designed by Panagiotis Kolkas and completed in 1874.[98]
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The Embassy of France in Vasilissis Sofias Avenue.
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The Italian Embassy in Vasilissis Sofias Avenue.
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Athens is twinned with:[99]
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Athens is the financial capital of Greece. According to data from 2014, Athens as a metropolitan economic area produced 130 billion US-dollars as GDP in PPP, which consists nearly a half of the production for the whole country. In the list with the strongest economic metropoles of the world Athens was ranked that year 102nd, while the GDP per capita for the same year was 32,000 US-dollars.[112]
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Athens is one of the major economic centres in south-eastern Europe and is considered as a regional economic power in Europe generally. The port of Piraeus, where big investments by COSCO have already been delivered during the recent decade, the completion of the new Cargo Centre in Thriasion,[113] the expansion of the Athens Metro and the Athens Tram, as well as the projected metropolitan park in Elliniko and other economic projects are the economic landmarks of the upcoming years.
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Important Greek companies such as Hellenic Aerospace Industry, Mytilineos Holdings, Titan Cement, Hellenic Petroleum, Papadopoulos E.J., Folli Follie, Jumbo S.A., OPAP, and COSMOTE have their headquarters in the metropolitan area of Athens. Multinational companies such as Ericsson, Sony, Siemens, Motorola, Samsung, Microsoft, Novartis, Mondelez, Coca-Cola, etc. have their regional research and development headquarters also there.
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The banking sector is represented by National Bank of Greece, Alpha Bank, Eurobank, and Piraeus Bank, while the Bank of Greece is also situated in the City Centre. The Athens Stock Exchange, the only in Greece, has been severely hit by the Greek government-debt crisis and the decision of the government to proceed into capital controls during summer 2015. As a whole the economy of Athens and Greece has been severely hit with today's data showing a change from long recession to growth of 1.4% in 2017.[114]
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Tourism is also a great contributor for the economy of the city, which is considered as one of the top destinations in Europe for city-break tourism and is also the gateway for excursions to the islands or the mainland. Greece attracted 26.5 million visitors in 2015, 30.1 million visitors in 2017 and over 33 million in 2018, making Greece one of the most visited countries in Europe and the world, and contributing 18% to the nation's Gross Domestic Product. Athens welcomed more than 5 million tourists in 2018 and 1,4 million of them were "city-breakers" (in 2013 the city-breakers were only 220.000).[115]
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Athens is the major transportation hub of Greece. They combines the largest Airport in Greece as also the largest port in Greece which also the largest port in Mediterranean in containers transport and the largest passenger port in Europe. Is also major national hub for Intercity (Ktel) and international buses as well as for domestic and international rail transport. Public transport is serviced by a variety of transportation means, forming the largest mass transit system of Greece. The Athens Mass Transit System consists of a large bus fleet, a trolleybus fleet that mainly serves Athens's city center, the city's Metro, a Proastiakos service[116] and a tram network, connecting the southern suburbs to the City Centre.[117]
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OSY (Greek: ΟΣΥ) (Odikes Sygkoinonies S.A.) which is subsidiary company of OASA (Athens urban transport organisation), is the main operator of buses and trolleybusses in Athens. Its network consists of about 300 bus lines and 22 trolleybus lines which span the Athens Metropolitan Area,[118] with a fleet of 1,839 buses and 366 trolleybuses.[119] Of those 1,839 buses 416 run on compressed natural gas,[119] making up the largest fleet of natural gas-powered buses in Europe[120] and all trolleybusses are equipped to enable them to run on diesel in case of power failure.[121]
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International links are provided by a number of private companies. National and regional bus links are provided by KTEL from two InterCity Bus Terminals, Kifissos Bus Terminal A and Liosion Bus Terminal B, both located in the north-western part of the city. Kifissos provides connections towards Peloponnese, North Greece, West Greece and some Ionian Islands, whereas Liosion is used for most of Central Greece.
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The Athens Metro is operated by STASY S.A (Greek: ΣΤΑΣΥ) (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation) and provides public transport throughout the Athens Urban Area. While its main purpose is transport, it also houses Greek artifacts found during construction of the system.[122] The Athens Metro has an operating staff of 387 and runs three metro lines, namely the line 1 (Green Line), line 2 (Red Line) and line 3 (Blue Line) lines, of which the first was constructed
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in 1869, and the other two largely during the 1990s, with the initial new sections opened in January 2000. Line 1 mostly runs at ground level and the other two (Line 2 & 3) routes run entirely underground. A fleet of 42 trains, using 252 carriages, operates on the network,[123] with a daily occupancy of 1,353,000 passengers.[124]
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The line 1 (Green Line) serves 24 stations, and is the oldest line of the Athens metro network. It runs from Piraeus station to Kifissia station and covers a distance of 25.6-kilometre (15.9 mi). There are transfer connections with the Blue Line 3 at Monastiraki station and with the Red Line 2 at Omonoia and Attiki stations.
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The line 2 (Red Line) runs from Anthoupoli station to Elliniko station and covers a distance of 17.5 km (10.9 mi).[123] The line connects the western suburbs of Athens with the southeast suburbs, passing through the center of Athens. The Red Line has transfer connections with the Green Line 1 at Attiki and Omonoia stations. There are also transfer connections with the Blue Line 3 at Syntagma station and with the tram at Syntagma, Syngrou Fix and Neos Kosmos stations.
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The line 3 (Blue Line) runs from Agia Marina and Egaleo stations, through the central Monastiraki and Syntagma stations to Doukissis Plakentias avenue in the northeastern suburb of Halandri, covering a distance of 16 km (10 mi).[123] It then ascends to ground level and continues to Eleftherios Venizelos International Airport, using the Suburban Railway infrastructure, extending its total length to 39 km (24 mi).[123] The spring 2007 extension from Monastiraki westwards to Egaleo connected some of the main night life hubs of the city, namely those of Gazi (Kerameikos station) with Psirri (Monastiraki station) and the City Centre (Syntagma station). Extensions are under construction to the west southwest suburbs of Athens, as far as the port of Piraeus. The new stations will be Agia Barvara, Koridallos, Nikaia, Maniatika, Piraeus and Dimotiko Theatro. The completed extension will be ready in 2022 (the first three stations will open in summer 2020), connecting the biggest port of Greece, Piraeus Port, with the biggest airport of Greece the Athens International Airport.
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The Athens commuter rail service, referred to as the "Proastiakós", connects Eleftherios Venizelos International Airport to the city of Kiato, 106 km (66 mi)[125] west of Athens, via Larissa station, the city's central rail station and the port of Piraeus. The length of Athens's commuter rail network extends to 120 km (75 mi),[125] and is expected to stretch to 281 km (175 mi) by 2010.[125]
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Athens tram is operated by STASY S.A (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation). It has a fleet of 35 Sirio type vehicles[126] which serve 48 stations,[126] employ 345 people with an average daily occupancy of 65,000 passengers.[126] The tram network spans a total length of 27 km (17 mi) and covers ten Athenian suburbs.[126] The network runs from Syntagma Square to the southwestern suburb of Palaio Faliro, where the line splits in two branches; the first runs along the Athens coastline toward the southern suburb of Voula, while the other heads toward Neo Faliro. The network covers the majority of the Athens coastline.[127] Further extensions are planned towards the major commercial port of Piraeus.[126] The expansion to Piraeus will include 12 new stations, increase the overall length of tram route by 5.4 km (3 mi), and increase the overall transportation network.[128]
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Athens is served by the Athens International Airport (ATH), located near the town of Spata, in the eastern Messoghia plain, some 35 km (22 mi) east of center of Athens.[129] The airport, awarded the "European Airport of the Year 2004" Award,[130] is intended as an expandable hub for air travel in southeastern Europe and was constructed in 51 months, costing 2.2 billion euros. It employs a staff of 14,000.[130]
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The airport is served by the Metro, the suburban rail, buses to Piraeus port, Athens' City Centre, Liosion and Kifisos Intercity bus stations and Elliniko metro's line 2 southern terminal, and also taxis. The airport accommodates 65 landings and take-offs per hour,[129] with its 24-passenger boarding bridges,[129] 144 check-in counters and broader 150,000 m2 (1,614,587 sq ft) main terminal;[129] and a commercial area of 7,000 m2 (75,347 sq ft) which includes cafés, duty-free shops,[130] and a small museum.
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In 2018, the airport handled 24,135,736 a huge increase over the last 4 years. In 2014, the airport handled 15,196,369 passengers, an increase of 21.2% over the previous year of 2013.[131] Of those 15,196,369 passengers, 5,267,593 passed through the airport for domestic flights,[132] and 9,970,006 passengers travelled through for international flights.[132] Beyond the dimensions of its passenger capacity, ATH handled 205,294 total flights in 2007, or approximately 562 flights per day.[133]
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Athens is the hub of the country's national railway system (OSE), connecting the capital with major cities across Greece and abroad (Istanbul, Sofia, Belgrade and Bucharest). The Port of Piraeus is the largest port in Greece and one of the largest in Europe. It is the largest container port in East Mediterranean Sea Basin. It is also the busiest passenger port in Europe and one of the largest passenger ports in the world. Connects Athens to the numerous Greek islands of the Aegean Sea, with ferries departing, while also serving the cruise ships that arrive.[134][135][136] Rafina and Lavrio act as alternative ports of Athens, connects the city with numerous Greek islands of the Aegean Sea, Evia and Cesme in Turkey,[137][138] while also serving the cruise ships that arrive.
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Two main motorways of Greece begin in Athens, namely the A1/E75, heading north towards Greece's second largest city, Thessaloniki; and the border crossing of Evzones and the A8/E94 heading west, towards Greece's third largest city, Patras, which incorporated the GR-8A. Before their completion much of the road traffic used the GR-1 and the GR-8.
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Athens' Metropolitan Area is served by the motorway network of the Attiki Odos toll-motorway (code: A6). Its main section extends from the western industrial suburb of Elefsina to Athens International Airport; while two beltways, namely the Aigaleo Beltway (A65) and the Hymettus Beltway (A64) serve parts of western and eastern Athens respectively. The span of the Attiki Odos in all its length is 65 km (40 mi),[139] making it the largest metropolitan motorway network in all of Greece.
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Located on Panepistimiou Street, the old campus of the University of Athens, the National Library, and the Athens Academy form the "Athens Trilogy" built in the mid-19th century. Most of the university's workings have been moved to a much larger, modern campus located in the eastern suburb of Zografou. The second higher education institution in the city is the Athens Polytechnic School, found in Patission Street. This was the location where on 17 November 1973, more than 13 students were killed and hundreds injured inside the university during the Athens Polytechnic uprising,[140] against the military junta that ruled the nation from 21 April 1967 until 23 July 1974.
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Other universities that lie within Athens are the Athens University of Economics and Business, the Panteion University, the Agricultural University of Athens and the University of Piraeus. There are overall ten state-supported Institutions of Higher (or Tertiary) education located in the Athens Urban Area, these are by chronological order: Athens School of Fine Arts (1837), National Technical University of Athens (1837), National and Kapodistrian University of Athens (1837), Agricultural University of Athens (1920), Athens University of Economics and Business (1920), Panteion University of Social and Political Sciences (1927), University of Piraeus (1938), Harokopio University of Athens (1990), School of Pedagogical and Technological Education (2002), University of West Attica (2018). There are also several other private colleges, as they called formally in Greece, as the establishment of private universities is prohibited by the constitution. Many of them are accredited by a foreign state or university such as the American College of Greece and the Athens Campus of the University of Indianapolis.[141]
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The city is a world centre of archaeological research. Along with national institutions, such as the Athens University and the Archaeological Society, there are multiple archaeological Museums including the National Archaeological Museum, the Cycladic Museum, the Epigraphic Museum, the Byzantine & Christian Museum, as well as museums at the ancient Agora, Acropolis, Kerameikos, and the Kerameikos Archaeological Museum. The city is also home to the Demokritos laboratory for Archaeometry, alongside regional and national archaeological authorities that form part of the Greek Department of Culture.
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Athens hosts 17 Foreign Archaeological Institutes which promote and facilitate research by scholars from their home countries. As a result, Athens has more than a dozen archaeological libraries and three specialized archaeological laboratories, and is the venue of several hundred specialized lectures, conferences and seminars, as well as dozens of archaeological exhibitions, each year. At any given time, hundreds of international scholars and researchers in all disciplines of archaeology are to be found in the city.
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Athens incorporates architectural styles ranging from Greco-Roman and Neoclassical to modern times. They are often to be found in the same areas, as Athens is not marked by a uniformity of architectural style. A visitor will quickly notice the absence of tall buildings: Athens has very strict height restriction laws in order to ensure the Acropolis hill is visible throughout the city. There is evidence of continuity in elements of the architectural environment through the city's history, although this is still a contested subject [142].
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For the greatest part of the 19th century Neoclassicism dominated Athens, as well as some deviations from it such as Eclecticism, especially in the early 20th century. Thus, the Old Royal Palace was the first important public building to be built, between 1836 and 1843. Later in the mid and late 19th century, Theophil Freiherr von Hansen and Ernst Ziller took part in the construction of many neoclassical buildings such as the Athens Academy and the Zappeion Hall. Ziller also designed many private mansions in the centre of Athens which gradually became public, usually through donations, such as Schliemann's Iliou Melathron.
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Beginning in the 1920s, modern architecture including Bauhaus and Art Deco began to exert an influence on almost all Greek architects, and buildings both public and private were constructed in accordance with these styles. Localities with a great number of such buildings include Kolonaki, and some areas of the centre of the city; neighbourhoods developed in this period include Kypseli.[143]
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In the 1950s and 1960s during the extension and development of Athens, other modern movements such as the International style played an important role. The centre of Athens was largely rebuilt, leading to the demolition of a number of neoclassical buildings. The architects of this era employed materials such as glass, marble and aluminium, and many blended modern and classical elements.[144] After World War II, internationally known architects to have designed and built in the city included Walter Gropius, with his design for the US Embassy, and, among others, Eero Saarinen, in his postwar design for the east terminal of the Ellinikon Airport.
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All over the city can be found several statues or busts. Apart from the neoclassicals by Leonidas Drosis at the Academy of Athens (Plato, Socrates, Apollo, Athena), other notable include the statue of Theseus by Georgios Fytalis at Thiseion, of philhellenes like Lord Byron, George Canning and William Gladstone, the equestrian statue of Theodoros Kolokotronis by Lazaros Sochos in front of the Old Parliament, statues of Ioannis Kapodistrias, Rigas Feraios and Adamantios Korais at the University, of Evangelos Zappas and Konstantinos Zappas at Zappeion, of Ioannis Varvakis at the National Garden, the "woodbreaker" by Dimitrios Filippotis, the equestrian statue of Alexandros Papagos at Papagou district and various busts of fighters of Greek independence at the Pedion tou Areos. A significant landmark is also the Tomb of the Unknown Soldier in Syntagma.
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Athens' most important museums include:
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Athens has been a destination for travellers since antiquity. Over the past decade, the city's infrastructure and social amenities have improved, in part because of its successful bid to stage the 2004 Olympic Games. The Greek Government, aided by the EU, has funded major infrastructure projects such as the state-of-the-art Eleftherios Venizelos International Airport,[145] the expansion of the Athens Metro system,[68] and the new Attiki Odos Motorway.[68]
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Athens was voted as the third best European city to visit in 2015 by European Best Destination. More than 240,000 people voted.
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Athens is home to 148 theatrical stages, more than any other city in the world, including the ancient Odeon of Herodes Atticus, home to the Athens Festival, which runs from May to October each year.[146][147] In addition to a large number of multiplexes, Athens plays host to open air garden cinemas. The city also supports music venues, including the Athens Concert Hall (Megaro Moussikis), which attracts world class artists.[148] The Athens Planetarium,[149] located in Andrea Syngrou Avenue, in Palaio Faliro[150] is one of the largest and best equipped digital planetaria in the world.[151] The Stavros Niarchos Foundation Cultural Center, inaugurated in 2016, will house the National Library of Greece and the Greek National Opera.[152]
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The most successful songs during the period 1870–1930 were the so-called Athenian serenades (Αθηναϊκές καντάδες), based on the Heptanesean kantádhes (καντάδες 'serenades'; sing.: καντάδα) and the songs performed on stage (επιθεωρησιακά τραγούδια 'theatrical revue songs') in revues, musical comedies, operettas and nocturnes that were dominating Athens' theatre scene.
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Notable composers of operettas or nocturnes were Kostas Giannidis, Dionysios Lavrangas, Nikos Hatziapostolou, while Theophrastos Sakellaridis' The Godson remains probably the most popular operetta. Despite the fact that the Athenian songs were not autonomous artistic creations (in contrast with the serenades) and despite their original connection with mainly dramatic forms of Art, they eventually became hits as independent songs. Notable actors of Greek operettas, who made also a series of melodies and songs popular at that time, include Orestis Makris, Kalouta sisters, Vasilis Avlonitis, Afroditi Laoutari, Eleni Papadaki, Marika Nezer, Marika Krevata and others. After 1930, wavering among American and European musical influences as well as the Greek musical tradition. Greek composers begin to write music using the tunes of the tango, waltz, swing, foxtrot, some times combined with melodies in the style of Athenian serenades' repertory. Nikos Gounaris was probably the most renowned composer and singer of the time.
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In 1923, after the population exchange between Greece and Turkey, many ethnic Greeks from Asia Minor fled to Athens as a result of the Greco-Turkish War. They settled in poor neighborhoods and brought with them Rebetiko music, making it popular also in Greece, which became later the base for the Laïko music. Other forms of song popular today in Greece are elafrolaika, entechno, dimotika, and skyladika.[153] Greece's most notable, and internationally famous, composers of Greek song, mainly of the entechno form, are Manos Hadjidakis and Mikis Theodorakis. Both composers have achieved fame abroad for their composition of film scores.[153]
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Athens has a long tradition in sports and sporting events, serving as home to the most important clubs in Greek sport and housing a large number of sports facilities. The city has also been host to sports events of international importance.
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Athens has hosted the Summer Olympic Games twice, in 1896 and 2004. The 2004 Summer Olympics required the development of the Athens Olympic Stadium, which has since gained a reputation as one of the most beautiful stadiums in the world, and one of its most interesting modern monuments.[154] The biggest stadium in the country, it hosted two finals of the UEFA Champions League, in 1994 and 2007. Athens' other major stadium, located in the Piraeus area, is the Karaiskakis Stadium, a sports and entertainment complex, host of the 1971 UEFA Cup Winners' Cup Final.
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Athens has hosted the EuroLeague final three times, the first in 1985 and second in 1993, both at the Peace and Friendship Stadium, most known as SEF, a large indoor arena,[155] and the third time in 2007 at the Olympic Indoor Hall. Events in other sports such as athletics, volleyball, water polo etc., have been hosted in the capital's venues.
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Athens is home to three European multi-sport clubs: Olympiacos, Panathinaikos, AEK Athens. In football, Olympiacos have dominated the domestic competitions, Panathinaikos made it to the 1971 European Cup Final, while AEK Athens is the other member of the big three. These clubs also have basketball teams; Panathinaikos and Olympiacos are among the top powers in European basketball, having won the Euroleague six times and three respectively, whilst AEK Athens was the first Greek team to win a European trophy in any team sport.
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Other notable clubs within Athens are Athinaikos, Panionios, Atromitos, Apollon, Panellinios, Ethnikos Piraeus, Maroussi BCE and Peristeri B.C.. Athenian clubs have also had domestic and international success in other sports.
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The Athens area encompasses a variety of terrain, notably hills and mountains rising around the city, and the capital is the only major city in Europe to be bisected by a mountain range. Four mountain ranges extend into city boundaries and thousands of miles of trails criss-cross the city and neighbouring areas, providing exercise and wilderness access on foot and bike.
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Beyond Athens and across the prefecture of Attica, outdoor activities include skiing, rock climbing, hang gliding and windsurfing. Numerous outdoor clubs serve these sports, including the Athens Chapter of the Sierra Club, which leads over 4,000 outings annually in the area.
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Beside the above clubs, inside the boundaries of Athens Municipality there are some more clubs with presence in national divisions or notable action for short periods. Some of them are PAO Rouf (Rouf) with earlier presence in Gamma Ethniki, Petralona F.C.(el) (Petralona), football club founded in 1963, with earlier presence in Beta Ethniki, Attikos F.C.(el) (Kolonos), football club founded in 1919 with short presence in Gamma Ethniki, Athinais Kypselis [es] (Kypseli), football club founded in 1938 with short presence in Gamma Ethniki, Gyziakos (Gyzi), basketball club founded in 1937 with short presence in Beta Ethniki basketball and Aetos B.C.(el) (Agios Panteleimonas), basketball club founded in 1992 with earlier presence in A2 Ethniki Basketball. Another important Athenian sport club is the Athens Tennis Club founded in 1895 with important offer for the Greek tennis.[156]
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The revival of the modern Olympic Games was brought forth in 1896, by Frenchman Pierre de Coubertin. Thanks to his efforts, Athens was awarded the first modern Olympic Games. In 1896, the city had a population of 123,000[93] and the event helped boost the city's international profile. Of the venues used for these Olympics, the Kallimarmaro Stadium, and Zappeion were most crucial. The Kallimarmaro is a replica of the ancient Athenian stadiums, and the only major stadium (in its capacity of 60,000) to be made entirely of white marble from Mount Penteli, the same material used for construction of the Parthenon.
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Fencing before the king of Greece at the 1896 Summer Olympics.
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The Panathenaic Stadium of Athens (Kallimarmaron) dates back to the 4th century BC and has hosted the first modern Olympic Games in 1896.
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The 1906 Summer Olympics, or the 1906 Intercalated games, were held in Athens. The intercalated competitions were intermediate games to the internationally organized Olympics, and were meant to be organized in Greece every four years, between the main Olympics. This idea later lost support from the IOC and these games were discontinued.
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Athens was awarded the 2004 Summer Olympics on 5 September 1997 in Lausanne, Switzerland, after having lost a previous bid to host the 1996 Summer Olympics, to Atlanta, United States.[17] It was to be the second time Athens would host the games, following the inaugural event of 1896. After an unsuccessful bid in 1990, the 1997 bid was radically improved, including an appeal to Greece's Olympic history. In the last round of voting, Athens defeated Rome with 66 votes to 41.[17] Prior to this round, the cities of Buenos Aires, Stockholm and Cape Town had been eliminated from competition, having received fewer votes.[17]
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During the first three years of preparations, the International Olympic Committee had expressed concern over the speed of construction progress for some of the new Olympic venues. In 2000 the Organising Committee's president was replaced by Gianna Angelopoulos-Daskalaki, who was the president of the original Bidding Committee in 1997. From that point forward, preparations continued at a highly accelerated, almost frenzied pace.
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Although the heavy cost was criticized, estimated at $1.5 billion, Athens was transformed into a more functional city that enjoys modern technology both in transportation and in modern urban development.[157] Some of the finest sporting venues in the world were created in the city, all of which were fully ready for the games. The games welcomed over 10,000 athletes from all 202 countries.[157]
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The 2004 Games were judged a success, as both security and organization worked well, and only a few visitors reported minor problems mainly concerning accommodation issues. The 2004 Olympic Games were described as Unforgettable, dream Games, by IOC President Jacques Rogge for their return to the birthplace of the Olympics, and for meeting the challenges of holding the Olympic Games.[157] The only observable problem was a somewhat sparse attendance of some early events. Eventually, however, a total of more than 3.5 million tickets were sold, which was higher than any other Olympics with the exception of Sydney (more than 5 million tickets were sold there in 2000).[158]
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In 2008 it was reported that most of the Olympic venues had fallen into disrepair: according to those reports, 21 of the 22 facilities built for the games had either been left abandoned or are in a state of dereliction, with several squatter camps having sprung up around certain facilities, and a number of venues afflicted by vandalism, graffiti or strewn with rubbish.[159][160] These claims, however, are disputed and likely to be inaccurate, as most of the facilities used for the Athens Olympics are either in use or in the process of being converted for post-Olympics use. The Greek Government has created a corporation, Olympic Properties SA, which is overseeing the post-Olympics management, development and conversion of these facilities, some of which will be sold off (or have already been sold off) to the private sector,[161] while other facilities are still in use just as during the Olympics, or have been converted for commercial use or modified for other sports.[162] Concerts and theatrical shows, such as those by the troupe Cirque du Soleil, have recently been held in the complex.[153]
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1 |
+
|
2 |
+
|
3 |
+
In physics, an orbit is the gravitationally curved trajectory of an object,[1] such as the trajectory of a planet around a star or a natural satellite around a planet. Normally, orbit refers to a regularly repeating trajectory, although it may also refer to a non-repeating trajectory. To a close approximation, planets and satellites follow elliptic orbits, with the center of mass being orbited at a focal point of the ellipse,[2] as described by Kepler's laws of planetary motion.
|
4 |
+
|
5 |
+
For most situations, orbital motion is adequately approximated by Newtonian mechanics, which explains gravity as a force obeying an inverse-square law.[3] However, Albert Einstein's general theory of relativity, which accounts for gravity as due to curvature of spacetime, with orbits following geodesics, provides a more accurate calculation and understanding of the exact mechanics of orbital motion.
|
6 |
+
|
7 |
+
Historically, the apparent motions of the planets were described by European and Arabic philosophers using the idea of celestial spheres. This model posited the existence of perfect moving spheres or rings to which the stars and planets were attached. It assumed the heavens were fixed apart from the motion of the spheres, and was developed without any understanding of gravity. After the planets' motions were more accurately measured, theoretical mechanisms such as deferent and epicycles were added. Although the model was capable of reasonably accurately predicting the planets' positions in the sky, more and more epicycles were required as the measurements became more accurate, hence the model became increasingly unwieldy. Originally geocentric, it was modified by Copernicus to place the Sun at the centre to help simplify the model. The model was further challenged during the 16th century, as comets were observed traversing the spheres.[4][5]
|
8 |
+
|
9 |
+
The basis for the modern understanding of orbits was first formulated by Johannes Kepler whose results are summarised in his three laws of planetary motion. First, he found that the orbits of the planets in our Solar System are elliptical, not circular (or epicyclic), as had previously been believed, and that the Sun is not located at the center of the orbits, but rather at one focus.[6] Second, he found that the orbital speed of each planet is not constant, as had previously been thought, but rather that the speed depends on the planet's distance from the Sun. Third, Kepler found a universal relationship between the orbital properties of all the planets orbiting the Sun. For the planets, the cubes of their distances from the Sun are proportional to the squares of their orbital periods. Jupiter and Venus, for example, are respectively about 5.2 and 0.723 AU distant from the Sun, their orbital periods respectively about 11.86 and 0.615 years. The proportionality is seen by the fact that the ratio for Jupiter, 5.23/11.862, is practically equal to that for Venus, 0.7233/0.6152, in accord with the relationship. Idealised orbits meeting these rules are known as Kepler orbits.
|
10 |
+
|
11 |
+
Isaac Newton demonstrated that Kepler's laws were derivable from his theory of gravitation and that, in general, the orbits of bodies subject to gravity were conic sections (this assumes that the force of gravity propagates instantaneously). Newton showed that, for a pair of bodies, the orbits' sizes are in inverse proportion to their masses, and that those bodies orbit their common center of mass. Where one body is much more massive than the other (as is the case of an artificial satellite orbiting a planet), it is a convenient approximation to take the center of mass as coinciding with the center of the more massive body.
|
12 |
+
|
13 |
+
Advances in Newtonian mechanics were then used to explore variations from the simple assumptions behind Kepler orbits, such as the perturbations due to other bodies, or the impact of spheroidal rather than spherical bodies. Lagrange (1736–1813) developed a new approach to Newtonian mechanics emphasizing energy more than force, and made progress on the three body problem, discovering the Lagrangian points. In a dramatic vindication of classical mechanics, in 1846 Urbain Le Verrier was able to predict the position of Neptune based on unexplained perturbations in the orbit of Uranus.
|
14 |
+
|
15 |
+
Albert Einstein (1879-1955) in his 1916 paper The Foundation of the General Theory of Relativity explained that gravity was due to curvature of space-time and removed Newton's assumption that changes propagate instantaneously. This led astronomers to recognize that Newtonian mechanics did not provide the highest accuracy in understanding orbits. In relativity theory, orbits follow geodesic trajectories which are usually approximated very well by the Newtonian predictions (except where there are very strong gravity fields and very high speeds) but the differences are measurable. Essentially all the experimental evidence that can distinguish between the theories agrees with relativity theory to within experimental measurement accuracy. The original vindication of general relativity is that it was able to account for the remaining unexplained amount in precession of Mercury's perihelion first noted by Le Verrier. However, Newton's solution is still used for most short term purposes since it is significantly easier to use and sufficiently accurate.
|
16 |
+
|
17 |
+
Within a planetary system, planets, dwarf planets, asteroids and other minor planets, comets, and space debris orbit the system's barycenter in elliptical orbits. A comet in a parabolic or hyperbolic orbit about a barycenter is not gravitationally bound to the star and therefore is not considered part of the star's planetary system. Bodies which are gravitationally bound to one of the planets in a planetary system, either natural or artificial satellites, follow orbits about a barycenter near or within that planet.
|
18 |
+
|
19 |
+
Owing to mutual gravitational perturbations, the eccentricities of the planetary orbits vary over time. Mercury, the smallest planet in the Solar System, has the most eccentric orbit. At the present epoch, Mars has the next largest eccentricity while the smallest orbital eccentricities are seen with Venus and Neptune.
|
20 |
+
|
21 |
+
As two objects orbit each other, the periapsis is that point at which the two objects are closest to each other and the apoapsis is that point at which they are the farthest. (More specific terms are used for specific bodies. For example, perigee and apogee are the lowest and highest parts of an orbit around Earth, while perihelion and aphelion are the closest and farthest points of an orbit around the Sun.)
|
22 |
+
|
23 |
+
In the case of planets orbiting a star, the mass of the star and all its satellites are calculated to be at a single point called the barycenter. The paths of all the star's satellites are elliptical orbits about that barycenter. Each satellite in that system will have its own elliptical orbit with the barycenter at one focal point of that ellipse. At any point along its orbit, any satellite will have a certain value of kinetic and potential energy with respect to the barycenter, and that energy is a constant value at every point along its orbit. As a result, as a planet approaches periapsis, the planet will increase in speed as its potential energy decreases; as a planet approaches apoapsis, its velocity will decrease as its potential energy increases.
|
24 |
+
|
25 |
+
There are a few common ways of understanding orbits:
|
26 |
+
|
27 |
+
As an illustration of an orbit around a planet, the Newton's cannonball model may prove useful (see image below). This is a 'thought experiment', in which a cannon on top of a tall mountain is able to fire a cannonball horizontally at any chosen muzzle speed. The effects of air friction on the cannonball are ignored (or perhaps the mountain is high enough that the cannon is above the Earth's atmosphere, which is the same thing).[7]
|
28 |
+
|
29 |
+
If the cannon fires its ball with a low initial speed, the trajectory of the ball curves downward and hits the ground (A). As the firing speed is increased, the cannonball hits the ground farther (B) away from the cannon, because while the ball is still falling towards the ground, the ground is increasingly curving away from it (see first point, above). All these motions are actually "orbits" in a technical sense – they are describing a portion of an elliptical path around the center of gravity – but the orbits are interrupted by striking the Earth.
|
30 |
+
|
31 |
+
If the cannonball is fired with sufficient speed, the ground curves away from the ball at least as much as the ball falls – so the ball never strikes the ground. It is now in what could be called a non-interrupted, or circumnavigating, orbit. For any specific combination of height above the center of gravity and mass of the planet, there is one specific firing speed (unaffected by the mass of the ball, which is assumed to be very small relative to the Earth's mass) that produces a circular orbit, as shown in (C).
|
32 |
+
|
33 |
+
As the firing speed is increased beyond this, non-interrupted elliptic orbits are produced; one is shown in (D). If the initial firing is above the surface of the Earth as shown, there will also be non-interrupted elliptical orbits at slower firing speed; these will come closest to the Earth at the point half an orbit beyond, and directly opposite the firing point, below the circular orbit.
|
34 |
+
|
35 |
+
At a specific horizontal firing speed called escape velocity, dependent on the mass of the planet, an open orbit (E) is achieved that has a parabolic path. At even greater speeds the object will follow a range of hyperbolic trajectories. In a practical sense, both of these trajectory types mean the object is "breaking free" of the planet's gravity, and "going off into space" never to return.
|
36 |
+
|
37 |
+
The velocity relationship of two moving objects with mass can thus be considered in four practical classes, with subtypes:
|
38 |
+
|
39 |
+
It is worth noting that orbital rockets are launched vertically at first to lift the rocket above the atmosphere (which causes frictional drag), and then slowly pitch over and finish firing the rocket engine parallel to the atmosphere to achieve orbit speed.
|
40 |
+
|
41 |
+
Once in orbit, their speed keeps them in orbit above the atmosphere. If e.g., an elliptical orbit dips into dense air, the object will lose speed and re-enter (i.e. fall). Occasionally a space craft will intentionally intercept the atmosphere, in an act commonly referred to as an aerobraking maneuver.
|
42 |
+
|
43 |
+
In most situations relativistic effects can be neglected, and Newton's laws give a sufficiently accurate description of motion. The acceleration of a body is equal to the sum of the forces acting on it, divided by its mass, and the gravitational force acting on a body is proportional to the product of the masses of the two attracting bodies and decreases inversely with the square of the distance between them. To this Newtonian approximation, for a system of two-point masses or spherical bodies, only influenced by their mutual gravitation (called a two-body problem), their trajectories can be exactly calculated. If the heavier body is much more massive than the smaller, as in the case of a satellite or small moon orbiting a planet or for the Earth orbiting the Sun, it is accurate enough and convenient to describe the motion in terms of a coordinate system that is centered on the heavier body, and we say that the lighter body is in orbit around the heavier. For the case where the masses of two bodies are comparable, an exact Newtonian solution is still sufficient and can be had by placing the coordinate system at the center of mass of the system.
|
44 |
+
|
45 |
+
Energy is associated with gravitational fields. A stationary body far from another can do external work if it is pulled towards it, and therefore has gravitational potential energy. Since work is required to separate two bodies against the pull of gravity, their gravitational potential energy increases as they are separated, and decreases as they approach one another. For point masses the gravitational energy decreases to zero as they approach zero separation. It is convenient and conventional to assign the potential energy as having zero value when they are an infinite distance apart, and hence it has a negative value (since it decreases from zero) for smaller finite distances.
|
46 |
+
|
47 |
+
When only two gravitational bodies interact, their orbits follow a conic section. The orbit can be open (implying the object never returns) or closed (returning). Which it is depends on the total energy (kinetic + potential energy) of the system. In the case of an open orbit, the speed at any position of the orbit is at least the escape velocity for that position, in the case of a closed orbit, the speed is always less than the escape velocity. Since the kinetic energy is never negative, if the common convention is adopted of taking the potential energy as zero at infinite separation, the bound orbits will have negative total energy, the parabolic trajectories zero total energy, and hyperbolic orbits positive total energy.
|
48 |
+
|
49 |
+
An open orbit will have a parabolic shape if it has velocity of exactly the escape velocity at that point in its trajectory, and it will have the shape of a hyperbola when its velocity is greater than the escape velocity. When bodies with escape velocity or greater approach each other, they will briefly curve around each other at the time of their closest approach, and then separate, forever.
|
50 |
+
|
51 |
+
All closed orbits have the shape of an ellipse. A circular orbit is a special case, wherein the foci of the ellipse coincide. The point where the orbiting body is closest to Earth is called the perigee, and is called the periapsis (less properly, "perifocus" or "pericentron") when the orbit is about a body other than Earth. The point where the satellite is farthest from Earth is called the apogee, apoapsis, or sometimes apifocus or apocentron. A line drawn from periapsis to apoapsis is the line-of-apsides. This is the major axis of the ellipse, the line through its longest part.
|
52 |
+
|
53 |
+
Bodies following closed orbits repeat their paths with a certain time called the period. This motion is described by the empirical laws of Kepler, which can be mathematically derived from Newton's laws. These can be
|
54 |
+
formulated as follows:
|
55 |
+
|
56 |
+
Note that while bound orbits of a point mass or a spherical body with a Newtonian gravitational field are closed ellipses, which repeat the same path exactly and indefinitely, any non-spherical or non-Newtonian effects (such as caused by the slight oblateness of the Earth, or by relativistic effects, thereby changing the gravitational field's behavior with distance) will cause the orbit's shape to depart from the closed ellipses characteristic of Newtonian two-body motion. The two-body solutions were published by Newton in Principia in 1687. In 1912, Karl Fritiof Sundman developed a converging infinite series that solves the three-body problem; however, it converges too slowly to be of much use. Except for special cases like the Lagrangian points, no method is known to solve the equations of motion for a system with four or more bodies.
|
57 |
+
|
58 |
+
Rather than an exact closed form solution, orbits with many bodies can be approximated with arbitrarily high accuracy. These approximations take two forms:
|
59 |
+
|
60 |
+
Differential simulations with large numbers of objects perform the calculations in a hierarchical pairwise fashion between centers of mass. Using this scheme, galaxies, star clusters and other large assemblages of objects have been simulated.[citation needed]
|
61 |
+
|
62 |
+
The Earth follows an ellipse round the sun.
|
63 |
+
But unlike the ellipse followed by a pendulum or an object attached to a spring, the sun is at a focal point of the ellipse and not at its centre.
|
64 |
+
|
65 |
+
The following derivation applies to such an elliptical orbit.
|
66 |
+
We start only with the Newtonian law of gravitation stating that the gravitational acceleration towards the central body is related to the inverse of the square of the distance between them, namely
|
67 |
+
|
68 |
+
where F2 is the force acting on the mass m2 caused by the gravitational attraction mass m1 has for m2, G is the universal gravitational constant, and r is the distance between the two masses centers.
|
69 |
+
|
70 |
+
From Newton's Second Law, the summation of the forces acting on m2 related to that bodies acceleration:
|
71 |
+
|
72 |
+
where A2 is the acceleration of m2 caused by the force of gravitational attraction F2 of m1 acting on m2.
|
73 |
+
|
74 |
+
Combining Eq 1 and 2:
|
75 |
+
|
76 |
+
Solving for the acceleration, A2:
|
77 |
+
|
78 |
+
where
|
79 |
+
|
80 |
+
|
81 |
+
|
82 |
+
μ
|
83 |
+
|
84 |
+
|
85 |
+
|
86 |
+
{\displaystyle \mu \,}
|
87 |
+
|
88 |
+
is the standard gravitational parameter, in this case
|
89 |
+
|
90 |
+
|
91 |
+
|
92 |
+
G
|
93 |
+
|
94 |
+
m
|
95 |
+
|
96 |
+
1
|
97 |
+
|
98 |
+
|
99 |
+
|
100 |
+
|
101 |
+
{\displaystyle Gm_{1}}
|
102 |
+
|
103 |
+
. It is understood that the system being described is m2, hence the subscripts can be dropped.
|
104 |
+
|
105 |
+
We assume that the central body is massive enough that it can be considered to be stationary and we ignore the more subtle effects of general relativity.
|
106 |
+
|
107 |
+
When a pendulum or an object attached to a spring swings in an ellipse,
|
108 |
+
the inward acceleration/force is proportional to the distance
|
109 |
+
|
110 |
+
|
111 |
+
|
112 |
+
A
|
113 |
+
=
|
114 |
+
F
|
115 |
+
|
116 |
+
/
|
117 |
+
|
118 |
+
m
|
119 |
+
=
|
120 |
+
−
|
121 |
+
k
|
122 |
+
r
|
123 |
+
.
|
124 |
+
|
125 |
+
|
126 |
+
{\displaystyle A=F/m=-kr.}
|
127 |
+
|
128 |
+
|
129 |
+
Due to the way vectors add, the component of the force in the
|
130 |
+
|
131 |
+
|
132 |
+
|
133 |
+
|
134 |
+
|
135 |
+
|
136 |
+
|
137 |
+
x
|
138 |
+
|
139 |
+
^
|
140 |
+
|
141 |
+
|
142 |
+
|
143 |
+
|
144 |
+
|
145 |
+
{\displaystyle {\hat {\mathbf {x} }}}
|
146 |
+
|
147 |
+
or in the
|
148 |
+
|
149 |
+
|
150 |
+
|
151 |
+
|
152 |
+
|
153 |
+
|
154 |
+
|
155 |
+
y
|
156 |
+
|
157 |
+
^
|
158 |
+
|
159 |
+
|
160 |
+
|
161 |
+
|
162 |
+
|
163 |
+
{\displaystyle {\hat {\mathbf {y} }}}
|
164 |
+
|
165 |
+
directions are also proportionate to the respective
|
166 |
+
components of the distances,
|
167 |
+
|
168 |
+
|
169 |
+
|
170 |
+
|
171 |
+
r
|
172 |
+
|
173 |
+
x
|
174 |
+
|
175 |
+
″
|
176 |
+
|
177 |
+
=
|
178 |
+
|
179 |
+
A
|
180 |
+
|
181 |
+
x
|
182 |
+
|
183 |
+
|
184 |
+
=
|
185 |
+
−
|
186 |
+
k
|
187 |
+
|
188 |
+
r
|
189 |
+
|
190 |
+
x
|
191 |
+
|
192 |
+
|
193 |
+
|
194 |
+
|
195 |
+
{\displaystyle r''_{x}=A_{x}=-kr_{x}}
|
196 |
+
|
197 |
+
. Hence, the entire analysis can be done separately in these dimensions. This results in the harmonic parabolic equations
|
198 |
+
|
199 |
+
|
200 |
+
|
201 |
+
x
|
202 |
+
=
|
203 |
+
A
|
204 |
+
cos
|
205 |
+
|
206 |
+
(
|
207 |
+
t
|
208 |
+
)
|
209 |
+
|
210 |
+
|
211 |
+
{\displaystyle x=A\cos(t)}
|
212 |
+
|
213 |
+
and
|
214 |
+
|
215 |
+
|
216 |
+
|
217 |
+
y
|
218 |
+
=
|
219 |
+
B
|
220 |
+
sin
|
221 |
+
|
222 |
+
(
|
223 |
+
t
|
224 |
+
)
|
225 |
+
|
226 |
+
|
227 |
+
{\displaystyle y=B\sin(t)}
|
228 |
+
|
229 |
+
of the ellipse. In contrast, with the decreasing relationship
|
230 |
+
|
231 |
+
|
232 |
+
|
233 |
+
A
|
234 |
+
=
|
235 |
+
μ
|
236 |
+
|
237 |
+
/
|
238 |
+
|
239 |
+
|
240 |
+
r
|
241 |
+
|
242 |
+
2
|
243 |
+
|
244 |
+
|
245 |
+
|
246 |
+
|
247 |
+
{\displaystyle A=\mu /r^{2}}
|
248 |
+
|
249 |
+
, the dimensions cannot be separated.[citation needed]
|
250 |
+
|
251 |
+
The location of the orbiting object at the current time
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
t
|
256 |
+
|
257 |
+
|
258 |
+
{\displaystyle t}
|
259 |
+
|
260 |
+
is located in the plane using
|
261 |
+
Vector calculus in polar coordinates both with the standard Euclidean basis and with the polar basis
|
262 |
+
with the origin coinciding with the center of force.
|
263 |
+
Let
|
264 |
+
|
265 |
+
|
266 |
+
|
267 |
+
r
|
268 |
+
|
269 |
+
|
270 |
+
{\displaystyle r}
|
271 |
+
|
272 |
+
be the distance between the object and the center and
|
273 |
+
|
274 |
+
|
275 |
+
|
276 |
+
|
277 |
+
θ
|
278 |
+
|
279 |
+
|
280 |
+
{\displaystyle \theta }
|
281 |
+
|
282 |
+
be the angle it has rotated.
|
283 |
+
Let
|
284 |
+
|
285 |
+
|
286 |
+
|
287 |
+
|
288 |
+
|
289 |
+
|
290 |
+
|
291 |
+
x
|
292 |
+
|
293 |
+
^
|
294 |
+
|
295 |
+
|
296 |
+
|
297 |
+
|
298 |
+
|
299 |
+
{\displaystyle {\hat {\mathbf {x} }}}
|
300 |
+
|
301 |
+
and
|
302 |
+
|
303 |
+
|
304 |
+
|
305 |
+
|
306 |
+
|
307 |
+
|
308 |
+
|
309 |
+
y
|
310 |
+
|
311 |
+
^
|
312 |
+
|
313 |
+
|
314 |
+
|
315 |
+
|
316 |
+
|
317 |
+
{\displaystyle {\hat {\mathbf {y} }}}
|
318 |
+
|
319 |
+
be the standard Euclidean bases and let
|
320 |
+
|
321 |
+
|
322 |
+
|
323 |
+
|
324 |
+
|
325 |
+
|
326 |
+
|
327 |
+
r
|
328 |
+
|
329 |
+
^
|
330 |
+
|
331 |
+
|
332 |
+
|
333 |
+
=
|
334 |
+
cos
|
335 |
+
|
336 |
+
(
|
337 |
+
θ
|
338 |
+
)
|
339 |
+
|
340 |
+
|
341 |
+
|
342 |
+
|
343 |
+
x
|
344 |
+
|
345 |
+
^
|
346 |
+
|
347 |
+
|
348 |
+
|
349 |
+
+
|
350 |
+
sin
|
351 |
+
|
352 |
+
(
|
353 |
+
θ
|
354 |
+
)
|
355 |
+
|
356 |
+
|
357 |
+
|
358 |
+
|
359 |
+
y
|
360 |
+
|
361 |
+
^
|
362 |
+
|
363 |
+
|
364 |
+
|
365 |
+
|
366 |
+
|
367 |
+
{\displaystyle {\hat {\mathbf {r} }}=\cos(\theta ){\hat {\mathbf {x} }}+\sin(\theta ){\hat {\mathbf {y} }}}
|
368 |
+
|
369 |
+
and
|
370 |
+
|
371 |
+
|
372 |
+
|
373 |
+
|
374 |
+
|
375 |
+
|
376 |
+
θ
|
377 |
+
^
|
378 |
+
|
379 |
+
|
380 |
+
|
381 |
+
=
|
382 |
+
−
|
383 |
+
sin
|
384 |
+
|
385 |
+
(
|
386 |
+
θ
|
387 |
+
)
|
388 |
+
|
389 |
+
|
390 |
+
|
391 |
+
|
392 |
+
x
|
393 |
+
|
394 |
+
^
|
395 |
+
|
396 |
+
|
397 |
+
|
398 |
+
+
|
399 |
+
cos
|
400 |
+
|
401 |
+
(
|
402 |
+
θ
|
403 |
+
)
|
404 |
+
|
405 |
+
|
406 |
+
|
407 |
+
|
408 |
+
y
|
409 |
+
|
410 |
+
^
|
411 |
+
|
412 |
+
|
413 |
+
|
414 |
+
|
415 |
+
|
416 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}=-\sin(\theta ){\hat {\mathbf {x} }}+\cos(\theta ){\hat {\mathbf {y} }}}
|
417 |
+
|
418 |
+
be the radial and transverse polar basis with the first being the unit vector pointing from the central body to the current location of the orbiting object and the second being the orthogonal unit vector pointing in the direction that the orbiting object would travel if orbiting in a counter clockwise circle. Then the vector to the orbiting object is
|
419 |
+
|
420 |
+
We use
|
421 |
+
|
422 |
+
|
423 |
+
|
424 |
+
|
425 |
+
|
426 |
+
|
427 |
+
r
|
428 |
+
˙
|
429 |
+
|
430 |
+
|
431 |
+
|
432 |
+
|
433 |
+
|
434 |
+
{\displaystyle {\dot {r}}}
|
435 |
+
|
436 |
+
and
|
437 |
+
|
438 |
+
|
439 |
+
|
440 |
+
|
441 |
+
|
442 |
+
|
443 |
+
θ
|
444 |
+
˙
|
445 |
+
|
446 |
+
|
447 |
+
|
448 |
+
|
449 |
+
|
450 |
+
{\displaystyle {\dot {\theta }}}
|
451 |
+
|
452 |
+
to denote the standard derivatives of how this distance and angle change over time. We take the derivative of a vector to see how it changes over time by subtracting its location at time
|
453 |
+
|
454 |
+
|
455 |
+
|
456 |
+
|
457 |
+
t
|
458 |
+
|
459 |
+
|
460 |
+
{\displaystyle t}
|
461 |
+
|
462 |
+
from that at time
|
463 |
+
|
464 |
+
|
465 |
+
|
466 |
+
t
|
467 |
+
+
|
468 |
+
δ
|
469 |
+
t
|
470 |
+
|
471 |
+
|
472 |
+
{\displaystyle t+\delta t}
|
473 |
+
|
474 |
+
and dividing by
|
475 |
+
|
476 |
+
|
477 |
+
|
478 |
+
δ
|
479 |
+
t
|
480 |
+
|
481 |
+
|
482 |
+
{\displaystyle \delta t}
|
483 |
+
|
484 |
+
. The result is also a vector. Because our basis vector
|
485 |
+
|
486 |
+
|
487 |
+
|
488 |
+
|
489 |
+
|
490 |
+
|
491 |
+
|
492 |
+
r
|
493 |
+
|
494 |
+
^
|
495 |
+
|
496 |
+
|
497 |
+
|
498 |
+
|
499 |
+
|
500 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
501 |
+
|
502 |
+
moves as the object orbits, we start by differentiating it.
|
503 |
+
From time
|
504 |
+
|
505 |
+
|
506 |
+
|
507 |
+
t
|
508 |
+
|
509 |
+
|
510 |
+
{\displaystyle t}
|
511 |
+
|
512 |
+
to
|
513 |
+
|
514 |
+
|
515 |
+
|
516 |
+
t
|
517 |
+
+
|
518 |
+
δ
|
519 |
+
t
|
520 |
+
|
521 |
+
|
522 |
+
{\displaystyle t+\delta t}
|
523 |
+
|
524 |
+
,
|
525 |
+
the vector
|
526 |
+
|
527 |
+
|
528 |
+
|
529 |
+
|
530 |
+
|
531 |
+
|
532 |
+
|
533 |
+
r
|
534 |
+
|
535 |
+
^
|
536 |
+
|
537 |
+
|
538 |
+
|
539 |
+
|
540 |
+
|
541 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
542 |
+
|
543 |
+
keeps its beginning at the origin and rotates from
|
544 |
+
angle
|
545 |
+
|
546 |
+
|
547 |
+
|
548 |
+
θ
|
549 |
+
|
550 |
+
|
551 |
+
{\displaystyle \theta }
|
552 |
+
|
553 |
+
to
|
554 |
+
|
555 |
+
|
556 |
+
|
557 |
+
θ
|
558 |
+
+
|
559 |
+
|
560 |
+
|
561 |
+
|
562 |
+
θ
|
563 |
+
˙
|
564 |
+
|
565 |
+
|
566 |
+
|
567 |
+
|
568 |
+
δ
|
569 |
+
t
|
570 |
+
|
571 |
+
|
572 |
+
{\displaystyle \theta +{\dot {\theta }}\ \delta t}
|
573 |
+
|
574 |
+
which moves its head a distance
|
575 |
+
|
576 |
+
|
577 |
+
|
578 |
+
|
579 |
+
|
580 |
+
|
581 |
+
θ
|
582 |
+
˙
|
583 |
+
|
584 |
+
|
585 |
+
|
586 |
+
|
587 |
+
δ
|
588 |
+
t
|
589 |
+
|
590 |
+
|
591 |
+
{\displaystyle {\dot {\theta }}\ \delta t}
|
592 |
+
|
593 |
+
in the perpendicular direction
|
594 |
+
|
595 |
+
|
596 |
+
|
597 |
+
|
598 |
+
|
599 |
+
|
600 |
+
θ
|
601 |
+
^
|
602 |
+
|
603 |
+
|
604 |
+
|
605 |
+
|
606 |
+
|
607 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}}
|
608 |
+
|
609 |
+
giving a derivative of
|
610 |
+
|
611 |
+
|
612 |
+
|
613 |
+
|
614 |
+
|
615 |
+
|
616 |
+
θ
|
617 |
+
˙
|
618 |
+
|
619 |
+
|
620 |
+
|
621 |
+
|
622 |
+
|
623 |
+
|
624 |
+
θ
|
625 |
+
^
|
626 |
+
|
627 |
+
|
628 |
+
|
629 |
+
|
630 |
+
|
631 |
+
{\displaystyle {\dot {\theta }}{\hat {\boldsymbol {\theta }}}}
|
632 |
+
|
633 |
+
.
|
634 |
+
|
635 |
+
We can now find the velocity and acceleration of our orbiting object.
|
636 |
+
|
637 |
+
The coefficients of
|
638 |
+
|
639 |
+
|
640 |
+
|
641 |
+
|
642 |
+
|
643 |
+
|
644 |
+
|
645 |
+
r
|
646 |
+
|
647 |
+
^
|
648 |
+
|
649 |
+
|
650 |
+
|
651 |
+
|
652 |
+
|
653 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
654 |
+
|
655 |
+
|
656 |
+
and
|
657 |
+
|
658 |
+
|
659 |
+
|
660 |
+
|
661 |
+
|
662 |
+
|
663 |
+
θ
|
664 |
+
^
|
665 |
+
|
666 |
+
|
667 |
+
|
668 |
+
|
669 |
+
|
670 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}}
|
671 |
+
|
672 |
+
give the accelerations in the radial and transverse directions.
|
673 |
+
As said, Newton gives this first due to gravity is
|
674 |
+
|
675 |
+
|
676 |
+
|
677 |
+
−
|
678 |
+
μ
|
679 |
+
|
680 |
+
/
|
681 |
+
|
682 |
+
|
683 |
+
r
|
684 |
+
|
685 |
+
2
|
686 |
+
|
687 |
+
|
688 |
+
|
689 |
+
|
690 |
+
{\displaystyle -\mu /r^{2}}
|
691 |
+
|
692 |
+
and the second is zero.
|
693 |
+
|
694 |
+
|
695 |
+
|
696 |
+
|
697 |
+
|
698 |
+
|
699 |
+
|
700 |
+
|
701 |
+
|
702 |
+
(1)
|
703 |
+
|
704 |
+
|
705 |
+
|
706 |
+
|
707 |
+
|
708 |
+
|
709 |
+
|
710 |
+
|
711 |
+
|
712 |
+
(2)
|
713 |
+
|
714 |
+
Equation (2) can be rearranged using integration by parts.
|
715 |
+
|
716 |
+
We can multiply through by
|
717 |
+
|
718 |
+
|
719 |
+
|
720 |
+
r
|
721 |
+
|
722 |
+
|
723 |
+
{\displaystyle r}
|
724 |
+
|
725 |
+
because it is not zero unless the orbiting object crashes.
|
726 |
+
Then having the derivative be zero gives that the function is a constant.
|
727 |
+
|
728 |
+
|
729 |
+
|
730 |
+
|
731 |
+
|
732 |
+
|
733 |
+
|
734 |
+
|
735 |
+
|
736 |
+
(3)
|
737 |
+
|
738 |
+
which is actually the theoretical proof of Kepler's second law (A line joining a planet and the Sun sweeps out equal areas during equal intervals of time). The constant of integration, h, is the angular momentum per unit mass.
|
739 |
+
|
740 |
+
In order to get an equation for the orbit from equation (1), we need to eliminate time.[8] (See also Binet equation.)
|
741 |
+
In polar coordinates, this would express the distance
|
742 |
+
|
743 |
+
|
744 |
+
|
745 |
+
r
|
746 |
+
|
747 |
+
|
748 |
+
{\displaystyle r}
|
749 |
+
|
750 |
+
of the orbiting object from the center as a function of its angle
|
751 |
+
|
752 |
+
|
753 |
+
|
754 |
+
θ
|
755 |
+
|
756 |
+
|
757 |
+
{\displaystyle \theta }
|
758 |
+
|
759 |
+
. However, it is easier to
|
760 |
+
introduce the auxiliary variable
|
761 |
+
|
762 |
+
|
763 |
+
|
764 |
+
u
|
765 |
+
=
|
766 |
+
1
|
767 |
+
|
768 |
+
/
|
769 |
+
|
770 |
+
r
|
771 |
+
|
772 |
+
|
773 |
+
{\displaystyle u=1/r}
|
774 |
+
|
775 |
+
and to express
|
776 |
+
|
777 |
+
|
778 |
+
|
779 |
+
u
|
780 |
+
|
781 |
+
|
782 |
+
{\displaystyle u}
|
783 |
+
|
784 |
+
as a function of
|
785 |
+
|
786 |
+
|
787 |
+
|
788 |
+
θ
|
789 |
+
|
790 |
+
|
791 |
+
{\displaystyle \theta }
|
792 |
+
|
793 |
+
. Derivatives of
|
794 |
+
|
795 |
+
|
796 |
+
|
797 |
+
r
|
798 |
+
|
799 |
+
|
800 |
+
{\displaystyle r}
|
801 |
+
|
802 |
+
with respect to time may be rewritten as derivatives of
|
803 |
+
|
804 |
+
|
805 |
+
|
806 |
+
u
|
807 |
+
|
808 |
+
|
809 |
+
{\displaystyle u}
|
810 |
+
|
811 |
+
with respect to angle.
|
812 |
+
|
813 |
+
Plugging these into (1) gives
|
814 |
+
|
815 |
+
So for the gravitational force – or, more generally, for any inverse square force law – the right hand side of the equation becomes a constant and the equation is seen to be the harmonic equation (up to a shift of origin of the dependent variable). The solution is:
|
816 |
+
|
817 |
+
where A and θ0 are arbitrary constants.
|
818 |
+
This resulting equation of the orbit of the object is that of an ellipse in Polar form relative to one of the focal points. This is put into a more standard form by
|
819 |
+
letting
|
820 |
+
|
821 |
+
|
822 |
+
|
823 |
+
e
|
824 |
+
≡
|
825 |
+
|
826 |
+
h
|
827 |
+
|
828 |
+
2
|
829 |
+
|
830 |
+
|
831 |
+
A
|
832 |
+
|
833 |
+
/
|
834 |
+
|
835 |
+
μ
|
836 |
+
|
837 |
+
|
838 |
+
{\displaystyle e\equiv h^{2}A/\mu }
|
839 |
+
|
840 |
+
be the eccentricity,
|
841 |
+
letting
|
842 |
+
|
843 |
+
|
844 |
+
|
845 |
+
a
|
846 |
+
≡
|
847 |
+
|
848 |
+
h
|
849 |
+
|
850 |
+
2
|
851 |
+
|
852 |
+
|
853 |
+
|
854 |
+
/
|
855 |
+
|
856 |
+
(
|
857 |
+
μ
|
858 |
+
(
|
859 |
+
1
|
860 |
+
−
|
861 |
+
|
862 |
+
e
|
863 |
+
|
864 |
+
2
|
865 |
+
|
866 |
+
|
867 |
+
)
|
868 |
+
)
|
869 |
+
|
870 |
+
|
871 |
+
{\displaystyle a\equiv h^{2}/(\mu (1-e^{2}))}
|
872 |
+
|
873 |
+
be the semi-major axis.
|
874 |
+
Finally, letting
|
875 |
+
|
876 |
+
|
877 |
+
|
878 |
+
|
879 |
+
θ
|
880 |
+
|
881 |
+
0
|
882 |
+
|
883 |
+
|
884 |
+
≡
|
885 |
+
0
|
886 |
+
|
887 |
+
|
888 |
+
{\displaystyle \theta _{0}\equiv 0}
|
889 |
+
|
890 |
+
so the long axis of the ellipse is along the positive x coordinate.
|
891 |
+
|
892 |
+
The above classical (Newtonian) analysis of orbital mechanics assumes that the more subtle effects of general relativity, such as frame dragging and gravitational time dilation are negligible. Relativistic effects cease to be negligible when near very massive bodies (as with the precession of Mercury's orbit about the Sun), or when extreme precision is needed (as with calculations of the orbital elements and time signal references for GPS satellites.[9]).
|
893 |
+
|
894 |
+
The analysis so far has been two dimensional; it turns out that an unperturbed orbit is two-dimensional in a plane fixed in space, and thus the extension to three dimensions requires simply rotating the two-dimensional plane into the required angle relative to the poles of the planetary body involved.
|
895 |
+
|
896 |
+
The rotation to do this in three dimensions requires three numbers to uniquely determine; traditionally these are expressed as three angles.
|
897 |
+
|
898 |
+
The orbital period is simply how long an orbiting body takes to complete one orbit.
|
899 |
+
|
900 |
+
Six parameters are required to specify a Keplerian orbit about a body. For example, the three numbers that specify the body's initial position, and the three values that specify its velocity will define a unique orbit that can be calculated forwards (or backwards) in time. However, traditionally the parameters used are slightly different.
|
901 |
+
|
902 |
+
The traditionally used set of orbital elements is called the set of Keplerian elements, after Johannes Kepler and his laws. The Keplerian elements are six:
|
903 |
+
|
904 |
+
In principle once the orbital elements are known for a body, its position can be calculated forward and backwards indefinitely in time. However, in practice, orbits are affected or perturbed, by other forces than simple gravity from an assumed point source (see the next section), and thus the orbital elements change over time.
|
905 |
+
|
906 |
+
An orbital perturbation is when a force or impulse which is much smaller than the overall force or average impulse of the main gravitating body and which is external to the two orbiting bodies causes an acceleration, which changes the parameters of the orbit over time.
|
907 |
+
|
908 |
+
A small radial impulse given to a body in orbit changes the eccentricity, but not the orbital period (to first order). A prograde or retrograde impulse (i.e. an impulse applied along the orbital motion) changes both the eccentricity and the orbital period. Notably, a prograde impulse at periapsis raises the altitude at apoapsis, and vice versa, and a retrograde impulse does the opposite. A transverse impulse (out of the orbital plane) causes rotation of the orbital plane without changing the period or eccentricity. In all instances, a closed orbit will still intersect the perturbation point.
|
909 |
+
|
910 |
+
If an orbit is about a planetary body with significant atmosphere, its orbit can decay because of drag. Particularly at each periapsis, the object experiences atmospheric drag, losing energy. Each time, the orbit grows less eccentric (more circular) because the object loses kinetic energy precisely when that energy is at its maximum. This is similar to the effect of slowing a pendulum at its lowest point; the highest point of the pendulum's swing becomes lower. With each successive slowing more of the orbit's path is affected by the atmosphere and the effect becomes more pronounced. Eventually, the effect becomes so great that the maximum kinetic energy is not enough to return the orbit above the limits of the atmospheric drag effect. When this happens the body will rapidly spiral down and intersect the central body.
|
911 |
+
|
912 |
+
The bounds of an atmosphere vary wildly. During a solar maximum, the Earth's atmosphere causes drag up to a hundred kilometres higher than during a solar minimum.
|
913 |
+
|
914 |
+
Some satellites with long conductive tethers can also experience orbital decay because of electromagnetic drag from the Earth's magnetic field. As the wire cuts the magnetic field it acts as a generator, moving electrons from one end to the other. The orbital energy is converted to heat in the wire.
|
915 |
+
|
916 |
+
Orbits can be artificially influenced through the use of rocket engines which change the kinetic energy of the body at some point in its path. This is the conversion of chemical or electrical energy to kinetic energy. In this way changes in the orbit shape or orientation can be facilitated.
|
917 |
+
|
918 |
+
Another method of artificially influencing an orbit is through the use of solar sails or magnetic sails. These forms of propulsion require no propellant or energy input other than that of the Sun, and so can be used indefinitely. See statite for one such proposed use.
|
919 |
+
|
920 |
+
Orbital decay can occur due to tidal forces for objects below the synchronous orbit for the body they're orbiting. The gravity of the orbiting object raises tidal bulges in the primary, and since below the synchronous orbit the orbiting object is moving faster than the body's surface the bulges lag a short angle behind it. The gravity of the bulges is slightly off of the primary-satellite axis and thus has a component along the satellite's motion. The near bulge slows the object more than the far bulge speeds it up, and as a result the orbit decays. Conversely, the gravity of the satellite on the bulges applies torque on the primary and speeds up its rotation. Artificial satellites are too small to have an appreciable tidal effect on the planets they orbit, but several moons in the Solar System are undergoing orbital decay by this mechanism. Mars' innermost moon Phobos is a prime example, and is expected to either impact Mars' surface or break up into a ring within 50 million years.
|
921 |
+
|
922 |
+
Orbits can decay via the emission of gravitational waves. This mechanism is extremely weak for most stellar objects, only becoming significant in cases where there is a combination of extreme mass and extreme acceleration, such as with black holes or neutron stars that are orbiting each other closely.
|
923 |
+
|
924 |
+
The standard analysis of orbiting bodies assumes that all bodies consist of uniform spheres, or more generally, concentric shells each of uniform density. It can be shown that such bodies are gravitationally equivalent to point sources.
|
925 |
+
|
926 |
+
However, in the real world, many bodies rotate, and this introduces oblateness and distorts the gravity field, and gives a quadrupole moment to the gravitational field which is significant at distances comparable to the radius of the body. In the general case, the gravitational potential of a rotating body such as, e.g., a planet is usually expanded in multipoles accounting for the departures of it from spherical symmetry. From the point of view of satellite dynamics, of particular relevance are the so-called even zonal harmonic coefficients, or even zonals, since they induce secular orbital perturbations which are cumulative over time spans longer than the orbital period.[10][11][12] They do depend on the orientation of the body's symmetry axis in the space, affecting, in general, the whole orbit, with the exception of the semimajor axis.
|
927 |
+
|
928 |
+
The effects of other gravitating bodies can be significant. For example, the orbit of the Moon cannot be accurately described without allowing for the action of the Sun's gravity as well as the Earth's. One approximate result is that bodies will usually have reasonably stable orbits around a heavier planet or moon, in spite of these perturbations, provided they are orbiting well within the heavier body's Hill sphere.
|
929 |
+
|
930 |
+
When there are more than two gravitating bodies it is referred to as an n-body problem. Most n-body problems have no closed form solution, although some special cases have been formulated.
|
931 |
+
|
932 |
+
For smaller bodies particularly, light and stellar wind can cause significant perturbations to the attitude and direction of motion of the body, and over time can be significant. Of the planetary bodies, the motion of asteroids is particularly affected over large periods when the asteroids are rotating relative to the Sun.
|
933 |
+
|
934 |
+
Mathematicians have discovered that it is possible in principle to have multiple bodies in non-elliptical orbits that repeat periodically, although most such orbits are not stable regarding small perturbations in mass, position, or velocity. However, some special stable cases have been identified, including a planar figure-eight orbit occupied by three moving bodies. Further studies have discovered that nonplanar orbits are also possible, including one involving 12 masses moving in 4 roughly circular, interlocking orbits topologically equivalent to the edges of a cuboctahedron.[13]
|
935 |
+
|
936 |
+
Finding such orbits naturally occurring in the universe is thought to be extremely unlikely, because of the improbability of the required conditions occurring by chance.[13]
|
937 |
+
|
938 |
+
Orbital mechanics or astrodynamics is the application of ballistics and celestial mechanics to the practical problems concerning the motion of rockets and other spacecraft. The motion of these objects is usually calculated from Newton's laws of motion and Newton's law of universal gravitation. It is a core discipline within space mission design and control. Celestial mechanics treats more broadly the orbital dynamics of systems under the influence of gravity, including spacecraft and natural astronomical bodies such as star systems, planets, moons, and comets. Orbital mechanics focuses on spacecraft trajectories, including orbital maneuvers, orbit plane changes, and interplanetary transfers, and is used by mission planners to predict the results of propulsive maneuvers. General relativity is a more exact theory than Newton's laws for calculating orbits, and is sometimes necessary for greater accuracy or in high-gravity situations (such as orbits close to the Sun).
|
939 |
+
|
940 |
+
The gravitational constant G has been calculated as:
|
941 |
+
|
942 |
+
Thus the constant has dimension density−1 time−2. This corresponds to the following properties.
|
943 |
+
|
944 |
+
Scaling of distances (including sizes of bodies, while keeping the densities the same) gives similar orbits without scaling the time: if for example distances are halved, masses are divided by 8, gravitational forces by 16 and gravitational accelerations by 2. Hence velocities are halved and orbital periods remain the same. Similarly, when an object is dropped from a tower, the time it takes to fall to the ground remains the same with a scale model of the tower on a scale model of the Earth.
|
945 |
+
|
946 |
+
Scaling of distances while keeping the masses the same (in the case of point masses, or by reducing the densities) gives similar orbits; if distances are multiplied by 4, gravitational forces and accelerations are divided by 16, velocities are halved and orbital periods are multiplied by 8.
|
947 |
+
|
948 |
+
When all densities are multiplied by 4, orbits are the same; gravitational forces are multiplied by 16 and accelerations by 4, velocities are doubled and orbital periods are halved.
|
949 |
+
|
950 |
+
When all densities are multiplied by 4, and all sizes are halved, orbits are similar; masses are divided by 2, gravitational forces are the same, gravitational accelerations are doubled. Hence velocities are the same and orbital periods are halved.
|
951 |
+
|
952 |
+
In all these cases of scaling. if densities are multiplied by 4, times are halved; if velocities are doubled, forces are multiplied by 16.
|
953 |
+
|
954 |
+
These properties are illustrated in the formula (derived from the formula for the orbital period)
|
955 |
+
|
956 |
+
for an elliptical orbit with semi-major axis a, of a small body around a spherical body with radius r and average density ρ, where T is the orbital period. See also Kepler's Third Law.
|
957 |
+
|
958 |
+
The application of certain orbits or orbital maneuvers to specific useful purposes have been the subject of patents.[17]
|
959 |
+
|
960 |
+
Some bodies are tidally locked with other bodies, meaning that one side of the celestial body is permanently facing its host object. This is the case for Earth-Moon and Pluto-Charon system.
|
961 |
+
|
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1 |
+
|
2 |
+
|
3 |
+
In physics, an orbit is the gravitationally curved trajectory of an object,[1] such as the trajectory of a planet around a star or a natural satellite around a planet. Normally, orbit refers to a regularly repeating trajectory, although it may also refer to a non-repeating trajectory. To a close approximation, planets and satellites follow elliptic orbits, with the center of mass being orbited at a focal point of the ellipse,[2] as described by Kepler's laws of planetary motion.
|
4 |
+
|
5 |
+
For most situations, orbital motion is adequately approximated by Newtonian mechanics, which explains gravity as a force obeying an inverse-square law.[3] However, Albert Einstein's general theory of relativity, which accounts for gravity as due to curvature of spacetime, with orbits following geodesics, provides a more accurate calculation and understanding of the exact mechanics of orbital motion.
|
6 |
+
|
7 |
+
Historically, the apparent motions of the planets were described by European and Arabic philosophers using the idea of celestial spheres. This model posited the existence of perfect moving spheres or rings to which the stars and planets were attached. It assumed the heavens were fixed apart from the motion of the spheres, and was developed without any understanding of gravity. After the planets' motions were more accurately measured, theoretical mechanisms such as deferent and epicycles were added. Although the model was capable of reasonably accurately predicting the planets' positions in the sky, more and more epicycles were required as the measurements became more accurate, hence the model became increasingly unwieldy. Originally geocentric, it was modified by Copernicus to place the Sun at the centre to help simplify the model. The model was further challenged during the 16th century, as comets were observed traversing the spheres.[4][5]
|
8 |
+
|
9 |
+
The basis for the modern understanding of orbits was first formulated by Johannes Kepler whose results are summarised in his three laws of planetary motion. First, he found that the orbits of the planets in our Solar System are elliptical, not circular (or epicyclic), as had previously been believed, and that the Sun is not located at the center of the orbits, but rather at one focus.[6] Second, he found that the orbital speed of each planet is not constant, as had previously been thought, but rather that the speed depends on the planet's distance from the Sun. Third, Kepler found a universal relationship between the orbital properties of all the planets orbiting the Sun. For the planets, the cubes of their distances from the Sun are proportional to the squares of their orbital periods. Jupiter and Venus, for example, are respectively about 5.2 and 0.723 AU distant from the Sun, their orbital periods respectively about 11.86 and 0.615 years. The proportionality is seen by the fact that the ratio for Jupiter, 5.23/11.862, is practically equal to that for Venus, 0.7233/0.6152, in accord with the relationship. Idealised orbits meeting these rules are known as Kepler orbits.
|
10 |
+
|
11 |
+
Isaac Newton demonstrated that Kepler's laws were derivable from his theory of gravitation and that, in general, the orbits of bodies subject to gravity were conic sections (this assumes that the force of gravity propagates instantaneously). Newton showed that, for a pair of bodies, the orbits' sizes are in inverse proportion to their masses, and that those bodies orbit their common center of mass. Where one body is much more massive than the other (as is the case of an artificial satellite orbiting a planet), it is a convenient approximation to take the center of mass as coinciding with the center of the more massive body.
|
12 |
+
|
13 |
+
Advances in Newtonian mechanics were then used to explore variations from the simple assumptions behind Kepler orbits, such as the perturbations due to other bodies, or the impact of spheroidal rather than spherical bodies. Lagrange (1736–1813) developed a new approach to Newtonian mechanics emphasizing energy more than force, and made progress on the three body problem, discovering the Lagrangian points. In a dramatic vindication of classical mechanics, in 1846 Urbain Le Verrier was able to predict the position of Neptune based on unexplained perturbations in the orbit of Uranus.
|
14 |
+
|
15 |
+
Albert Einstein (1879-1955) in his 1916 paper The Foundation of the General Theory of Relativity explained that gravity was due to curvature of space-time and removed Newton's assumption that changes propagate instantaneously. This led astronomers to recognize that Newtonian mechanics did not provide the highest accuracy in understanding orbits. In relativity theory, orbits follow geodesic trajectories which are usually approximated very well by the Newtonian predictions (except where there are very strong gravity fields and very high speeds) but the differences are measurable. Essentially all the experimental evidence that can distinguish between the theories agrees with relativity theory to within experimental measurement accuracy. The original vindication of general relativity is that it was able to account for the remaining unexplained amount in precession of Mercury's perihelion first noted by Le Verrier. However, Newton's solution is still used for most short term purposes since it is significantly easier to use and sufficiently accurate.
|
16 |
+
|
17 |
+
Within a planetary system, planets, dwarf planets, asteroids and other minor planets, comets, and space debris orbit the system's barycenter in elliptical orbits. A comet in a parabolic or hyperbolic orbit about a barycenter is not gravitationally bound to the star and therefore is not considered part of the star's planetary system. Bodies which are gravitationally bound to one of the planets in a planetary system, either natural or artificial satellites, follow orbits about a barycenter near or within that planet.
|
18 |
+
|
19 |
+
Owing to mutual gravitational perturbations, the eccentricities of the planetary orbits vary over time. Mercury, the smallest planet in the Solar System, has the most eccentric orbit. At the present epoch, Mars has the next largest eccentricity while the smallest orbital eccentricities are seen with Venus and Neptune.
|
20 |
+
|
21 |
+
As two objects orbit each other, the periapsis is that point at which the two objects are closest to each other and the apoapsis is that point at which they are the farthest. (More specific terms are used for specific bodies. For example, perigee and apogee are the lowest and highest parts of an orbit around Earth, while perihelion and aphelion are the closest and farthest points of an orbit around the Sun.)
|
22 |
+
|
23 |
+
In the case of planets orbiting a star, the mass of the star and all its satellites are calculated to be at a single point called the barycenter. The paths of all the star's satellites are elliptical orbits about that barycenter. Each satellite in that system will have its own elliptical orbit with the barycenter at one focal point of that ellipse. At any point along its orbit, any satellite will have a certain value of kinetic and potential energy with respect to the barycenter, and that energy is a constant value at every point along its orbit. As a result, as a planet approaches periapsis, the planet will increase in speed as its potential energy decreases; as a planet approaches apoapsis, its velocity will decrease as its potential energy increases.
|
24 |
+
|
25 |
+
There are a few common ways of understanding orbits:
|
26 |
+
|
27 |
+
As an illustration of an orbit around a planet, the Newton's cannonball model may prove useful (see image below). This is a 'thought experiment', in which a cannon on top of a tall mountain is able to fire a cannonball horizontally at any chosen muzzle speed. The effects of air friction on the cannonball are ignored (or perhaps the mountain is high enough that the cannon is above the Earth's atmosphere, which is the same thing).[7]
|
28 |
+
|
29 |
+
If the cannon fires its ball with a low initial speed, the trajectory of the ball curves downward and hits the ground (A). As the firing speed is increased, the cannonball hits the ground farther (B) away from the cannon, because while the ball is still falling towards the ground, the ground is increasingly curving away from it (see first point, above). All these motions are actually "orbits" in a technical sense – they are describing a portion of an elliptical path around the center of gravity – but the orbits are interrupted by striking the Earth.
|
30 |
+
|
31 |
+
If the cannonball is fired with sufficient speed, the ground curves away from the ball at least as much as the ball falls – so the ball never strikes the ground. It is now in what could be called a non-interrupted, or circumnavigating, orbit. For any specific combination of height above the center of gravity and mass of the planet, there is one specific firing speed (unaffected by the mass of the ball, which is assumed to be very small relative to the Earth's mass) that produces a circular orbit, as shown in (C).
|
32 |
+
|
33 |
+
As the firing speed is increased beyond this, non-interrupted elliptic orbits are produced; one is shown in (D). If the initial firing is above the surface of the Earth as shown, there will also be non-interrupted elliptical orbits at slower firing speed; these will come closest to the Earth at the point half an orbit beyond, and directly opposite the firing point, below the circular orbit.
|
34 |
+
|
35 |
+
At a specific horizontal firing speed called escape velocity, dependent on the mass of the planet, an open orbit (E) is achieved that has a parabolic path. At even greater speeds the object will follow a range of hyperbolic trajectories. In a practical sense, both of these trajectory types mean the object is "breaking free" of the planet's gravity, and "going off into space" never to return.
|
36 |
+
|
37 |
+
The velocity relationship of two moving objects with mass can thus be considered in four practical classes, with subtypes:
|
38 |
+
|
39 |
+
It is worth noting that orbital rockets are launched vertically at first to lift the rocket above the atmosphere (which causes frictional drag), and then slowly pitch over and finish firing the rocket engine parallel to the atmosphere to achieve orbit speed.
|
40 |
+
|
41 |
+
Once in orbit, their speed keeps them in orbit above the atmosphere. If e.g., an elliptical orbit dips into dense air, the object will lose speed and re-enter (i.e. fall). Occasionally a space craft will intentionally intercept the atmosphere, in an act commonly referred to as an aerobraking maneuver.
|
42 |
+
|
43 |
+
In most situations relativistic effects can be neglected, and Newton's laws give a sufficiently accurate description of motion. The acceleration of a body is equal to the sum of the forces acting on it, divided by its mass, and the gravitational force acting on a body is proportional to the product of the masses of the two attracting bodies and decreases inversely with the square of the distance between them. To this Newtonian approximation, for a system of two-point masses or spherical bodies, only influenced by their mutual gravitation (called a two-body problem), their trajectories can be exactly calculated. If the heavier body is much more massive than the smaller, as in the case of a satellite or small moon orbiting a planet or for the Earth orbiting the Sun, it is accurate enough and convenient to describe the motion in terms of a coordinate system that is centered on the heavier body, and we say that the lighter body is in orbit around the heavier. For the case where the masses of two bodies are comparable, an exact Newtonian solution is still sufficient and can be had by placing the coordinate system at the center of mass of the system.
|
44 |
+
|
45 |
+
Energy is associated with gravitational fields. A stationary body far from another can do external work if it is pulled towards it, and therefore has gravitational potential energy. Since work is required to separate two bodies against the pull of gravity, their gravitational potential energy increases as they are separated, and decreases as they approach one another. For point masses the gravitational energy decreases to zero as they approach zero separation. It is convenient and conventional to assign the potential energy as having zero value when they are an infinite distance apart, and hence it has a negative value (since it decreases from zero) for smaller finite distances.
|
46 |
+
|
47 |
+
When only two gravitational bodies interact, their orbits follow a conic section. The orbit can be open (implying the object never returns) or closed (returning). Which it is depends on the total energy (kinetic + potential energy) of the system. In the case of an open orbit, the speed at any position of the orbit is at least the escape velocity for that position, in the case of a closed orbit, the speed is always less than the escape velocity. Since the kinetic energy is never negative, if the common convention is adopted of taking the potential energy as zero at infinite separation, the bound orbits will have negative total energy, the parabolic trajectories zero total energy, and hyperbolic orbits positive total energy.
|
48 |
+
|
49 |
+
An open orbit will have a parabolic shape if it has velocity of exactly the escape velocity at that point in its trajectory, and it will have the shape of a hyperbola when its velocity is greater than the escape velocity. When bodies with escape velocity or greater approach each other, they will briefly curve around each other at the time of their closest approach, and then separate, forever.
|
50 |
+
|
51 |
+
All closed orbits have the shape of an ellipse. A circular orbit is a special case, wherein the foci of the ellipse coincide. The point where the orbiting body is closest to Earth is called the perigee, and is called the periapsis (less properly, "perifocus" or "pericentron") when the orbit is about a body other than Earth. The point where the satellite is farthest from Earth is called the apogee, apoapsis, or sometimes apifocus or apocentron. A line drawn from periapsis to apoapsis is the line-of-apsides. This is the major axis of the ellipse, the line through its longest part.
|
52 |
+
|
53 |
+
Bodies following closed orbits repeat their paths with a certain time called the period. This motion is described by the empirical laws of Kepler, which can be mathematically derived from Newton's laws. These can be
|
54 |
+
formulated as follows:
|
55 |
+
|
56 |
+
Note that while bound orbits of a point mass or a spherical body with a Newtonian gravitational field are closed ellipses, which repeat the same path exactly and indefinitely, any non-spherical or non-Newtonian effects (such as caused by the slight oblateness of the Earth, or by relativistic effects, thereby changing the gravitational field's behavior with distance) will cause the orbit's shape to depart from the closed ellipses characteristic of Newtonian two-body motion. The two-body solutions were published by Newton in Principia in 1687. In 1912, Karl Fritiof Sundman developed a converging infinite series that solves the three-body problem; however, it converges too slowly to be of much use. Except for special cases like the Lagrangian points, no method is known to solve the equations of motion for a system with four or more bodies.
|
57 |
+
|
58 |
+
Rather than an exact closed form solution, orbits with many bodies can be approximated with arbitrarily high accuracy. These approximations take two forms:
|
59 |
+
|
60 |
+
Differential simulations with large numbers of objects perform the calculations in a hierarchical pairwise fashion between centers of mass. Using this scheme, galaxies, star clusters and other large assemblages of objects have been simulated.[citation needed]
|
61 |
+
|
62 |
+
The Earth follows an ellipse round the sun.
|
63 |
+
But unlike the ellipse followed by a pendulum or an object attached to a spring, the sun is at a focal point of the ellipse and not at its centre.
|
64 |
+
|
65 |
+
The following derivation applies to such an elliptical orbit.
|
66 |
+
We start only with the Newtonian law of gravitation stating that the gravitational acceleration towards the central body is related to the inverse of the square of the distance between them, namely
|
67 |
+
|
68 |
+
where F2 is the force acting on the mass m2 caused by the gravitational attraction mass m1 has for m2, G is the universal gravitational constant, and r is the distance between the two masses centers.
|
69 |
+
|
70 |
+
From Newton's Second Law, the summation of the forces acting on m2 related to that bodies acceleration:
|
71 |
+
|
72 |
+
where A2 is the acceleration of m2 caused by the force of gravitational attraction F2 of m1 acting on m2.
|
73 |
+
|
74 |
+
Combining Eq 1 and 2:
|
75 |
+
|
76 |
+
Solving for the acceleration, A2:
|
77 |
+
|
78 |
+
where
|
79 |
+
|
80 |
+
|
81 |
+
|
82 |
+
μ
|
83 |
+
|
84 |
+
|
85 |
+
|
86 |
+
{\displaystyle \mu \,}
|
87 |
+
|
88 |
+
is the standard gravitational parameter, in this case
|
89 |
+
|
90 |
+
|
91 |
+
|
92 |
+
G
|
93 |
+
|
94 |
+
m
|
95 |
+
|
96 |
+
1
|
97 |
+
|
98 |
+
|
99 |
+
|
100 |
+
|
101 |
+
{\displaystyle Gm_{1}}
|
102 |
+
|
103 |
+
. It is understood that the system being described is m2, hence the subscripts can be dropped.
|
104 |
+
|
105 |
+
We assume that the central body is massive enough that it can be considered to be stationary and we ignore the more subtle effects of general relativity.
|
106 |
+
|
107 |
+
When a pendulum or an object attached to a spring swings in an ellipse,
|
108 |
+
the inward acceleration/force is proportional to the distance
|
109 |
+
|
110 |
+
|
111 |
+
|
112 |
+
A
|
113 |
+
=
|
114 |
+
F
|
115 |
+
|
116 |
+
/
|
117 |
+
|
118 |
+
m
|
119 |
+
=
|
120 |
+
−
|
121 |
+
k
|
122 |
+
r
|
123 |
+
.
|
124 |
+
|
125 |
+
|
126 |
+
{\displaystyle A=F/m=-kr.}
|
127 |
+
|
128 |
+
|
129 |
+
Due to the way vectors add, the component of the force in the
|
130 |
+
|
131 |
+
|
132 |
+
|
133 |
+
|
134 |
+
|
135 |
+
|
136 |
+
|
137 |
+
x
|
138 |
+
|
139 |
+
^
|
140 |
+
|
141 |
+
|
142 |
+
|
143 |
+
|
144 |
+
|
145 |
+
{\displaystyle {\hat {\mathbf {x} }}}
|
146 |
+
|
147 |
+
or in the
|
148 |
+
|
149 |
+
|
150 |
+
|
151 |
+
|
152 |
+
|
153 |
+
|
154 |
+
|
155 |
+
y
|
156 |
+
|
157 |
+
^
|
158 |
+
|
159 |
+
|
160 |
+
|
161 |
+
|
162 |
+
|
163 |
+
{\displaystyle {\hat {\mathbf {y} }}}
|
164 |
+
|
165 |
+
directions are also proportionate to the respective
|
166 |
+
components of the distances,
|
167 |
+
|
168 |
+
|
169 |
+
|
170 |
+
|
171 |
+
r
|
172 |
+
|
173 |
+
x
|
174 |
+
|
175 |
+
″
|
176 |
+
|
177 |
+
=
|
178 |
+
|
179 |
+
A
|
180 |
+
|
181 |
+
x
|
182 |
+
|
183 |
+
|
184 |
+
=
|
185 |
+
−
|
186 |
+
k
|
187 |
+
|
188 |
+
r
|
189 |
+
|
190 |
+
x
|
191 |
+
|
192 |
+
|
193 |
+
|
194 |
+
|
195 |
+
{\displaystyle r''_{x}=A_{x}=-kr_{x}}
|
196 |
+
|
197 |
+
. Hence, the entire analysis can be done separately in these dimensions. This results in the harmonic parabolic equations
|
198 |
+
|
199 |
+
|
200 |
+
|
201 |
+
x
|
202 |
+
=
|
203 |
+
A
|
204 |
+
cos
|
205 |
+
|
206 |
+
(
|
207 |
+
t
|
208 |
+
)
|
209 |
+
|
210 |
+
|
211 |
+
{\displaystyle x=A\cos(t)}
|
212 |
+
|
213 |
+
and
|
214 |
+
|
215 |
+
|
216 |
+
|
217 |
+
y
|
218 |
+
=
|
219 |
+
B
|
220 |
+
sin
|
221 |
+
|
222 |
+
(
|
223 |
+
t
|
224 |
+
)
|
225 |
+
|
226 |
+
|
227 |
+
{\displaystyle y=B\sin(t)}
|
228 |
+
|
229 |
+
of the ellipse. In contrast, with the decreasing relationship
|
230 |
+
|
231 |
+
|
232 |
+
|
233 |
+
A
|
234 |
+
=
|
235 |
+
μ
|
236 |
+
|
237 |
+
/
|
238 |
+
|
239 |
+
|
240 |
+
r
|
241 |
+
|
242 |
+
2
|
243 |
+
|
244 |
+
|
245 |
+
|
246 |
+
|
247 |
+
{\displaystyle A=\mu /r^{2}}
|
248 |
+
|
249 |
+
, the dimensions cannot be separated.[citation needed]
|
250 |
+
|
251 |
+
The location of the orbiting object at the current time
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
t
|
256 |
+
|
257 |
+
|
258 |
+
{\displaystyle t}
|
259 |
+
|
260 |
+
is located in the plane using
|
261 |
+
Vector calculus in polar coordinates both with the standard Euclidean basis and with the polar basis
|
262 |
+
with the origin coinciding with the center of force.
|
263 |
+
Let
|
264 |
+
|
265 |
+
|
266 |
+
|
267 |
+
r
|
268 |
+
|
269 |
+
|
270 |
+
{\displaystyle r}
|
271 |
+
|
272 |
+
be the distance between the object and the center and
|
273 |
+
|
274 |
+
|
275 |
+
|
276 |
+
|
277 |
+
θ
|
278 |
+
|
279 |
+
|
280 |
+
{\displaystyle \theta }
|
281 |
+
|
282 |
+
be the angle it has rotated.
|
283 |
+
Let
|
284 |
+
|
285 |
+
|
286 |
+
|
287 |
+
|
288 |
+
|
289 |
+
|
290 |
+
|
291 |
+
x
|
292 |
+
|
293 |
+
^
|
294 |
+
|
295 |
+
|
296 |
+
|
297 |
+
|
298 |
+
|
299 |
+
{\displaystyle {\hat {\mathbf {x} }}}
|
300 |
+
|
301 |
+
and
|
302 |
+
|
303 |
+
|
304 |
+
|
305 |
+
|
306 |
+
|
307 |
+
|
308 |
+
|
309 |
+
y
|
310 |
+
|
311 |
+
^
|
312 |
+
|
313 |
+
|
314 |
+
|
315 |
+
|
316 |
+
|
317 |
+
{\displaystyle {\hat {\mathbf {y} }}}
|
318 |
+
|
319 |
+
be the standard Euclidean bases and let
|
320 |
+
|
321 |
+
|
322 |
+
|
323 |
+
|
324 |
+
|
325 |
+
|
326 |
+
|
327 |
+
r
|
328 |
+
|
329 |
+
^
|
330 |
+
|
331 |
+
|
332 |
+
|
333 |
+
=
|
334 |
+
cos
|
335 |
+
|
336 |
+
(
|
337 |
+
θ
|
338 |
+
)
|
339 |
+
|
340 |
+
|
341 |
+
|
342 |
+
|
343 |
+
x
|
344 |
+
|
345 |
+
^
|
346 |
+
|
347 |
+
|
348 |
+
|
349 |
+
+
|
350 |
+
sin
|
351 |
+
|
352 |
+
(
|
353 |
+
θ
|
354 |
+
)
|
355 |
+
|
356 |
+
|
357 |
+
|
358 |
+
|
359 |
+
y
|
360 |
+
|
361 |
+
^
|
362 |
+
|
363 |
+
|
364 |
+
|
365 |
+
|
366 |
+
|
367 |
+
{\displaystyle {\hat {\mathbf {r} }}=\cos(\theta ){\hat {\mathbf {x} }}+\sin(\theta ){\hat {\mathbf {y} }}}
|
368 |
+
|
369 |
+
and
|
370 |
+
|
371 |
+
|
372 |
+
|
373 |
+
|
374 |
+
|
375 |
+
|
376 |
+
θ
|
377 |
+
^
|
378 |
+
|
379 |
+
|
380 |
+
|
381 |
+
=
|
382 |
+
−
|
383 |
+
sin
|
384 |
+
|
385 |
+
(
|
386 |
+
θ
|
387 |
+
)
|
388 |
+
|
389 |
+
|
390 |
+
|
391 |
+
|
392 |
+
x
|
393 |
+
|
394 |
+
^
|
395 |
+
|
396 |
+
|
397 |
+
|
398 |
+
+
|
399 |
+
cos
|
400 |
+
|
401 |
+
(
|
402 |
+
θ
|
403 |
+
)
|
404 |
+
|
405 |
+
|
406 |
+
|
407 |
+
|
408 |
+
y
|
409 |
+
|
410 |
+
^
|
411 |
+
|
412 |
+
|
413 |
+
|
414 |
+
|
415 |
+
|
416 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}=-\sin(\theta ){\hat {\mathbf {x} }}+\cos(\theta ){\hat {\mathbf {y} }}}
|
417 |
+
|
418 |
+
be the radial and transverse polar basis with the first being the unit vector pointing from the central body to the current location of the orbiting object and the second being the orthogonal unit vector pointing in the direction that the orbiting object would travel if orbiting in a counter clockwise circle. Then the vector to the orbiting object is
|
419 |
+
|
420 |
+
We use
|
421 |
+
|
422 |
+
|
423 |
+
|
424 |
+
|
425 |
+
|
426 |
+
|
427 |
+
r
|
428 |
+
˙
|
429 |
+
|
430 |
+
|
431 |
+
|
432 |
+
|
433 |
+
|
434 |
+
{\displaystyle {\dot {r}}}
|
435 |
+
|
436 |
+
and
|
437 |
+
|
438 |
+
|
439 |
+
|
440 |
+
|
441 |
+
|
442 |
+
|
443 |
+
θ
|
444 |
+
˙
|
445 |
+
|
446 |
+
|
447 |
+
|
448 |
+
|
449 |
+
|
450 |
+
{\displaystyle {\dot {\theta }}}
|
451 |
+
|
452 |
+
to denote the standard derivatives of how this distance and angle change over time. We take the derivative of a vector to see how it changes over time by subtracting its location at time
|
453 |
+
|
454 |
+
|
455 |
+
|
456 |
+
|
457 |
+
t
|
458 |
+
|
459 |
+
|
460 |
+
{\displaystyle t}
|
461 |
+
|
462 |
+
from that at time
|
463 |
+
|
464 |
+
|
465 |
+
|
466 |
+
t
|
467 |
+
+
|
468 |
+
δ
|
469 |
+
t
|
470 |
+
|
471 |
+
|
472 |
+
{\displaystyle t+\delta t}
|
473 |
+
|
474 |
+
and dividing by
|
475 |
+
|
476 |
+
|
477 |
+
|
478 |
+
δ
|
479 |
+
t
|
480 |
+
|
481 |
+
|
482 |
+
{\displaystyle \delta t}
|
483 |
+
|
484 |
+
. The result is also a vector. Because our basis vector
|
485 |
+
|
486 |
+
|
487 |
+
|
488 |
+
|
489 |
+
|
490 |
+
|
491 |
+
|
492 |
+
r
|
493 |
+
|
494 |
+
^
|
495 |
+
|
496 |
+
|
497 |
+
|
498 |
+
|
499 |
+
|
500 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
501 |
+
|
502 |
+
moves as the object orbits, we start by differentiating it.
|
503 |
+
From time
|
504 |
+
|
505 |
+
|
506 |
+
|
507 |
+
t
|
508 |
+
|
509 |
+
|
510 |
+
{\displaystyle t}
|
511 |
+
|
512 |
+
to
|
513 |
+
|
514 |
+
|
515 |
+
|
516 |
+
t
|
517 |
+
+
|
518 |
+
δ
|
519 |
+
t
|
520 |
+
|
521 |
+
|
522 |
+
{\displaystyle t+\delta t}
|
523 |
+
|
524 |
+
,
|
525 |
+
the vector
|
526 |
+
|
527 |
+
|
528 |
+
|
529 |
+
|
530 |
+
|
531 |
+
|
532 |
+
|
533 |
+
r
|
534 |
+
|
535 |
+
^
|
536 |
+
|
537 |
+
|
538 |
+
|
539 |
+
|
540 |
+
|
541 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
542 |
+
|
543 |
+
keeps its beginning at the origin and rotates from
|
544 |
+
angle
|
545 |
+
|
546 |
+
|
547 |
+
|
548 |
+
θ
|
549 |
+
|
550 |
+
|
551 |
+
{\displaystyle \theta }
|
552 |
+
|
553 |
+
to
|
554 |
+
|
555 |
+
|
556 |
+
|
557 |
+
θ
|
558 |
+
+
|
559 |
+
|
560 |
+
|
561 |
+
|
562 |
+
θ
|
563 |
+
˙
|
564 |
+
|
565 |
+
|
566 |
+
|
567 |
+
|
568 |
+
δ
|
569 |
+
t
|
570 |
+
|
571 |
+
|
572 |
+
{\displaystyle \theta +{\dot {\theta }}\ \delta t}
|
573 |
+
|
574 |
+
which moves its head a distance
|
575 |
+
|
576 |
+
|
577 |
+
|
578 |
+
|
579 |
+
|
580 |
+
|
581 |
+
θ
|
582 |
+
˙
|
583 |
+
|
584 |
+
|
585 |
+
|
586 |
+
|
587 |
+
δ
|
588 |
+
t
|
589 |
+
|
590 |
+
|
591 |
+
{\displaystyle {\dot {\theta }}\ \delta t}
|
592 |
+
|
593 |
+
in the perpendicular direction
|
594 |
+
|
595 |
+
|
596 |
+
|
597 |
+
|
598 |
+
|
599 |
+
|
600 |
+
θ
|
601 |
+
^
|
602 |
+
|
603 |
+
|
604 |
+
|
605 |
+
|
606 |
+
|
607 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}}
|
608 |
+
|
609 |
+
giving a derivative of
|
610 |
+
|
611 |
+
|
612 |
+
|
613 |
+
|
614 |
+
|
615 |
+
|
616 |
+
θ
|
617 |
+
˙
|
618 |
+
|
619 |
+
|
620 |
+
|
621 |
+
|
622 |
+
|
623 |
+
|
624 |
+
θ
|
625 |
+
^
|
626 |
+
|
627 |
+
|
628 |
+
|
629 |
+
|
630 |
+
|
631 |
+
{\displaystyle {\dot {\theta }}{\hat {\boldsymbol {\theta }}}}
|
632 |
+
|
633 |
+
.
|
634 |
+
|
635 |
+
We can now find the velocity and acceleration of our orbiting object.
|
636 |
+
|
637 |
+
The coefficients of
|
638 |
+
|
639 |
+
|
640 |
+
|
641 |
+
|
642 |
+
|
643 |
+
|
644 |
+
|
645 |
+
r
|
646 |
+
|
647 |
+
^
|
648 |
+
|
649 |
+
|
650 |
+
|
651 |
+
|
652 |
+
|
653 |
+
{\displaystyle {\hat {\mathbf {r} }}}
|
654 |
+
|
655 |
+
|
656 |
+
and
|
657 |
+
|
658 |
+
|
659 |
+
|
660 |
+
|
661 |
+
|
662 |
+
|
663 |
+
θ
|
664 |
+
^
|
665 |
+
|
666 |
+
|
667 |
+
|
668 |
+
|
669 |
+
|
670 |
+
{\displaystyle {\hat {\boldsymbol {\theta }}}}
|
671 |
+
|
672 |
+
give the accelerations in the radial and transverse directions.
|
673 |
+
As said, Newton gives this first due to gravity is
|
674 |
+
|
675 |
+
|
676 |
+
|
677 |
+
−
|
678 |
+
μ
|
679 |
+
|
680 |
+
/
|
681 |
+
|
682 |
+
|
683 |
+
r
|
684 |
+
|
685 |
+
2
|
686 |
+
|
687 |
+
|
688 |
+
|
689 |
+
|
690 |
+
{\displaystyle -\mu /r^{2}}
|
691 |
+
|
692 |
+
and the second is zero.
|
693 |
+
|
694 |
+
|
695 |
+
|
696 |
+
|
697 |
+
|
698 |
+
|
699 |
+
|
700 |
+
|
701 |
+
|
702 |
+
(1)
|
703 |
+
|
704 |
+
|
705 |
+
|
706 |
+
|
707 |
+
|
708 |
+
|
709 |
+
|
710 |
+
|
711 |
+
|
712 |
+
(2)
|
713 |
+
|
714 |
+
Equation (2) can be rearranged using integration by parts.
|
715 |
+
|
716 |
+
We can multiply through by
|
717 |
+
|
718 |
+
|
719 |
+
|
720 |
+
r
|
721 |
+
|
722 |
+
|
723 |
+
{\displaystyle r}
|
724 |
+
|
725 |
+
because it is not zero unless the orbiting object crashes.
|
726 |
+
Then having the derivative be zero gives that the function is a constant.
|
727 |
+
|
728 |
+
|
729 |
+
|
730 |
+
|
731 |
+
|
732 |
+
|
733 |
+
|
734 |
+
|
735 |
+
|
736 |
+
(3)
|
737 |
+
|
738 |
+
which is actually the theoretical proof of Kepler's second law (A line joining a planet and the Sun sweeps out equal areas during equal intervals of time). The constant of integration, h, is the angular momentum per unit mass.
|
739 |
+
|
740 |
+
In order to get an equation for the orbit from equation (1), we need to eliminate time.[8] (See also Binet equation.)
|
741 |
+
In polar coordinates, this would express the distance
|
742 |
+
|
743 |
+
|
744 |
+
|
745 |
+
r
|
746 |
+
|
747 |
+
|
748 |
+
{\displaystyle r}
|
749 |
+
|
750 |
+
of the orbiting object from the center as a function of its angle
|
751 |
+
|
752 |
+
|
753 |
+
|
754 |
+
θ
|
755 |
+
|
756 |
+
|
757 |
+
{\displaystyle \theta }
|
758 |
+
|
759 |
+
. However, it is easier to
|
760 |
+
introduce the auxiliary variable
|
761 |
+
|
762 |
+
|
763 |
+
|
764 |
+
u
|
765 |
+
=
|
766 |
+
1
|
767 |
+
|
768 |
+
/
|
769 |
+
|
770 |
+
r
|
771 |
+
|
772 |
+
|
773 |
+
{\displaystyle u=1/r}
|
774 |
+
|
775 |
+
and to express
|
776 |
+
|
777 |
+
|
778 |
+
|
779 |
+
u
|
780 |
+
|
781 |
+
|
782 |
+
{\displaystyle u}
|
783 |
+
|
784 |
+
as a function of
|
785 |
+
|
786 |
+
|
787 |
+
|
788 |
+
θ
|
789 |
+
|
790 |
+
|
791 |
+
{\displaystyle \theta }
|
792 |
+
|
793 |
+
. Derivatives of
|
794 |
+
|
795 |
+
|
796 |
+
|
797 |
+
r
|
798 |
+
|
799 |
+
|
800 |
+
{\displaystyle r}
|
801 |
+
|
802 |
+
with respect to time may be rewritten as derivatives of
|
803 |
+
|
804 |
+
|
805 |
+
|
806 |
+
u
|
807 |
+
|
808 |
+
|
809 |
+
{\displaystyle u}
|
810 |
+
|
811 |
+
with respect to angle.
|
812 |
+
|
813 |
+
Plugging these into (1) gives
|
814 |
+
|
815 |
+
So for the gravitational force – or, more generally, for any inverse square force law – the right hand side of the equation becomes a constant and the equation is seen to be the harmonic equation (up to a shift of origin of the dependent variable). The solution is:
|
816 |
+
|
817 |
+
where A and θ0 are arbitrary constants.
|
818 |
+
This resulting equation of the orbit of the object is that of an ellipse in Polar form relative to one of the focal points. This is put into a more standard form by
|
819 |
+
letting
|
820 |
+
|
821 |
+
|
822 |
+
|
823 |
+
e
|
824 |
+
≡
|
825 |
+
|
826 |
+
h
|
827 |
+
|
828 |
+
2
|
829 |
+
|
830 |
+
|
831 |
+
A
|
832 |
+
|
833 |
+
/
|
834 |
+
|
835 |
+
μ
|
836 |
+
|
837 |
+
|
838 |
+
{\displaystyle e\equiv h^{2}A/\mu }
|
839 |
+
|
840 |
+
be the eccentricity,
|
841 |
+
letting
|
842 |
+
|
843 |
+
|
844 |
+
|
845 |
+
a
|
846 |
+
≡
|
847 |
+
|
848 |
+
h
|
849 |
+
|
850 |
+
2
|
851 |
+
|
852 |
+
|
853 |
+
|
854 |
+
/
|
855 |
+
|
856 |
+
(
|
857 |
+
μ
|
858 |
+
(
|
859 |
+
1
|
860 |
+
−
|
861 |
+
|
862 |
+
e
|
863 |
+
|
864 |
+
2
|
865 |
+
|
866 |
+
|
867 |
+
)
|
868 |
+
)
|
869 |
+
|
870 |
+
|
871 |
+
{\displaystyle a\equiv h^{2}/(\mu (1-e^{2}))}
|
872 |
+
|
873 |
+
be the semi-major axis.
|
874 |
+
Finally, letting
|
875 |
+
|
876 |
+
|
877 |
+
|
878 |
+
|
879 |
+
θ
|
880 |
+
|
881 |
+
0
|
882 |
+
|
883 |
+
|
884 |
+
≡
|
885 |
+
0
|
886 |
+
|
887 |
+
|
888 |
+
{\displaystyle \theta _{0}\equiv 0}
|
889 |
+
|
890 |
+
so the long axis of the ellipse is along the positive x coordinate.
|
891 |
+
|
892 |
+
The above classical (Newtonian) analysis of orbital mechanics assumes that the more subtle effects of general relativity, such as frame dragging and gravitational time dilation are negligible. Relativistic effects cease to be negligible when near very massive bodies (as with the precession of Mercury's orbit about the Sun), or when extreme precision is needed (as with calculations of the orbital elements and time signal references for GPS satellites.[9]).
|
893 |
+
|
894 |
+
The analysis so far has been two dimensional; it turns out that an unperturbed orbit is two-dimensional in a plane fixed in space, and thus the extension to three dimensions requires simply rotating the two-dimensional plane into the required angle relative to the poles of the planetary body involved.
|
895 |
+
|
896 |
+
The rotation to do this in three dimensions requires three numbers to uniquely determine; traditionally these are expressed as three angles.
|
897 |
+
|
898 |
+
The orbital period is simply how long an orbiting body takes to complete one orbit.
|
899 |
+
|
900 |
+
Six parameters are required to specify a Keplerian orbit about a body. For example, the three numbers that specify the body's initial position, and the three values that specify its velocity will define a unique orbit that can be calculated forwards (or backwards) in time. However, traditionally the parameters used are slightly different.
|
901 |
+
|
902 |
+
The traditionally used set of orbital elements is called the set of Keplerian elements, after Johannes Kepler and his laws. The Keplerian elements are six:
|
903 |
+
|
904 |
+
In principle once the orbital elements are known for a body, its position can be calculated forward and backwards indefinitely in time. However, in practice, orbits are affected or perturbed, by other forces than simple gravity from an assumed point source (see the next section), and thus the orbital elements change over time.
|
905 |
+
|
906 |
+
An orbital perturbation is when a force or impulse which is much smaller than the overall force or average impulse of the main gravitating body and which is external to the two orbiting bodies causes an acceleration, which changes the parameters of the orbit over time.
|
907 |
+
|
908 |
+
A small radial impulse given to a body in orbit changes the eccentricity, but not the orbital period (to first order). A prograde or retrograde impulse (i.e. an impulse applied along the orbital motion) changes both the eccentricity and the orbital period. Notably, a prograde impulse at periapsis raises the altitude at apoapsis, and vice versa, and a retrograde impulse does the opposite. A transverse impulse (out of the orbital plane) causes rotation of the orbital plane without changing the period or eccentricity. In all instances, a closed orbit will still intersect the perturbation point.
|
909 |
+
|
910 |
+
If an orbit is about a planetary body with significant atmosphere, its orbit can decay because of drag. Particularly at each periapsis, the object experiences atmospheric drag, losing energy. Each time, the orbit grows less eccentric (more circular) because the object loses kinetic energy precisely when that energy is at its maximum. This is similar to the effect of slowing a pendulum at its lowest point; the highest point of the pendulum's swing becomes lower. With each successive slowing more of the orbit's path is affected by the atmosphere and the effect becomes more pronounced. Eventually, the effect becomes so great that the maximum kinetic energy is not enough to return the orbit above the limits of the atmospheric drag effect. When this happens the body will rapidly spiral down and intersect the central body.
|
911 |
+
|
912 |
+
The bounds of an atmosphere vary wildly. During a solar maximum, the Earth's atmosphere causes drag up to a hundred kilometres higher than during a solar minimum.
|
913 |
+
|
914 |
+
Some satellites with long conductive tethers can also experience orbital decay because of electromagnetic drag from the Earth's magnetic field. As the wire cuts the magnetic field it acts as a generator, moving electrons from one end to the other. The orbital energy is converted to heat in the wire.
|
915 |
+
|
916 |
+
Orbits can be artificially influenced through the use of rocket engines which change the kinetic energy of the body at some point in its path. This is the conversion of chemical or electrical energy to kinetic energy. In this way changes in the orbit shape or orientation can be facilitated.
|
917 |
+
|
918 |
+
Another method of artificially influencing an orbit is through the use of solar sails or magnetic sails. These forms of propulsion require no propellant or energy input other than that of the Sun, and so can be used indefinitely. See statite for one such proposed use.
|
919 |
+
|
920 |
+
Orbital decay can occur due to tidal forces for objects below the synchronous orbit for the body they're orbiting. The gravity of the orbiting object raises tidal bulges in the primary, and since below the synchronous orbit the orbiting object is moving faster than the body's surface the bulges lag a short angle behind it. The gravity of the bulges is slightly off of the primary-satellite axis and thus has a component along the satellite's motion. The near bulge slows the object more than the far bulge speeds it up, and as a result the orbit decays. Conversely, the gravity of the satellite on the bulges applies torque on the primary and speeds up its rotation. Artificial satellites are too small to have an appreciable tidal effect on the planets they orbit, but several moons in the Solar System are undergoing orbital decay by this mechanism. Mars' innermost moon Phobos is a prime example, and is expected to either impact Mars' surface or break up into a ring within 50 million years.
|
921 |
+
|
922 |
+
Orbits can decay via the emission of gravitational waves. This mechanism is extremely weak for most stellar objects, only becoming significant in cases where there is a combination of extreme mass and extreme acceleration, such as with black holes or neutron stars that are orbiting each other closely.
|
923 |
+
|
924 |
+
The standard analysis of orbiting bodies assumes that all bodies consist of uniform spheres, or more generally, concentric shells each of uniform density. It can be shown that such bodies are gravitationally equivalent to point sources.
|
925 |
+
|
926 |
+
However, in the real world, many bodies rotate, and this introduces oblateness and distorts the gravity field, and gives a quadrupole moment to the gravitational field which is significant at distances comparable to the radius of the body. In the general case, the gravitational potential of a rotating body such as, e.g., a planet is usually expanded in multipoles accounting for the departures of it from spherical symmetry. From the point of view of satellite dynamics, of particular relevance are the so-called even zonal harmonic coefficients, or even zonals, since they induce secular orbital perturbations which are cumulative over time spans longer than the orbital period.[10][11][12] They do depend on the orientation of the body's symmetry axis in the space, affecting, in general, the whole orbit, with the exception of the semimajor axis.
|
927 |
+
|
928 |
+
The effects of other gravitating bodies can be significant. For example, the orbit of the Moon cannot be accurately described without allowing for the action of the Sun's gravity as well as the Earth's. One approximate result is that bodies will usually have reasonably stable orbits around a heavier planet or moon, in spite of these perturbations, provided they are orbiting well within the heavier body's Hill sphere.
|
929 |
+
|
930 |
+
When there are more than two gravitating bodies it is referred to as an n-body problem. Most n-body problems have no closed form solution, although some special cases have been formulated.
|
931 |
+
|
932 |
+
For smaller bodies particularly, light and stellar wind can cause significant perturbations to the attitude and direction of motion of the body, and over time can be significant. Of the planetary bodies, the motion of asteroids is particularly affected over large periods when the asteroids are rotating relative to the Sun.
|
933 |
+
|
934 |
+
Mathematicians have discovered that it is possible in principle to have multiple bodies in non-elliptical orbits that repeat periodically, although most such orbits are not stable regarding small perturbations in mass, position, or velocity. However, some special stable cases have been identified, including a planar figure-eight orbit occupied by three moving bodies. Further studies have discovered that nonplanar orbits are also possible, including one involving 12 masses moving in 4 roughly circular, interlocking orbits topologically equivalent to the edges of a cuboctahedron.[13]
|
935 |
+
|
936 |
+
Finding such orbits naturally occurring in the universe is thought to be extremely unlikely, because of the improbability of the required conditions occurring by chance.[13]
|
937 |
+
|
938 |
+
Orbital mechanics or astrodynamics is the application of ballistics and celestial mechanics to the practical problems concerning the motion of rockets and other spacecraft. The motion of these objects is usually calculated from Newton's laws of motion and Newton's law of universal gravitation. It is a core discipline within space mission design and control. Celestial mechanics treats more broadly the orbital dynamics of systems under the influence of gravity, including spacecraft and natural astronomical bodies such as star systems, planets, moons, and comets. Orbital mechanics focuses on spacecraft trajectories, including orbital maneuvers, orbit plane changes, and interplanetary transfers, and is used by mission planners to predict the results of propulsive maneuvers. General relativity is a more exact theory than Newton's laws for calculating orbits, and is sometimes necessary for greater accuracy or in high-gravity situations (such as orbits close to the Sun).
|
939 |
+
|
940 |
+
The gravitational constant G has been calculated as:
|
941 |
+
|
942 |
+
Thus the constant has dimension density−1 time−2. This corresponds to the following properties.
|
943 |
+
|
944 |
+
Scaling of distances (including sizes of bodies, while keeping the densities the same) gives similar orbits without scaling the time: if for example distances are halved, masses are divided by 8, gravitational forces by 16 and gravitational accelerations by 2. Hence velocities are halved and orbital periods remain the same. Similarly, when an object is dropped from a tower, the time it takes to fall to the ground remains the same with a scale model of the tower on a scale model of the Earth.
|
945 |
+
|
946 |
+
Scaling of distances while keeping the masses the same (in the case of point masses, or by reducing the densities) gives similar orbits; if distances are multiplied by 4, gravitational forces and accelerations are divided by 16, velocities are halved and orbital periods are multiplied by 8.
|
947 |
+
|
948 |
+
When all densities are multiplied by 4, orbits are the same; gravitational forces are multiplied by 16 and accelerations by 4, velocities are doubled and orbital periods are halved.
|
949 |
+
|
950 |
+
When all densities are multiplied by 4, and all sizes are halved, orbits are similar; masses are divided by 2, gravitational forces are the same, gravitational accelerations are doubled. Hence velocities are the same and orbital periods are halved.
|
951 |
+
|
952 |
+
In all these cases of scaling. if densities are multiplied by 4, times are halved; if velocities are doubled, forces are multiplied by 16.
|
953 |
+
|
954 |
+
These properties are illustrated in the formula (derived from the formula for the orbital period)
|
955 |
+
|
956 |
+
for an elliptical orbit with semi-major axis a, of a small body around a spherical body with radius r and average density ρ, where T is the orbital period. See also Kepler's Third Law.
|
957 |
+
|
958 |
+
The application of certain orbits or orbital maneuvers to specific useful purposes have been the subject of patents.[17]
|
959 |
+
|
960 |
+
Some bodies are tidally locked with other bodies, meaning that one side of the celestial body is permanently facing its host object. This is the case for Earth-Moon and Pluto-Charon system.
|
961 |
+
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1 |
+
An orchestra (/ˈɔːrkɪstrə/; Italian: [orˈkɛstra]) is a large instrumental ensemble typical of classical music, which combines instruments from different families, including bowed string instruments such as the violin, viola, cello, and double bass, brass instruments such as the horn, trumpet, trombone and tuba, woodwinds such as the flute, oboe, clarinet and bassoon, and percussion instruments such as the timpani, bass drum, triangle, snare drum, cymbals, and mallet percussion instruments each grouped in sections. Other instruments such as the piano and celesta may sometimes appear in a fifth keyboard section or may stand alone, as may the concert harp and, for performances of some modern compositions, electronic instruments.
|
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|
3 |
+
A full-size Western orchestra may sometimes be called a symphony orchestra or philharmonic orchestra (from Greek phil-, "loving", and "harmonic"). The actual number of musicians employed in a given performance may vary from seventy to over one hundred musicians, depending on the work being played and the size of the venue. A chamber orchestra (sometimes concert orchestra) is a smaller-sized ensemble of about fifty musicians or fewer.[citation needed] Orchestras that specialize in the Baroque music of, for example, Johann Sebastian Bach and George Frideric Handel, or Classical repertoire, such as that of Haydn and Mozart, tend to be smaller than orchestras performing a Romantic music repertoire,[citation needed] such as the symphonies of Johannes Brahms. The typical orchestra grew in size throughout the 18th and 19th centuries, reaching a peak with the large orchestras (of as many as 120 players) called for in the works of Richard Wagner, and later, Gustav Mahler.
|
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|
5 |
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Orchestras are usually led by a conductor who directs the performance with movements of the hands and arms, often made easier for the musicians to see by use of a conductor's baton. The conductor unifies the orchestra, sets the tempo and shapes the sound of the ensemble.[1] The conductor also prepares the orchestra by leading rehearsals before the public concert, in which the conductor provides instructions to the musicians on their interpretation of the music being performed.
|
6 |
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|
7 |
+
The leader of the first violin section, commonly called the concertmaster, also plays an important role in leading the musicians. In the Baroque music era (1600–1750), orchestras were often led by the concertmaster or by a chord-playing musician performing the basso continuo parts on a harpsichord or pipe organ, a tradition that some 20th-century and 21st-century early music ensembles continue. Orchestras play a wide range of repertoire, including symphonies, opera and ballet overtures, concertos for solo instruments, and as pit ensembles for operas, ballets, and some types of musical theatre (e.g., Gilbert and Sullivan operettas).
|
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|
9 |
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Amateur orchestras include those made up of students from an elementary school or a high school, youth orchestras, and community orchestras; the latter two typically being made up of amateur musicians from a particular city or region.
|
10 |
+
|
11 |
+
The term orchestra derives from the Greek ὀρχήστρα (orchestra), the name for the area in front of a stage in ancient Greek theatre reserved for the Greek chorus.[2]
|
12 |
+
|
13 |
+
The invention of the piston and rotary valve by Heinrich Stölzel and Friedrich Blühmel, both Silesians, in 1815, was the first in a series of innovations which impacted the orchestra, including the development of modern keywork for the flute by Theobald Boehm and the innovations of Adolphe Sax in the woodwinds, notably the invention of the saxophone. These advances would lead Hector Berlioz to write a landmark book on instrumentation, which was the first systematic treatise on the use of instrumental sound as an expressive element of music.[3]
|
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|
15 |
+
The next major expansion of symphonic practice came from Richard Wagner's Bayreuth orchestra, founded to accompany his musical dramas. Wagner's works for the stage were scored with unprecedented scope and complexity: indeed, his score to Das Rheingold calls for six harps. Thus, Wagner envisioned an ever-more-demanding role for the conductor of the theatre orchestra, as he elaborated in his influential work On Conducting.[4] This brought about a revolution in orchestral composition, and set the style for orchestral performance for the next eighty years. Wagner's theories re-examined the importance of tempo, dynamics, bowing of string instruments and the role of principals in the orchestra.
|
16 |
+
|
17 |
+
At the beginning of the 20th century, symphony orchestras were larger, better funded, and better trained than previously; consequently, composers could compose larger and more ambitious works. The works of Gustav Mahler were particularly innovative; in his later symphonies, such as the mammoth Symphony No. 8, Mahler pushes the furthest boundaries of orchestral size, employing huge forces. By the late Romantic era, orchestras could support the most enormous forms of symphonic expression, with huge string sections, massive brass sections and an expanded range of percussion instruments. With the recording era beginning, the standards of performance were pushed to a new level, because a recorded symphony could be listened to closely and even minor errors in intonation or ensemble, which might not be noticeable in a live performance, could be heard by critics. As recording technologies improved over the 20th and 21st centuries, eventually small errors in a recording could be "fixed" by audio editing or overdubbing. Some older conductors and composers could remember a time when simply "getting through" the music as best as possible was the standard. Combined with the wider audience made possible by recording, this led to a renewed focus on particular star conductors and on a high standard of orchestral execution.[5]
|
18 |
+
|
19 |
+
The typical symphony orchestra consists of four groups of related musical instruments called the woodwinds, brass, percussion, and strings. Other instruments such as the piano and celesta may sometimes be grouped into a fifth section such as a keyboard section or may stand alone, as may the concert harp and electric and electronic instruments. The orchestra, depending on the size, contains almost all of the standard instruments in each group.
|
20 |
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|
21 |
+
In the history of the orchestra, its instrumentation has been expanded over time, often agreed to have been standardized by the classical period[6] and Ludwig van Beethoven's influence on the classical model.[7] In the 20th and 21st century, new repertory demands expanded the instrumentation of the orchestra, resulting in a flexible use of the classical-model instruments and newly developed electric and electronic instruments in various combinations.
|
22 |
+
|
23 |
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The terms symphony orchestra and philharmonic orchestra may be used to distinguish different ensembles from the same locality, such as the London Symphony Orchestra and the London Philharmonic Orchestra. A symphony orchestra will usually have over eighty musicians on its roster, in some cases over a hundred, but the actual number of musicians employed in a particular performance may vary according to the work being played and the size of the venue.
|
24 |
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|
25 |
+
Chamber orchestra usually refers to smaller-sized ensembles; a major chamber orchestra might employ as many as fifty musicians; some are much smaller than that. The term concert orchestra may also be used, as in the BBC Concert Orchestra and the RTÉ Concert Orchestra.
|
26 |
+
|
27 |
+
The so-called "standard complement" of doubled winds and brass in the orchestra from the first half of the 19th century is generally attributed to the forces called for by Beethoven.[citation needed] The composer's instrumentation almost always included paired flutes, oboes, clarinets, bassoons, horns and trumpets. The exceptions to this are his Symphony No. 4, Violin Concerto, and Piano Concerto No. 4, which each specify a single flute. Beethoven carefully calculated the expansion of this particular timbral "palette" in Symphonies 3, 5, 6, and 9 for an innovative effect. The third horn in the "Eroica" Symphony arrives to provide not only some harmonic flexibility, but also the effect of "choral" brass in the Trio movement. Piccolo, contrabassoon, and trombones add to the triumphal finale of his Symphony No. 5. A piccolo and a pair of trombones help deliver the effect of storm and sunshine in the Sixth, also known as the Pastoral Symphony. The Ninth asks for a second pair of horns, for reasons similar to the "Eroica" (four horns has since become standard); Beethoven's use of piccolo, contrabassoon, trombones, and untuned percussion—plus chorus and vocal soloists—in his finale, are his earliest suggestion that the timbral boundaries of symphony might be expanded. For several decades after his death, symphonic instrumentation was faithful to Beethoven's well-established model, with few exceptions.[citation needed]
|
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|
29 |
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Apart from the core orchestral complement, various other instruments are called for occasionally.[8] These include the flugelhorn and cornet. Saxophones and classical guitars, for example, appear in some 19th- through 21st-century scores. While appearing only as featured solo instruments in some works, for example Maurice Ravel's orchestration of Modest Mussorgsky's Pictures at an Exhibition and Sergei Rachmaninoff's Symphonic Dances, the saxophone is included in other works, such as Ravel's Boléro, Sergei Prokofiev's Romeo and Juliet Suites 1 and 2, Vaughan Williams' Symphonies No.6 and 9 and William Walton's Belshazzar's Feast, and many other works as a member of the orchestral ensemble. The euphonium is featured in a few late Romantic and 20th-century works, usually playing parts marked "tenor tuba", including Gustav Holst's The Planets, and Richard Strauss's Ein Heldenleben. The Wagner tuba, a modified member of the horn family, appears in Richard Wagner's cycle Der Ring des Nibelungen and several other works by Strauss, Béla Bartók, and others; it has a prominent role in Anton Bruckner's Symphony No. 7 in E Major.[9] Cornets appear in Pyotr Ilyich Tchaikovsky's ballet Swan Lake, Claude Debussy's La Mer, and several orchestral works by Hector Berlioz. Unless these instruments are played by members "doubling" on another instrument (for example, a trombone player changing to euphonium or a bassoon player switching to contrabassoon for a certain passage), orchestras typically hire freelance musicians to augment their regular ensemble.
|
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|
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The 20th-century orchestra was far more flexible than its predecessors.[10] In Beethoven's and Felix Mendelssohn's time, the orchestra was composed of a fairly standard core of instruments, which was very rarely modified by composers. As time progressed, and as the Romantic period saw changes in accepted modification with composers such as Berlioz and Mahler; some composers used multiple harps and sound effect such as the wind machine. During the 20th century, the modern orchestra was generally standardized with the modern instrumentation listed below. Nevertheless, by the mid- to late 20th century, with the development of contemporary classical music, instrumentation could practically be hand-picked by the composer (e.g., to add electric instruments such as electric guitar, electronic instruments such as synthesizers, non-Western instruments, or other instruments not traditionally used in orchestra).
|
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|
33 |
+
With this history in mind, the orchestra can be analysed in five eras: the Baroque era, the Classical era, early/mid-Romantic music era, late-Romantic era and combined Modern/Postmodern eras. The first is a Baroque orchestra (i.e., J.S. Bach, Handel, Vivaldi), which generally had a smaller number of performers, and in which one or more chord-playing instruments, the basso continuo group (e.g., harpsichord or pipe organ and assorted bass instruments to perform the bassline), played an important role; the second is a typical classical period orchestra (e.g., early Beethoven along with Mozart and Haydn), which used a smaller group of performers than a Romantic music orchestra and a fairly standardized instrumentation; the third is typical of an early/mid-Romantic era (e.g., Schubert, Berlioz, Schumann, Brahms); the fourth is a late-Romantic/early 20th-century orchestra (e.g., Wagner, Mahler, Stravinsky), to the common complement of a 2010-era modern orchestra (e.g., Adams, Barber, Aaron Copland, Glass, Penderecki).
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Among the instrument groups and within each group of instruments, there is a generally accepted hierarchy. Every instrumental group (or section) has a principal who is generally responsible for leading the group and playing orchestral solos. The violins are divided into two groups, first violin and second violin, with the second violins playing in lower registers than the first violins, playing an accompaniment part, or harmonizing the melody played by the first violins. The principal first violin is called the concertmaster (or "leader" in the UK) and is not only considered the leader of the string section, but the second-in-command of the entire orchestra, behind only the conductor. The concertmaster leads the pre-concert tuning and handles musical aspects of orchestra management, such as determining the bowings for the violins or for all of the string section. The concertmaster usually sits to the conductor's left, closest to the audience. There is also a principal second violin, a principal viola, a principal cello and a principal bass.
|
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|
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+
The principal trombone is considered the leader of the low brass section, while the principal trumpet is generally considered the leader of the entire brass section. While the oboe often provides the tuning note for the orchestra (due to 300-year-old convention), there is generally no designated principal of the woodwind section (though in woodwind ensembles, the flute is often considered the leader.)[11] Instead, each principal confers with the others as equals in the case of musical differences of opinion. Most sections also have an assistant principal (or co-principal or associate principal), or in the case of the first violins, an assistant concertmaster, who often plays a tutti part in addition to replacing the principal in his or her absence.
|
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A section string player plays in unison with the rest of the section, except in the case of divided (divisi) parts, where upper and lower parts in the music are often assigned to "outside" (nearer the audience) and "inside" seated players. Where a solo part is called for in a string section, the section leader invariably plays that part. The section leader (or principal) of a string section is also responsible for determining the bowings, often based on the bowings set out by the concertmaster. In some cases, the principal of a string section may use a slightly different bowing than the concertmaster, to accommodate the requirements of playing their instrument (e.g., the double-bass section). Principals of a string section will also lead entrances for their section, typically by lifting the bow before the entrance, to ensure the section plays together. Tutti wind and brass players generally play a unique but non-solo part. Section percussionists play parts assigned to them by the principal percussionist.
|
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|
65 |
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In modern times, the musicians are usually directed by a conductor, although early orchestras did not have one, giving this role instead to the concertmaster or the harpsichordist playing the continuo. Some modern orchestras also do without conductors, particularly smaller orchestras and those specializing in historically accurate (so-called "period") performances of baroque and earlier music.
|
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|
67 |
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The most frequently performed repertoire for a symphony orchestra is Western classical music or opera. However, orchestras are used sometimes in popular music (e.g., to accompany a rock or pop band in a concert), extensively in film music, and increasingly often in video game music. Orchestras are also used in the symphonic metal genre. The term "orchestra" can also be applied to a jazz ensemble, for example in the performance of big-band music.
|
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|
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In the 2000s, all tenured members of a professional orchestra normally audition for positions in the ensemble. Performers typically play one or more solo pieces of the auditionee's choice, such as a movement of a concerto, a solo Bach movement, and a variety of excerpts from the orchestral literature that are advertised in the audition poster (so the auditionees can prepare). The excerpts are typically the most technically challenging parts and solos from the orchestral literature. Orchestral auditions are typically held in front of a panel that includes the conductor, the concertmaster, the principal player of the section for which the auditionee is applying, and possibly other principal players.
|
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|
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The most promising candidates from the first round of auditions are invited to return for a second or third round of auditions, which allows the conductor and the panel to compare the best candidates. Performers may be asked to sight read orchestral music. The final stage of the audition process in some orchestras is a test week, in which the performer plays with the orchestra for a week or two, which allows the conductor and principal players to see if the individual can function well in an actual rehearsal and performance setting.
|
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|
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There are a range of different employment arrangements. The most sought-after positions are permanent, tenured positions in the orchestra. Orchestras also hire musicians on contracts, ranging in length from a single concert to a full season or more. Contract performers may be hired for individual concerts when the orchestra is doing an exceptionally large late-Romantic era orchestral work, or to substitute for a permanent member who is sick. A professional musician who is hired to perform for a single concert is sometimes called a "sub". Some contract musicians may be hired to replace permanent members for the period that the permanent member is on parental leave or disability leave.
|
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|
75 |
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Historically, major professional orchestras have been mostly or entirely composed of male musicians. The first female members hired in professional orchestras have been harpists. The Vienna Philharmonic, for example, did not accept women to permanent membership until 1997, far later than comparable orchestras (the other orchestras ranked among the world’s top five by Gramophone in 2008).[12] The last major orchestra to appoint a woman to a permanent position was the Berlin Philharmonic.[13] In February 1996, the Vienna Philharmonic's principal flute, Dieter Flury, told Westdeutscher Rundfunk that accepting women would be "gambling with the emotional unity (emotionelle Geschlossenheit) that this organism currently has".[14] In April 1996, the orchestra’s press secretary wrote that "compensating for the expected leaves of absence" of maternity leave would be a problem.[15]
|
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|
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In 1997, the Vienna Philharmonic was "facing protests during a [US] tour" by the National Organization for Women and the International Alliance for Women in Music. Finally, "after being held up to increasing ridicule even in socially conservative Austria, members of the orchestra gathered [on 28 February 1997] in an extraordinary meeting on the eve of their departure and agreed to admit a woman, Anna Lelkes, as harpist."[16] As of 2013, the orchestra has six female members; one of them, violinist Albena Danailova, became one of the orchestra’s concertmasters in 2008, the first woman to hold that position.[17] In 2012, women made up 6% of the orchestra's membership. VPO president Clemens Hellsberg said the VPO now uses completely screened blind auditions.[18]
|
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|
79 |
+
In 2013, an article in Mother Jones stated that while "[m]any prestigious orchestras have significant female membership—women outnumber men in the New York Philharmonic's violin section—and several renowned ensembles, including the National Symphony Orchestra, the Detroit Symphony, and the Minnesota Symphony, are led by women violinists", the double bass, brass, and percussion sections of major orchestras "...are still predominantly male."[19] A 2014 BBC article stated that the "...introduction of ‘blind’ auditions, where a prospective instrumentalist performs behind a screen so that the judging panel can exercise no gender or racial prejudice, has seen the gender balance of traditionally male-dominated symphony orchestras gradually shift."[20]
|
80 |
+
|
81 |
+
There are also a variety of amateur orchestras:
|
82 |
+
|
83 |
+
Orchestras play a wide range of repertoire ranging from 17th-century dance suites, 18th-century divertimentos to 20th-century film scores and 21st-century symphonies. Orchestras have become synonymous with the symphony, an extended musical composition in Western classical music that typically contains multiple movements which provide contrasting keys and tempos. Symphonies are notated in a musical score, which contains all the instrument parts. The conductor uses the score to study the symphony before rehearsals and decide on their interpretation (e.g., tempos, articulation, phrasing, etc.), and to follow the music during rehearsals and concerts, while leading the ensemble. Orchestral musicians play from parts containing just the notated music for their instrument. A small number of symphonies also contain vocal parts (e.g., Beethoven's Ninth Symphony).
|
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+
|
85 |
+
Orchestras also perform overtures, a term originally applied to the instrumental introduction to an opera.[21] During the early Romantic era, composers such as Beethoven and Mendelssohn began to use the term to refer to independent, self-existing instrumental, programmatic works that presaged genres such as the symphonic poem, a form devised by Franz Liszt in several works that began as dramatic overtures. These were "at first undoubtedly intended to be played at the head of a programme".[21] In the 1850s the concert overture began to be supplanted by the symphonic poem.
|
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+
|
87 |
+
Orchestras also play with instrumental soloists in concertos. During concertos, the orchestra plays an accompaniment role to the soloist (e.g., a solo violinist or pianist) and, at times, introduces musical themes or interludes while the soloist is not playing. Orchestras also play during operas, ballets, some musical theatre works and some choral works (both sacred works such as Masses and secular works). In operas and ballets, the orchestra accompanies the singers and dancers, respectively, and plays overtures and interludes where the melodies played by the orchestra take centre stage.
|
88 |
+
|
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+
In the Baroque era, orchestras performed in a range of venues, including at the fine houses of aristocrats, in opera halls and in churches. Some wealthy aristocrats had an orchestra in residence at their estate, to entertain them and their guests with performances. During the Classical era, as composers increasingly sought out financial support from the general public, orchestra concerts were increasingly held in public concert halls, where music lovers could buy tickets to hear the orchestra. Aristocratic patronage of orchestras continued during the Classical era, but this went on alongside public concerts. In the 20th and 21st century, orchestras found a new patron: governments. Many orchestras in North America and Europe receive part of their funding from national, regional level governments (e.g., state governments in the U.S.) or city governments. These government subsidies make up part of orchestra revenue, along with ticket sales, charitable donations (if the orchestra is registered as a charity) and other fundraising activities. With the invention of successive technologies, including sound recording, radio broadcasting, television broadcasting and Internet-based streaming and downloading of concert videos, orchestras have been able to find new revenue sources.
|
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+
|
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+
One of the "great unmentionable [topics] of orchestral playing" is "faking", the process by which an orchestral musician gives the "...impression of playing every note as written", typically for a very challenging passage that is very high or very fast, while not actually playing the notes that are in the printed music part.[22] An article in The Strad states that all orchestral musicians, even those in the top orchestras, occasionally "fake" certain passages.[22] One reason that musicians "fake" is because there are not enough rehearsals.[22] Another factor is the extreme challenges in 20th-century and 21st-century contemporary pieces; some professionals said "faking" was "necessary in anything from ten to almost ninety per cent of some modern works.[22] Professional players who were interviewed were of a consensus that faking may be acceptable when a part is not written well for the instrument, but faking "just because you haven’t practised" the music is not acceptable.[22]
|
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+
With the advent of the early music movement, smaller orchestras where players worked on execution of works in styles derived from the study of older treatises on playing became common. These include the Orchestra of the Age of Enlightenment, the London Classical Players under the direction of Sir Roger Norrington and the Academy of Ancient Music under Christopher Hogwood, among others.[citation needed]
|
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+
|
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+
In the United States, the late 20th century saw a crisis of funding and support for orchestras. The size and cost of a symphony orchestra, compared to the size of the base of supporters, became an issue that struck at the core of the institution. Few orchestras could fill auditoriums, and the time-honored season-subscription system became increasingly anachronistic, as more and more listeners would buy tickets on an ad hoc basis for individual events. Orchestral endowments and—more centrally to the daily operation of American orchestras—orchestral donors have seen investment portfolios shrink or produce lower yields, reducing the ability of donors to contribute; further, there has been a trend toward donors finding other social causes more compelling. While government funding is less central to American than European orchestras, cuts in such funding are still significant for American ensembles. Finally, the drastic falling-off of revenues from recording, tied to no small extent to changes in the recording industry itself, began a period of change that has yet to reach its conclusion.
|
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+
|
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+
U.S. orchestras that have gone into Chapter 11 bankruptcy include the Philadelphia Orchestra (in April 2011), and the Louisville Orchestra, in December 2010; orchestras that have gone into Chapter 7 bankruptcy and have ceased operations include the Northwest Chamber Orchestra in 2006, the Honolulu Orchestra in March 2011, the New Mexico Symphony Orchestra in April 2011, and the Syracuse Symphony in June 2011. The Festival of Orchestras in Orlando, Florida ceased operations at the end of March, 2011.
|
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One source of financial difficulties that received notice and criticism was high salaries for music directors of US orchestras,[23] which led several high-profile conductors to take pay cuts in recent years.[24][25][26] Music administrators such as Michael Tilson Thomas and Esa-Pekka Salonen argued that new music, new means of presenting it, and a renewed relationship with the community could revitalize the symphony orchestra. The American critic Greg Sandow has argued in detail that orchestras must revise their approach to music, performance, the concert experience, marketing, public relations, community involvement, and presentation to bring them in line with the expectations of 21st-century audiences immersed in popular culture.
|
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It is not uncommon for contemporary composers to use unconventional instruments, including various synthesizers, to achieve desired effects. Many, however, find more conventional orchestral configuration to provide better possibilities for color and depth. Composers like John Adams often employ Romantic-size orchestras, as in Adams' opera Nixon in China; Philip Glass and others may be more free, yet still identify size-boundaries. Glass in particular has recently turned to conventional orchestras in works like the Concerto for Cello and Orchestra and the Violin Concerto No. 2.
|
102 |
+
|
103 |
+
Along with a decrease in funding, some U.S. orchestras have reduced their overall personnel, as well as the number of players appearing in performances. The reduced numbers in performance are usually confined to the string section, since the numbers here have traditionally been flexible (as multiple players typically play from the same part).
|
104 |
+
|
105 |
+
Conducting is the art of directing a musical performance, such as an orchestral or choral concert. The primary duties of the conductor are to set the tempo, ensure correct entries by various members of the ensemble, and "shape" the phrasing where appropriate.[1] To convey their ideas and interpretation, a conductor communicates with their musicians primarily through hand gestures, typically though not invariably with the aid of a baton, and may use other gestures or signals, such as eye contact with relevant performers.[27] A conductor's directions will almost invariably be supplemented or reinforced by verbal instructions or suggestions to their musicians in rehearsal prior to a performance.[27]
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+
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+
The conductor typically stands on a raised podium with a large music stand for the full score, which contains the musical notation for all the instruments and voices. Since the mid-18th century, most conductors have not played an instrument when conducting,[citation needed] although in earlier periods of classical music history, leading an ensemble while playing an instrument was common. In Baroque music from the 1600s to the 1750s, the group would typically be led by the harpsichordist or first violinist (see concertmaster), an approach that in modern times has been revived by several music directors for music from this period. Conducting while playing a piano or synthesizer may also be done with musical theatre pit orchestras. Communication is typically non-verbal during a performance (this is strictly the case in art music, but in jazz big bands or large pop ensembles, there may be occasional spoken instructions, such as a "count in"). However, in rehearsals, frequent interruptions allow the conductor to give verbal directions as to how the music should be played or sung.
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+
Conductors act as guides to the orchestras or choirs they conduct. They choose the works to be performed and study their scores, to which they may make certain adjustments (e.g., regarding tempo, articulation, phrasing, repetitions of sections, and so on), work out their interpretation, and relay their vision to the performers. They may also attend to organizational matters, such as scheduling rehearsals,[28] planning a concert season, hearing auditions and selecting members, and promoting their ensemble in the media. Orchestras, choirs, concert bands and other sizable musical ensembles such as big bands are usually led by conductors.
|
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+
In the Baroque music era (1600–1750), most orchestras were led by one of the musicians, typically the principal first violin, called the concertmaster. The concertmaster would lead the tempo of pieces by lifting his or her bow in a rhythmic manner. Leadership might also be provided by one of the chord-playing instrumentalists playing the basso continuo part which was the core of most Baroque instrumental ensemble pieces. Typically, this would be a harpsichord player, a pipe organist or a luteist or theorbo player. A keyboard player could lead the ensemble with his or her head, or by taking one of the hands off the keyboard to lead a more difficult tempo change. A lutenist or theorbo player could lead by lifting the instrument neck up and down to indicate the tempo of a piece, or to lead a ritard during a cadence or ending. In some works which combined choirs and instrumental ensembles, two leaders were sometimes used: a concertmaster to lead the instrumentalists and a chord-playing performer to lead the singers. During the Classical music period (ca. 1720–1800), the practice of using chordal instruments to play basso continuo was gradually phased out, and it disappeared completely by 1800. Instead, ensembles began to use conductors to lead the orchestra's tempos and playing style, while the concertmaster played an additional leadership role for the musicians, especially the string players, who imitate the bowstroke and playing style of the concertmaster, to the degree that is feasible for the different stringed instruments.
|
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In 1922, the idea of a conductor-less orchestra was revived in post-revolutionary Soviet Union. The symphony orchestra Persimfans was formed without a conductor, because the founders believed that the ensemble should be modeled on the ideal Marxist state, in which all people are equal. As such, its members felt that there was no need to be led by the dictatorial baton of a conductor; instead they were led by a committee, which determined tempos and playing styles. Although it was a partial success within the Soviet Union, the principal difficulty with the concept was in changing tempo during performances, because even if the committee had issued a decree about where a tempo change should take place, there was no leader in the ensemble to guide this tempo change. The orchestra survived for ten years before Stalin's cultural politics disbanded it by taking away its funding.[29]
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In Western nations, some ensembles, such as the Orpheus Chamber Orchestra, based in New York City, have had more success with conductorless orchestras, although decisions are likely to be deferred to some sense of leadership within the ensemble (for example, the principal wind and string players, notably the concertmaster). Others have returned to the tradition of a principal player, usually a violinist, being the artistic director and running rehearsal and leading concerts. Examples include the Australian Chamber Orchestra, Amsterdam Sinfonietta & Candida Thompson and the New Century Chamber Orchestra. As well, as part of the early music movement, some 20th- and 21st-century orchestras have revived the Baroque practice of having no conductor on the podium for Baroque pieces, using the concertmaster or a chord-playing basso continuo performer (e.g., harpsichord or organ) to lead the group.
|
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Some orchestral works specify that an offstage trumpet should be used or that other instruments from the orchestra should be positioned off-stage or behind the stage, to create a haunted, mystical effect. To ensure that the offstage instrumentalist(s) play in time, sometimes a sub-conductor will be stationed offstage with a clear view of the principal conductor. Examples include the ending of "Neptune" from Gustav Holst's The Planets. The principal conductor leads the large orchestra, and the sub-conductor relays the principal conductor's tempo and gestures to the offstage musician (or musicians). One of the challenges with using two conductors is that the second conductor may get out of synchronization with the main conductor, or may mis-convey (or misunderstand) the principal conductor's gestures, which can lead to the offstage instruments being out of time. In the late 20th century and early 21st century, some orchestras use a video camera pointed at the principal conductor and a closed-circuit TV set in front of the offstage performer(s), instead of using two conductors.
|
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The techniques of polystylism and polytempo[30] music have led a few 20th- and 21st-century composers to write music where multiple orchestras or ensembles perform simultaneously. These trends have brought about the phenomenon of polyconductor music, wherein separate sub-conductors conduct each group of musicians. Usually, one principal conductor conducts the sub-conductors, thereby shaping the overall performance. In Percy Grainger's "The Warriors" which includes three conductors: the primary conductor of the orchestra, a secondary conductor directing an off-stage brass ensemble, and a tertiary conductor directing percussion and harp. One example in the late-century orchestral music is Karlheinz Stockhausen's Gruppen, for three orchestras, which are placed around the audience. This way, the "sound masses" could be spacialized, as in an electroacoustic work. Gruppen was premiered in Cologne, in 1958, conducted by Stockhausen, Bruno Maderna and Pierre Boulez. It has been performed in 1996 by Simon Rattle, John Carewe and Daniel Harding.[31]
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1 |
+
In set theory, an ordinal number, or ordinal, is one generalization of the concept of a natural number that is used to describe a way to arrange a (possibly infinite) collection of objects in order, one after another. Any finite collection of objects can be put in order just by the process of counting: labeling the objects with distinct natural numbers. Ordinal numbers are thus the "labels" needed to arrange collections of objects in order.
|
2 |
+
|
3 |
+
An ordinal number is used to describe the order type of a well-ordered set (though this does not work for a well-ordered proper class). A well-ordered set is a set with a relation > such that
|
4 |
+
|
5 |
+
Two well-ordered sets have the same order type if and only if there is a bijection from one set to the other that converts the relation in the first set to the relation in the second set.
|
6 |
+
|
7 |
+
Whereas ordinals are useful for ordering the objects in a collection, they are distinct from cardinal numbers, which are useful for saying how many objects are in a collection. Although the distinction between ordinals and cardinals is not always apparent in finite sets (one can go from one to the other just by counting labels), different infinite ordinals can correspond to the same cardinal. Like other kinds of numbers, ordinals can be added, multiplied, and exponentiated, although none of these operations are commutative.
|
8 |
+
|
9 |
+
Ordinals were introduced by Georg Cantor in 1883[1] in order to accommodate infinite sequences and classify derived sets, which he had previously introduced in 1872 while studying the uniqueness of trigonometric series.[2]
|
10 |
+
|
11 |
+
A natural number (which, in this context, includes the number 0) can be used for two purposes: to describe the size of a set, or to describe the position of an element in a sequence. When restricted to finite sets these two concepts coincide, there is only one way to put a finite set into a linear sequence, up to isomorphism. When dealing with infinite sets one has to distinguish between the notion of size, which leads to cardinal numbers, and the notion of position, which is generalized by the ordinal numbers described here. This is because while any set has only one size (its cardinality), there are many nonisomorphic well-orderings of any infinite set, as explained below.
|
12 |
+
|
13 |
+
Whereas the notion of cardinal number is associated with a set with no particular structure on it, the ordinals are intimately linked with the special kind of sets that are called well-ordered (so intimately linked, in fact, that some mathematicians make no distinction between the two concepts). A well-ordered set is a totally ordered set (given any two elements one defines a smaller and a larger one in a coherent way) in which there is no infinite decreasing sequence (however, there may be infinite increasing sequences); equivalently, every non-empty subset of the set has a least element. Ordinals may be used to label the elements of any given well-ordered set (the smallest element being labelled 0, the one after that 1, the next one 2, "and so on") and to measure the "length" of the whole set by the least ordinal that is not a label for an element of the set. This "length" is called the order type of the set.
|
14 |
+
|
15 |
+
Any ordinal is defined by the set of ordinals that precede it: in fact, the most common definition of ordinals identifies each ordinal as the set of ordinals that precede it. For example, the ordinal 42 is the order type of the ordinals less than it, i.e., the ordinals from 0 (the smallest of all ordinals) to 41 (the immediate predecessor of 42), and it is generally identified as the set {0,1,2,…,41}. Conversely, any set S of ordinals that is downward-closed — meaning that for any ordinal α in S and any ordinal β < α, β is also in S — is (or can be identified with) an ordinal.
|
16 |
+
|
17 |
+
There are infinite ordinals as well: the smallest infinite ordinal is ω, which is the order type of the natural numbers (finite ordinals) and that can even be identified with the set of natural numbers (indeed, the set of natural numbers is well-ordered—as is any set of ordinals—and since it is downward closed it can be identified with the ordinal associated with it, which is exactly how ω is defined).
|
18 |
+
|
19 |
+
Perhaps a clearer intuition of ordinals can be formed by examining a first few of them: as mentioned above, they start with the natural numbers, 0, 1, 2, 3, 4, 5, … After all natural numbers comes the first infinite ordinal, ω, and after that come ω+1, ω+2, ω+3, and so on. (Exactly what addition means will be defined later on: just consider them as names.) After all of these come ω·2 (which is ω+ω), ω·2+1, ω·2+2, and so on, then ω·3, and then later on ω·4. Now the set of ordinals formed in this way (the ω·m+n, where m and n are natural numbers) must itself have an ordinal associated with it: and that is ω2. Further on, there will be ω3, then ω4, and so on, and ωω, then ωωω, then later ωωωω, and even later ε0 (epsilon nought) (to give a few examples of relatively small—countable—ordinals). This can be continued indefinitely far ("indefinitely far" is exactly what ordinals are good at: every time one says "and so on" when enumerating ordinals, it defines a larger ordinal). The smallest uncountable ordinal is the set of all countable ordinals, expressed as ω1.
|
20 |
+
|
21 |
+
In a well-ordered set, every non-empty subset contains a distinct smallest element. Given the axiom of dependent choice, this is equivalent to just saying that the set is totally ordered and there is no infinite decreasing sequence, something perhaps easier to visualize. In practice, the importance of well-ordering is justified by the possibility of applying transfinite induction, which says, essentially, that any property that passes on from the predecessors of an element to that element itself must be true of all elements (of the given well-ordered set). If the states of a computation (computer program or game) can be well-ordered in such a way that each step is followed by a "lower" step, then the computation will terminate.
|
22 |
+
|
23 |
+
It is inappropriate to distinguish between two well-ordered sets if they only differ in the "labeling of their elements", or more formally: if the elements of the first set can be paired off with the elements of the second set such that if one element is smaller than another in the first set, then the partner of the first element is smaller than the partner of the second element in the second set, and vice versa. Such a one-to-one correspondence is called an order isomorphism and the two well-ordered sets are said to be order-isomorphic, or similar (obviously this is an equivalence relation). Formally, if a partial order ≤ is defined on the set S, and a partial order ≤' is defined on the set S' , then the posets (S,≤) and (S' ,≤') are order isomorphic if there is a bijection f that preserves the ordering. That is, f(a) ≤' f(b) if and only if a ≤ b. Provided there exists an order isomorphism between two well-ordered sets, the order isomorphism is unique: this makes it quite justifiable to consider the two sets as essentially identical, and to seek a "canonical" representative of the isomorphism type (class). This is exactly what the ordinals provide, and it also provides a canonical labeling of the elements of any well-ordered set. Every well-ordered set (S,<) is order-isomorphic to the set of ordinals less than one specific ordinal number under their natural ordering. This canonical set is the order type of (S,<).
|
24 |
+
|
25 |
+
Essentially, an ordinal is intended to be defined as an isomorphism class of well-ordered sets: that is, as an equivalence class for the equivalence relation of "being order-isomorphic". There is a technical difficulty involved, however, in the fact that the equivalence class is too large to be a set in the usual Zermelo–Fraenkel (ZF) formalization of set theory. But this is not a serious difficulty. The ordinal can be said to be the order type of any set in the class.
|
26 |
+
|
27 |
+
The original definition of ordinal numbers, found for example in the Principia Mathematica, defines the order type of a well-ordering as the set of all well-orderings similar (order-isomorphic) to that well-ordering: in other words, an ordinal number is genuinely an equivalence class of well-ordered sets. This definition must be abandoned in ZF and related systems of axiomatic set theory because these equivalence classes are too large to form a set. However, this definition still can be used in type theory and in Quine's axiomatic set theory New Foundations and related systems (where it affords a rather surprising alternative solution to the Burali-Forti paradox of the largest ordinal).
|
28 |
+
|
29 |
+
Rather than defining an ordinal as an equivalence class of well-ordered sets, it will be defined as a particular well-ordered set that (canonically) represents the class. Thus, an ordinal number will be a well-ordered set; and every well-ordered set will be order-isomorphic to exactly one ordinal number.
|
30 |
+
|
31 |
+
For each well-ordered set
|
32 |
+
|
33 |
+
|
34 |
+
|
35 |
+
T
|
36 |
+
|
37 |
+
|
38 |
+
{\displaystyle T}
|
39 |
+
|
40 |
+
,
|
41 |
+
|
42 |
+
|
43 |
+
|
44 |
+
a
|
45 |
+
↦
|
46 |
+
|
47 |
+
T
|
48 |
+
|
49 |
+
<
|
50 |
+
a
|
51 |
+
|
52 |
+
|
53 |
+
|
54 |
+
|
55 |
+
{\displaystyle a\mapsto T_{<a}}
|
56 |
+
|
57 |
+
defines an order isomorphism between
|
58 |
+
|
59 |
+
|
60 |
+
|
61 |
+
T
|
62 |
+
|
63 |
+
|
64 |
+
{\displaystyle T}
|
65 |
+
|
66 |
+
and the set of all subsets of
|
67 |
+
|
68 |
+
|
69 |
+
|
70 |
+
T
|
71 |
+
|
72 |
+
|
73 |
+
{\displaystyle T}
|
74 |
+
|
75 |
+
having the form
|
76 |
+
|
77 |
+
|
78 |
+
|
79 |
+
|
80 |
+
T
|
81 |
+
|
82 |
+
<
|
83 |
+
a
|
84 |
+
|
85 |
+
|
86 |
+
:=
|
87 |
+
{
|
88 |
+
x
|
89 |
+
∈
|
90 |
+
T
|
91 |
+
∣
|
92 |
+
x
|
93 |
+
<
|
94 |
+
a
|
95 |
+
}
|
96 |
+
|
97 |
+
|
98 |
+
{\displaystyle T_{<a}:=\{x\in T\mid x<a\}}
|
99 |
+
|
100 |
+
ordered by inclusion. This motivates the standard definition, suggested by John von Neumann, now called definition of von Neumann ordinals: "each ordinal is the well-ordered set of all smaller ordinals." In symbols, λ = [0,λ).[3][4] Formally:
|
101 |
+
|
102 |
+
The natural numbers are thus ordinals by this definition. For instance, 2 is an element of 4 = {0, 1, 2, 3}, and 2 is equal to {0, 1} and so it is a subset of {0, 1, 2, 3}.
|
103 |
+
|
104 |
+
It can be shown by transfinite induction that every well-ordered set is order-isomorphic to exactly one of these ordinals, that is, there is an order preserving bijective function between them.
|
105 |
+
|
106 |
+
Furthermore, the elements of every ordinal are ordinals themselves. Given two ordinals S and T, S is an element of T if and only if S is a proper subset of T. Moreover, either S is an element of T, or T is an element of S, or they are equal. So every set of ordinals is totally ordered. Further, every set of ordinals is well-ordered. This generalizes the fact that every set of natural numbers is well-ordered.
|
107 |
+
|
108 |
+
Consequently, every ordinal S is a set having as elements precisely the ordinals smaller than S. For example, every set of ordinals has a supremum, the ordinal obtained by taking the union of all the ordinals in the set. This union exists regardless of the set's size, by the axiom of union.
|
109 |
+
|
110 |
+
The class of all ordinals is not a set. If it were a set, one could show that it was an ordinal and thus a member of itself, which would contradict its strict ordering by membership. This is the Burali-Forti paradox. The class of all ordinals is variously called "Ord", "ON", or "∞".
|
111 |
+
|
112 |
+
An ordinal is finite if and only if the opposite order is also well-ordered, which is the case if and only if each of its non-empty subsets has a maximum.
|
113 |
+
|
114 |
+
There are other modern formulations of the definition of ordinal. For example, assuming the axiom of regularity, the following are equivalent for a set x:
|
115 |
+
|
116 |
+
These definitions cannot be used in non-well-founded set theories. In set theories with urelements, one has to further make sure that the definition excludes urelements from appearing in ordinals.
|
117 |
+
|
118 |
+
If α is any ordinal and X is a set, an α-indexed sequence of elements of X is a function from α to X. This concept, a transfinite sequence (if α is infinite) or ordinal-indexed sequence, is a generalization of the concept of a sequence. An ordinary sequence corresponds to the case α = ω, while a finite α corresponds to a tuple (mathematics), a.k.a. string (computer science).
|
119 |
+
|
120 |
+
Transfinite induction holds in any well-ordered set, but it is so important in relation to ordinals that it is worth restating here.
|
121 |
+
|
122 |
+
That is, if P(α) is true whenever P(β) is true for all β < α, then P(α) is true for all α. Or, more practically: in order to prove a property P for all ordinals α, one can assume that it is already known for all smaller β < α.
|
123 |
+
|
124 |
+
Transfinite induction can be used not only to prove things, but also to define them. Such a definition is normally said to be by transfinite recursion – the proof that the result is well-defined uses transfinite induction. Let F denote a (class) function F to be defined on the ordinals. The idea now is that, in defining F(α) for an unspecified ordinal α, one may assume that F(β) is already defined for all β < α and thus give a formula for F(α) in terms of these F(β). It then follows by transfinite induction that there is one and only one function satisfying the recursion formula up to and including α.
|
125 |
+
|
126 |
+
Here is an example of definition by transfinite recursion on the ordinals (more will be given later): define function F by letting F(α) be the smallest ordinal not in the set {F(β) | β < α}, that is, the set consisting of all F(β) for β < α. This definition assumes the F(β) known in the very process of defining F; this apparent vicious circle is exactly what definition by transfinite recursion permits. In fact, F(0) makes sense since there is no ordinal β < 0, and the set {F(β) | β < 0} is empty. So F(0) is equal to 0 (the smallest ordinal of all). Now that F(0) is known, the definition applied to F(1) makes sense (it is the smallest ordinal not in the singleton set {F(0)} = {0}), and so on (the and so on is exactly transfinite induction). It turns out that this example is not very exciting, since provably F(α) = α for all ordinals α, which can be shown, precisely, by transfinite induction.
|
127 |
+
|
128 |
+
Any nonzero ordinal has the minimum element, zero. It may or may not have a maximum element. For example, 42 has maximum 41 and ω+6 has maximum ω+5. On the other hand, ω does not have a maximum since there is no largest natural number. If an ordinal has a maximum α, then it is the next ordinal after α, and it is called a successor ordinal, namely the successor of α, written α+1. In the von Neumann definition of ordinals, the successor of α is
|
129 |
+
|
130 |
+
|
131 |
+
|
132 |
+
α
|
133 |
+
∪
|
134 |
+
{
|
135 |
+
α
|
136 |
+
}
|
137 |
+
|
138 |
+
|
139 |
+
{\displaystyle \alpha \cup \{\alpha \}}
|
140 |
+
|
141 |
+
since its elements are those of α and α itself.[3]
|
142 |
+
|
143 |
+
A nonzero ordinal that is not a successor is called a limit ordinal. One justification for this term is that a limit ordinal is the limit in a topological sense of all smaller ordinals (under the order topology).
|
144 |
+
|
145 |
+
When
|
146 |
+
|
147 |
+
|
148 |
+
|
149 |
+
⟨
|
150 |
+
|
151 |
+
α
|
152 |
+
|
153 |
+
ι
|
154 |
+
|
155 |
+
|
156 |
+
|
157 |
+
|
|
158 |
+
|
159 |
+
ι
|
160 |
+
<
|
161 |
+
γ
|
162 |
+
⟩
|
163 |
+
|
164 |
+
|
165 |
+
{\displaystyle \langle \alpha _{\iota }|\iota <\gamma \rangle }
|
166 |
+
|
167 |
+
is an ordinal-indexed sequence, indexed by a limit γ and the sequence is increasing, i.e.
|
168 |
+
|
169 |
+
|
170 |
+
|
171 |
+
|
172 |
+
α
|
173 |
+
|
174 |
+
ι
|
175 |
+
|
176 |
+
|
177 |
+
<
|
178 |
+
|
179 |
+
α
|
180 |
+
|
181 |
+
ρ
|
182 |
+
|
183 |
+
|
184 |
+
|
185 |
+
|
186 |
+
|
187 |
+
{\displaystyle \alpha _{\iota }<\alpha _{\rho }\!}
|
188 |
+
|
189 |
+
whenever
|
190 |
+
|
191 |
+
|
192 |
+
|
193 |
+
ι
|
194 |
+
<
|
195 |
+
ρ
|
196 |
+
,
|
197 |
+
|
198 |
+
|
199 |
+
|
200 |
+
{\displaystyle \iota <\rho ,\!}
|
201 |
+
|
202 |
+
its limit is defined as the least upper bound of the set
|
203 |
+
|
204 |
+
|
205 |
+
|
206 |
+
{
|
207 |
+
|
208 |
+
α
|
209 |
+
|
210 |
+
ι
|
211 |
+
|
212 |
+
|
213 |
+
|
214 |
+
|
|
215 |
+
|
216 |
+
ι
|
217 |
+
<
|
218 |
+
γ
|
219 |
+
}
|
220 |
+
,
|
221 |
+
|
222 |
+
|
223 |
+
|
224 |
+
{\displaystyle \{\alpha _{\iota }|\iota <\gamma \},\!}
|
225 |
+
|
226 |
+
that is, the smallest ordinal (it always exists) greater than any term of the sequence. In this sense, a limit ordinal is the limit of all smaller ordinals (indexed by itself). Put more directly, it is the supremum of the set of smaller ordinals.
|
227 |
+
|
228 |
+
Another way of defining a limit ordinal is to say that α is a limit ordinal if and only if:
|
229 |
+
|
230 |
+
So in the following sequence:
|
231 |
+
|
232 |
+
ω is a limit ordinal because for any smaller ordinal (in this example, a natural number) there is another ordinal (natural number) larger than it, but still less than ω.
|
233 |
+
|
234 |
+
Thus, every ordinal is either zero, or a successor (of a well-defined predecessor), or a limit. This distinction is important, because many definitions by transfinite induction rely upon it. Very often, when defining a function F by transfinite induction on all ordinals, one defines F(0), and F(α+1) assuming F(α) is defined, and then, for limit ordinals δ one defines F(δ) as the limit of the F(β) for all β<δ (either in the sense of ordinal limits, as previously explained, or for some other notion of limit if F does not take ordinal values). Thus, the interesting step in the definition is the successor step, not the limit ordinals. Such functions (especially for F nondecreasing and taking ordinal values) are called continuous. Ordinal addition, multiplication and exponentiation are continuous as functions of their second argument.
|
235 |
+
|
236 |
+
Any well-ordered set is similar (order-isomorphic) to a unique ordinal number
|
237 |
+
|
238 |
+
|
239 |
+
|
240 |
+
α
|
241 |
+
|
242 |
+
|
243 |
+
{\displaystyle \alpha }
|
244 |
+
|
245 |
+
; in other words, its elements can be indexed in increasing fashion by the ordinals less than
|
246 |
+
|
247 |
+
|
248 |
+
|
249 |
+
α
|
250 |
+
|
251 |
+
|
252 |
+
{\displaystyle \alpha }
|
253 |
+
|
254 |
+
. This applies, in particular, to any set of ordinals: any set of ordinals is naturally indexed by the ordinals less than some
|
255 |
+
|
256 |
+
|
257 |
+
|
258 |
+
α
|
259 |
+
|
260 |
+
|
261 |
+
{\displaystyle \alpha }
|
262 |
+
|
263 |
+
. The same holds, with a slight modification, for classes of ordinals (a collection of ordinals, possibly too large to form a set, defined by some property): any class of ordinals can be indexed by ordinals (and, when the class is unbounded in the class of all ordinals, this puts it in class-bijection with the class of all ordinals). So the
|
264 |
+
|
265 |
+
|
266 |
+
|
267 |
+
γ
|
268 |
+
|
269 |
+
|
270 |
+
{\displaystyle \gamma }
|
271 |
+
|
272 |
+
-th element in the class (with the convention that the "0-th" is the smallest, the "1-st" is the next smallest, and so on) can be freely spoken of. Formally, the definition is by transfinite induction: the
|
273 |
+
|
274 |
+
|
275 |
+
|
276 |
+
γ
|
277 |
+
|
278 |
+
|
279 |
+
{\displaystyle \gamma }
|
280 |
+
|
281 |
+
-th element of the class is defined (provided it has already been defined for all
|
282 |
+
|
283 |
+
|
284 |
+
|
285 |
+
β
|
286 |
+
<
|
287 |
+
γ
|
288 |
+
|
289 |
+
|
290 |
+
{\displaystyle \beta <\gamma }
|
291 |
+
|
292 |
+
), as the smallest element greater than the
|
293 |
+
|
294 |
+
|
295 |
+
|
296 |
+
β
|
297 |
+
|
298 |
+
|
299 |
+
{\displaystyle \beta }
|
300 |
+
|
301 |
+
-th element for all
|
302 |
+
|
303 |
+
|
304 |
+
|
305 |
+
β
|
306 |
+
<
|
307 |
+
γ
|
308 |
+
|
309 |
+
|
310 |
+
{\displaystyle \beta <\gamma }
|
311 |
+
|
312 |
+
.
|
313 |
+
|
314 |
+
This could be applied, for example, to the class of limit ordinals: the
|
315 |
+
|
316 |
+
|
317 |
+
|
318 |
+
γ
|
319 |
+
|
320 |
+
|
321 |
+
{\displaystyle \gamma }
|
322 |
+
|
323 |
+
-th ordinal, which is either a limit or zero is
|
324 |
+
|
325 |
+
|
326 |
+
|
327 |
+
ω
|
328 |
+
⋅
|
329 |
+
γ
|
330 |
+
|
331 |
+
|
332 |
+
{\displaystyle \omega \cdot \gamma }
|
333 |
+
|
334 |
+
(see ordinal arithmetic for the definition of multiplication of ordinals). Similarly, one can consider additively indecomposable ordinals (meaning a nonzero ordinal that is not the sum of two strictly smaller ordinals): the
|
335 |
+
|
336 |
+
|
337 |
+
|
338 |
+
γ
|
339 |
+
|
340 |
+
|
341 |
+
{\displaystyle \gamma }
|
342 |
+
|
343 |
+
-th additively indecomposable ordinal is indexed as
|
344 |
+
|
345 |
+
|
346 |
+
|
347 |
+
|
348 |
+
ω
|
349 |
+
|
350 |
+
γ
|
351 |
+
|
352 |
+
|
353 |
+
|
354 |
+
|
355 |
+
{\displaystyle \omega ^{\gamma }}
|
356 |
+
|
357 |
+
. The technique of indexing classes of ordinals is often useful in the context of fixed points: for example, the
|
358 |
+
|
359 |
+
|
360 |
+
|
361 |
+
γ
|
362 |
+
|
363 |
+
|
364 |
+
{\displaystyle \gamma }
|
365 |
+
|
366 |
+
-th ordinal
|
367 |
+
|
368 |
+
|
369 |
+
|
370 |
+
α
|
371 |
+
|
372 |
+
|
373 |
+
{\displaystyle \alpha }
|
374 |
+
|
375 |
+
such that
|
376 |
+
|
377 |
+
|
378 |
+
|
379 |
+
|
380 |
+
ω
|
381 |
+
|
382 |
+
α
|
383 |
+
|
384 |
+
|
385 |
+
=
|
386 |
+
α
|
387 |
+
|
388 |
+
|
389 |
+
{\displaystyle \omega ^{\alpha }=\alpha }
|
390 |
+
|
391 |
+
is written
|
392 |
+
|
393 |
+
|
394 |
+
|
395 |
+
|
396 |
+
ε
|
397 |
+
|
398 |
+
γ
|
399 |
+
|
400 |
+
|
401 |
+
|
402 |
+
|
403 |
+
{\displaystyle \varepsilon _{\gamma }}
|
404 |
+
|
405 |
+
. These are called the "epsilon numbers".
|
406 |
+
|
407 |
+
A class
|
408 |
+
|
409 |
+
|
410 |
+
|
411 |
+
C
|
412 |
+
|
413 |
+
|
414 |
+
{\displaystyle C}
|
415 |
+
|
416 |
+
of ordinals is said to be unbounded, or cofinal, when given any ordinal
|
417 |
+
|
418 |
+
|
419 |
+
|
420 |
+
α
|
421 |
+
|
422 |
+
|
423 |
+
{\displaystyle \alpha }
|
424 |
+
|
425 |
+
, there is a
|
426 |
+
|
427 |
+
|
428 |
+
|
429 |
+
β
|
430 |
+
|
431 |
+
|
432 |
+
{\displaystyle \beta }
|
433 |
+
|
434 |
+
in
|
435 |
+
|
436 |
+
|
437 |
+
|
438 |
+
C
|
439 |
+
|
440 |
+
|
441 |
+
{\displaystyle C}
|
442 |
+
|
443 |
+
such that
|
444 |
+
|
445 |
+
|
446 |
+
|
447 |
+
α
|
448 |
+
<
|
449 |
+
β
|
450 |
+
|
451 |
+
|
452 |
+
{\displaystyle \alpha <\beta }
|
453 |
+
|
454 |
+
(then the class must be a proper class, i.e., it cannot be a set). It is said to be closed when the limit of a sequence of ordinals in the class is again in the class: or, equivalently, when the indexing (class-)function
|
455 |
+
|
456 |
+
|
457 |
+
|
458 |
+
F
|
459 |
+
|
460 |
+
|
461 |
+
{\displaystyle F}
|
462 |
+
|
463 |
+
is continuous in the sense that, for
|
464 |
+
|
465 |
+
|
466 |
+
|
467 |
+
δ
|
468 |
+
|
469 |
+
|
470 |
+
{\displaystyle \delta }
|
471 |
+
|
472 |
+
a limit ordinal,
|
473 |
+
|
474 |
+
|
475 |
+
|
476 |
+
F
|
477 |
+
(
|
478 |
+
δ
|
479 |
+
)
|
480 |
+
|
481 |
+
|
482 |
+
{\displaystyle F(\delta )}
|
483 |
+
|
484 |
+
(the
|
485 |
+
|
486 |
+
|
487 |
+
|
488 |
+
δ
|
489 |
+
|
490 |
+
|
491 |
+
{\displaystyle \delta }
|
492 |
+
|
493 |
+
-th ordinal in the class) is the limit of all
|
494 |
+
|
495 |
+
|
496 |
+
|
497 |
+
F
|
498 |
+
(
|
499 |
+
γ
|
500 |
+
)
|
501 |
+
|
502 |
+
|
503 |
+
{\displaystyle F(\gamma )}
|
504 |
+
|
505 |
+
for
|
506 |
+
|
507 |
+
|
508 |
+
|
509 |
+
γ
|
510 |
+
<
|
511 |
+
δ
|
512 |
+
|
513 |
+
|
514 |
+
{\displaystyle \gamma <\delta }
|
515 |
+
|
516 |
+
; this is also the same as being closed, in the topological sense, for the order topology (to avoid talking of topology on proper classes, one can demand that the intersection of the class with any given ordinal is closed for the order topology on that ordinal, this is again equivalent).
|
517 |
+
|
518 |
+
Of particular importance are those classes of ordinals that are closed and unbounded, sometimes called clubs. For example, the class of all limit ordinals is closed and unbounded: this translates the fact that there is always a limit ordinal greater than a given ordinal, and that a limit of limit ordinals is a limit ordinal (a fortunate fact if the terminology is to make any sense at all!). The class of additively indecomposable ordinals, or the class of
|
519 |
+
|
520 |
+
|
521 |
+
|
522 |
+
|
523 |
+
ε
|
524 |
+
|
525 |
+
⋅
|
526 |
+
|
527 |
+
|
528 |
+
|
529 |
+
|
530 |
+
{\displaystyle \varepsilon _{\cdot }}
|
531 |
+
|
532 |
+
ordinals, or the class of cardinals, are all closed unbounded; the set of regular cardinals, however, is unbounded but not closed, and any finite set of ordinals is closed but not unbounded.
|
533 |
+
|
534 |
+
A class is stationary if it has a nonempty intersection with every closed unbounded class. All superclasses of closed unbounded classes are stationary, and stationary classes are unbounded, but there are stationary classes that are not closed and stationary classes that have no closed unbounded subclass (such as the class of all limit ordinals with countable cofinality). Since the intersection of two closed unbounded classes is closed and unbounded, the intersection of a stationary class and a closed unbounded class is stationary. But the intersection of two stationary classes may be empty, e.g. the class of ordinals with cofinality ω with the class of ordinals with uncountable cofinality.
|
535 |
+
|
536 |
+
Rather than formulating these definitions for (proper) classes of ordinals, one can formulate them for sets of ordinals below a given ordinal
|
537 |
+
|
538 |
+
|
539 |
+
|
540 |
+
α
|
541 |
+
|
542 |
+
|
543 |
+
{\displaystyle \alpha }
|
544 |
+
|
545 |
+
: A subset of a limit ordinal
|
546 |
+
|
547 |
+
|
548 |
+
|
549 |
+
α
|
550 |
+
|
551 |
+
|
552 |
+
{\displaystyle \alpha }
|
553 |
+
|
554 |
+
is said to be unbounded (or cofinal) under
|
555 |
+
|
556 |
+
|
557 |
+
|
558 |
+
α
|
559 |
+
|
560 |
+
|
561 |
+
{\displaystyle \alpha }
|
562 |
+
|
563 |
+
provided any ordinal less than
|
564 |
+
|
565 |
+
|
566 |
+
|
567 |
+
α
|
568 |
+
|
569 |
+
|
570 |
+
{\displaystyle \alpha }
|
571 |
+
|
572 |
+
is less than some ordinal in the set. More generally, one can call a subset of any ordinal
|
573 |
+
|
574 |
+
|
575 |
+
|
576 |
+
α
|
577 |
+
|
578 |
+
|
579 |
+
{\displaystyle \alpha }
|
580 |
+
|
581 |
+
cofinal in
|
582 |
+
|
583 |
+
|
584 |
+
|
585 |
+
α
|
586 |
+
|
587 |
+
|
588 |
+
{\displaystyle \alpha }
|
589 |
+
|
590 |
+
provided every ordinal less than
|
591 |
+
|
592 |
+
|
593 |
+
|
594 |
+
α
|
595 |
+
|
596 |
+
|
597 |
+
{\displaystyle \alpha }
|
598 |
+
|
599 |
+
is less than or equal to some ordinal in the set. The subset is said to be closed under
|
600 |
+
|
601 |
+
|
602 |
+
|
603 |
+
α
|
604 |
+
|
605 |
+
|
606 |
+
{\displaystyle \alpha }
|
607 |
+
|
608 |
+
provided it is closed for the order topology in
|
609 |
+
|
610 |
+
|
611 |
+
|
612 |
+
α
|
613 |
+
|
614 |
+
|
615 |
+
{\displaystyle \alpha }
|
616 |
+
|
617 |
+
, i.e. a limit of ordinals in the set is either in the set or equal to
|
618 |
+
|
619 |
+
|
620 |
+
|
621 |
+
α
|
622 |
+
|
623 |
+
|
624 |
+
{\displaystyle \alpha }
|
625 |
+
|
626 |
+
itself.
|
627 |
+
|
628 |
+
There are three usual operations on ordinals: addition, multiplication, and (ordinal) exponentiation. Each can be defined in essentially two different ways: either by constructing an explicit well-ordered set that represents the operation or by using transfinite recursion. The Cantor normal form provides a standardized way of writing ordinals. It uniquely represents each ordinal as a finite sum of ordinal powers of ω. However, this cannot form the basis of a universal ordinal notation due to such self-referential representations as ε0 = ωε0. The so-called "natural" arithmetical operations retain commutativity at the expense of continuity.
|
629 |
+
|
630 |
+
Interpreted as nimbers, ordinals are also subject to nimber arithmetic operations.
|
631 |
+
|
632 |
+
Each ordinal associates with one cardinal, its cardinality. If there is a bijection between two ordinals (e.g. ω = 1 + ω and ω + 1 > ω), then they associate with the same cardinal. Any well-ordered set having an ordinal as its order-type has the same cardinality as that ordinal. The least ordinal associated with a given cardinal is called the initial ordinal of that cardinal. Every finite ordinal (natural number) is initial, and no other ordinal associates with its cardinal. But most infinite ordinals are not initial, as many infinite ordinals associate with the same cardinal. The axiom of choice is equivalent to the statement that every set can be well-ordered, i.e. that every cardinal has an initial ordinal. In theories with the axiom of choice, the cardinal number of any set has an initial ordinal, and one may employ the Von Neumann cardinal assignment as the cardinal's representation. In set theories without the axiom of choice, a cardinal may be represented by the set of sets with that cardinality having minimal rank (see Scott's trick).
|
633 |
+
|
634 |
+
The α-th infinite initial ordinal is written
|
635 |
+
|
636 |
+
|
637 |
+
|
638 |
+
|
639 |
+
ω
|
640 |
+
|
641 |
+
α
|
642 |
+
|
643 |
+
|
644 |
+
|
645 |
+
|
646 |
+
{\displaystyle \omega _{\alpha }}
|
647 |
+
|
648 |
+
, it is always a limit ordinal. Its cardinality is written
|
649 |
+
|
650 |
+
|
651 |
+
|
652 |
+
|
653 |
+
ℵ
|
654 |
+
|
655 |
+
α
|
656 |
+
|
657 |
+
|
658 |
+
|
659 |
+
|
660 |
+
{\displaystyle \aleph _{\alpha }}
|
661 |
+
|
662 |
+
. For example, the cardinality of ω0 = ω is
|
663 |
+
|
664 |
+
|
665 |
+
|
666 |
+
|
667 |
+
ℵ
|
668 |
+
|
669 |
+
0
|
670 |
+
|
671 |
+
|
672 |
+
|
673 |
+
|
674 |
+
{\displaystyle \aleph _{0}}
|
675 |
+
|
676 |
+
, which is also the cardinality of ω2 or ε0 (all are countable ordinals). So ω can be identified with
|
677 |
+
|
678 |
+
|
679 |
+
|
680 |
+
|
681 |
+
ℵ
|
682 |
+
|
683 |
+
0
|
684 |
+
|
685 |
+
|
686 |
+
|
687 |
+
|
688 |
+
{\displaystyle \aleph _{0}}
|
689 |
+
|
690 |
+
, except that the notation
|
691 |
+
|
692 |
+
|
693 |
+
|
694 |
+
|
695 |
+
ℵ
|
696 |
+
|
697 |
+
0
|
698 |
+
|
699 |
+
|
700 |
+
|
701 |
+
|
702 |
+
{\displaystyle \aleph _{0}}
|
703 |
+
|
704 |
+
is used when writing cardinals, and ω when writing ordinals (this is important since, for example,
|
705 |
+
|
706 |
+
|
707 |
+
|
708 |
+
|
709 |
+
ℵ
|
710 |
+
|
711 |
+
0
|
712 |
+
|
713 |
+
|
714 |
+
2
|
715 |
+
|
716 |
+
|
717 |
+
|
718 |
+
|
719 |
+
{\displaystyle \aleph _{0}^{2}}
|
720 |
+
|
721 |
+
=
|
722 |
+
|
723 |
+
|
724 |
+
|
725 |
+
|
726 |
+
ℵ
|
727 |
+
|
728 |
+
0
|
729 |
+
|
730 |
+
|
731 |
+
|
732 |
+
|
733 |
+
{\displaystyle \aleph _{0}}
|
734 |
+
|
735 |
+
whereas
|
736 |
+
|
737 |
+
|
738 |
+
|
739 |
+
|
740 |
+
ω
|
741 |
+
|
742 |
+
2
|
743 |
+
|
744 |
+
|
745 |
+
>
|
746 |
+
ω
|
747 |
+
|
748 |
+
|
749 |
+
{\displaystyle \omega ^{2}>\omega }
|
750 |
+
|
751 |
+
). Also,
|
752 |
+
|
753 |
+
|
754 |
+
|
755 |
+
|
756 |
+
ω
|
757 |
+
|
758 |
+
1
|
759 |
+
|
760 |
+
|
761 |
+
|
762 |
+
|
763 |
+
{\displaystyle \omega _{1}}
|
764 |
+
|
765 |
+
is the smallest uncountable ordinal (to see that it exists, consider the set of equivalence classes of well-orderings of the natural numbers: each such well-ordering defines a countable ordinal, and
|
766 |
+
|
767 |
+
|
768 |
+
|
769 |
+
|
770 |
+
ω
|
771 |
+
|
772 |
+
1
|
773 |
+
|
774 |
+
|
775 |
+
|
776 |
+
|
777 |
+
{\displaystyle \omega _{1}}
|
778 |
+
|
779 |
+
is the order type of that set),
|
780 |
+
|
781 |
+
|
782 |
+
|
783 |
+
|
784 |
+
ω
|
785 |
+
|
786 |
+
2
|
787 |
+
|
788 |
+
|
789 |
+
|
790 |
+
|
791 |
+
{\displaystyle \omega _{2}}
|
792 |
+
|
793 |
+
is the smallest ordinal whose cardinality is greater than
|
794 |
+
|
795 |
+
|
796 |
+
|
797 |
+
|
798 |
+
ℵ
|
799 |
+
|
800 |
+
1
|
801 |
+
|
802 |
+
|
803 |
+
|
804 |
+
|
805 |
+
{\displaystyle \aleph _{1}}
|
806 |
+
|
807 |
+
, and so on, and
|
808 |
+
|
809 |
+
|
810 |
+
|
811 |
+
|
812 |
+
ω
|
813 |
+
|
814 |
+
ω
|
815 |
+
|
816 |
+
|
817 |
+
|
818 |
+
|
819 |
+
{\displaystyle \omega _{\omega }}
|
820 |
+
|
821 |
+
is the limit of the
|
822 |
+
|
823 |
+
|
824 |
+
|
825 |
+
|
826 |
+
ω
|
827 |
+
|
828 |
+
n
|
829 |
+
|
830 |
+
|
831 |
+
|
832 |
+
|
833 |
+
{\displaystyle \omega _{n}}
|
834 |
+
|
835 |
+
for natural numbers n (any limit of cardinals is a cardinal, so this limit is indeed the first cardinal after all the
|
836 |
+
|
837 |
+
|
838 |
+
|
839 |
+
|
840 |
+
ω
|
841 |
+
|
842 |
+
n
|
843 |
+
|
844 |
+
|
845 |
+
|
846 |
+
|
847 |
+
{\displaystyle \omega _{n}}
|
848 |
+
|
849 |
+
).
|
850 |
+
|
851 |
+
The cofinality of an ordinal
|
852 |
+
|
853 |
+
|
854 |
+
|
855 |
+
α
|
856 |
+
|
857 |
+
|
858 |
+
{\displaystyle \alpha }
|
859 |
+
|
860 |
+
is the smallest ordinal
|
861 |
+
|
862 |
+
|
863 |
+
|
864 |
+
δ
|
865 |
+
|
866 |
+
|
867 |
+
{\displaystyle \delta }
|
868 |
+
|
869 |
+
that is the order type of a cofinal subset of
|
870 |
+
|
871 |
+
|
872 |
+
|
873 |
+
α
|
874 |
+
|
875 |
+
|
876 |
+
{\displaystyle \alpha }
|
877 |
+
|
878 |
+
. Notice that a number of authors define cofinality or use it only for limit ordinals. The cofinality of a set of ordinals or any other well-ordered set is the cofinality of the order type of that set.
|
879 |
+
|
880 |
+
Thus for a limit ordinal, there exists a
|
881 |
+
|
882 |
+
|
883 |
+
|
884 |
+
δ
|
885 |
+
|
886 |
+
|
887 |
+
{\displaystyle \delta }
|
888 |
+
|
889 |
+
-indexed strictly increasing sequence with limit
|
890 |
+
|
891 |
+
|
892 |
+
|
893 |
+
α
|
894 |
+
|
895 |
+
|
896 |
+
{\displaystyle \alpha }
|
897 |
+
|
898 |
+
. For example, the cofinality of ω² is ω, because the sequence ω·m (where m ranges over the natural numbers) tends to ω²; but, more generally, any countable limit ordinal has cofinality ω. An uncountable limit ordinal may have either cofinality ω as does
|
899 |
+
|
900 |
+
|
901 |
+
|
902 |
+
|
903 |
+
ω
|
904 |
+
|
905 |
+
ω
|
906 |
+
|
907 |
+
|
908 |
+
|
909 |
+
|
910 |
+
{\displaystyle \omega _{\omega }}
|
911 |
+
|
912 |
+
or an uncountable cofinality.
|
913 |
+
|
914 |
+
The cofinality of 0 is 0. And the cofinality of any successor ordinal is 1. The cofinality of any limit ordinal is at least
|
915 |
+
|
916 |
+
|
917 |
+
|
918 |
+
ω
|
919 |
+
|
920 |
+
|
921 |
+
{\displaystyle \omega }
|
922 |
+
|
923 |
+
.
|
924 |
+
|
925 |
+
An ordinal that is equal to its cofinality is called regular and it is always an initial ordinal. Any limit of regular ordinals is a limit of initial ordinals and thus is also initial even if it is not regular, which it usually is not. If the Axiom of Choice, then
|
926 |
+
|
927 |
+
|
928 |
+
|
929 |
+
|
930 |
+
ω
|
931 |
+
|
932 |
+
α
|
933 |
+
+
|
934 |
+
1
|
935 |
+
|
936 |
+
|
937 |
+
|
938 |
+
|
939 |
+
{\displaystyle \omega _{\alpha +1}}
|
940 |
+
|
941 |
+
is regular for each α. In this case, the ordinals 0, 1,
|
942 |
+
|
943 |
+
|
944 |
+
|
945 |
+
ω
|
946 |
+
|
947 |
+
|
948 |
+
{\displaystyle \omega }
|
949 |
+
|
950 |
+
,
|
951 |
+
|
952 |
+
|
953 |
+
|
954 |
+
|
955 |
+
ω
|
956 |
+
|
957 |
+
1
|
958 |
+
|
959 |
+
|
960 |
+
|
961 |
+
|
962 |
+
{\displaystyle \omega _{1}}
|
963 |
+
|
964 |
+
, and
|
965 |
+
|
966 |
+
|
967 |
+
|
968 |
+
|
969 |
+
ω
|
970 |
+
|
971 |
+
2
|
972 |
+
|
973 |
+
|
974 |
+
|
975 |
+
|
976 |
+
{\displaystyle \omega _{2}}
|
977 |
+
|
978 |
+
are regular, whereas 2, 3,
|
979 |
+
|
980 |
+
|
981 |
+
|
982 |
+
|
983 |
+
ω
|
984 |
+
|
985 |
+
ω
|
986 |
+
|
987 |
+
|
988 |
+
|
989 |
+
|
990 |
+
{\displaystyle \omega _{\omega }}
|
991 |
+
|
992 |
+
, and ωω·2 are initial ordinals that are not regular.
|
993 |
+
|
994 |
+
The cofinality of any ordinal α is a regular ordinal, i.e. the cofinality of the cofinality of α is the same as the cofinality of α. So the cofinality operation is idempotent.
|
995 |
+
|
996 |
+
As mentioned above (see Cantor normal form) the ordinal ε0, which is the smallest satisfying the equation
|
997 |
+
|
998 |
+
|
999 |
+
|
1000 |
+
|
1001 |
+
ω
|
1002 |
+
|
1003 |
+
α
|
1004 |
+
|
1005 |
+
|
1006 |
+
=
|
1007 |
+
α
|
1008 |
+
|
1009 |
+
|
1010 |
+
{\displaystyle \omega ^{\alpha }=\alpha }
|
1011 |
+
|
1012 |
+
, so it is the limit of the sequence 0, 1,
|
1013 |
+
|
1014 |
+
|
1015 |
+
|
1016 |
+
ω
|
1017 |
+
|
1018 |
+
|
1019 |
+
{\displaystyle \omega }
|
1020 |
+
|
1021 |
+
,
|
1022 |
+
|
1023 |
+
|
1024 |
+
|
1025 |
+
|
1026 |
+
ω
|
1027 |
+
|
1028 |
+
ω
|
1029 |
+
|
1030 |
+
|
1031 |
+
|
1032 |
+
|
1033 |
+
{\displaystyle \omega ^{\omega }}
|
1034 |
+
|
1035 |
+
,
|
1036 |
+
|
1037 |
+
|
1038 |
+
|
1039 |
+
|
1040 |
+
ω
|
1041 |
+
|
1042 |
+
|
1043 |
+
ω
|
1044 |
+
|
1045 |
+
ω
|
1046 |
+
|
1047 |
+
|
1048 |
+
|
1049 |
+
|
1050 |
+
|
1051 |
+
|
1052 |
+
{\displaystyle \omega ^{\omega ^{\omega }}}
|
1053 |
+
|
1054 |
+
, etc. Many ordinals can be defined in such a manner as fixed points of certain ordinal functions (the
|
1055 |
+
|
1056 |
+
|
1057 |
+
|
1058 |
+
ι
|
1059 |
+
|
1060 |
+
|
1061 |
+
{\displaystyle \iota }
|
1062 |
+
|
1063 |
+
-th ordinal such that
|
1064 |
+
|
1065 |
+
|
1066 |
+
|
1067 |
+
|
1068 |
+
ω
|
1069 |
+
|
1070 |
+
α
|
1071 |
+
|
1072 |
+
|
1073 |
+
=
|
1074 |
+
α
|
1075 |
+
|
1076 |
+
|
1077 |
+
{\displaystyle \omega ^{\alpha }=\alpha }
|
1078 |
+
|
1079 |
+
is called
|
1080 |
+
|
1081 |
+
|
1082 |
+
|
1083 |
+
|
1084 |
+
ε
|
1085 |
+
|
1086 |
+
ι
|
1087 |
+
|
1088 |
+
|
1089 |
+
|
1090 |
+
|
1091 |
+
{\displaystyle \varepsilon _{\iota }}
|
1092 |
+
|
1093 |
+
, then one could go on trying to find the
|
1094 |
+
|
1095 |
+
|
1096 |
+
|
1097 |
+
ι
|
1098 |
+
|
1099 |
+
|
1100 |
+
{\displaystyle \iota }
|
1101 |
+
|
1102 |
+
-th ordinal such that
|
1103 |
+
|
1104 |
+
|
1105 |
+
|
1106 |
+
|
1107 |
+
ε
|
1108 |
+
|
1109 |
+
α
|
1110 |
+
|
1111 |
+
|
1112 |
+
=
|
1113 |
+
α
|
1114 |
+
|
1115 |
+
|
1116 |
+
{\displaystyle \varepsilon _{\alpha }=\alpha }
|
1117 |
+
|
1118 |
+
, "and so on", but all the subtlety lies in the "and so on"). One could try to do this systematically, but no matter what system is used to define and construct ordinals, there is always an ordinal that lies just above all the ordinals constructed by the system. Perhaps the most important ordinal that limits a system of construction in this manner is the Church–Kleene ordinal,
|
1119 |
+
|
1120 |
+
|
1121 |
+
|
1122 |
+
|
1123 |
+
ω
|
1124 |
+
|
1125 |
+
1
|
1126 |
+
|
1127 |
+
|
1128 |
+
|
1129 |
+
C
|
1130 |
+
K
|
1131 |
+
|
1132 |
+
|
1133 |
+
|
1134 |
+
|
1135 |
+
|
1136 |
+
{\displaystyle \omega _{1}^{\mathrm {CK} }}
|
1137 |
+
|
1138 |
+
(despite the
|
1139 |
+
|
1140 |
+
|
1141 |
+
|
1142 |
+
|
1143 |
+
ω
|
1144 |
+
|
1145 |
+
1
|
1146 |
+
|
1147 |
+
|
1148 |
+
|
1149 |
+
|
1150 |
+
{\displaystyle \omega _{1}}
|
1151 |
+
|
1152 |
+
in the name, this ordinal is countable), which is the smallest ordinal that cannot in any way be represented by a computable function (this can be made rigorous, of course). Considerably large ordinals can be defined below
|
1153 |
+
|
1154 |
+
|
1155 |
+
|
1156 |
+
|
1157 |
+
ω
|
1158 |
+
|
1159 |
+
1
|
1160 |
+
|
1161 |
+
|
1162 |
+
|
1163 |
+
C
|
1164 |
+
K
|
1165 |
+
|
1166 |
+
|
1167 |
+
|
1168 |
+
|
1169 |
+
|
1170 |
+
{\displaystyle \omega _{1}^{\mathrm {CK} }}
|
1171 |
+
|
1172 |
+
, however, which measure the "proof-theoretic strength" of certain formal systems (for example,
|
1173 |
+
|
1174 |
+
|
1175 |
+
|
1176 |
+
|
1177 |
+
ε
|
1178 |
+
|
1179 |
+
0
|
1180 |
+
|
1181 |
+
|
1182 |
+
|
1183 |
+
|
1184 |
+
{\displaystyle \varepsilon _{0}}
|
1185 |
+
|
1186 |
+
measures the strength of Peano arithmetic). Large countable ordinals such as countable admissible ordinals can also be defined above the Church-Kleene ordinal, which are of interest in various parts of logic.[citation needed]
|
1187 |
+
|
1188 |
+
Any ordinal number can be made into a topological space by endowing it with the order topology; this topology is discrete if and only if the ordinal is a countable cardinal, i.e. at most ω. A subset of ω + 1 is open in the order topology if and only if either it is cofinite or it does not contain ω as an element.
|
1189 |
+
|
1190 |
+
See the Topology and ordinals section of the "Order topology" article.
|
1191 |
+
|
1192 |
+
A set is downward closed if anything less than an element of the set is also in the set. If a set of ordinals is downward closed, then that set is an ordinal—the least ordinal not in the set.
|
1193 |
+
|
1194 |
+
Examples:
|
1195 |
+
|
1196 |
+
The transfinite ordinal numbers, which first appeared in 1883,[5] originated in Cantor's work with derived sets. If P is a set of real numbers, the derived set P' is the set of limit points of P. In 1872, Cantor generated the sets P(n) by applying the derived set operation n times to P. In 1880, he pointed out that these sets form the sequence P' ⊇ ··· ⊇ P(n) ⊇ P(n + 1) ⊇ ···, and he continued the derivation process by defining P(∞) as the intersection of these sets. Then he iterated the derived set operation and intersections to extend his sequence of sets into the infinite: P(∞) ⊇ P(∞ + 1) ⊇ P(∞ + 2) ⊇ ··· ⊇ P(2∞) ⊇ ··· ⊇ P(∞2) ⊇ ···.[6] The superscripts containing ∞ are just indices defined by the derivation process.[7]
|
1197 |
+
|
1198 |
+
Cantor used these sets in the theorems: (1) If P(α) = ∅ for some index α, then P' is countable; (2) Conversely, if P' is countable, then there is an index α such that P(α) = ∅. These theorems are proved by partitioning P' into pairwise disjoint sets: P' = (P' ∖ P(2)) ∪ (P(2) ∖ P(3)) ∪ ··· ∪ (P(∞) ∖ P(∞ + 1)) ∪ ··· ∪ P(α). For β < α: since P(β + 1) contains the limit points of P(β), the sets P(β) ∖ P(β + 1) have no limit points. Hence, they are discrete sets, so they are countable. Proof of first theorem: If P(α) = ∅ for some index α, then P' is the countable union of countable sets. Therefore, P' is countable.[8]
|
1199 |
+
|
1200 |
+
The second theorem requires proving the existence of an α such that P(α) = ∅. To prove this, Cantor considered the set of all α having countably many predecessors. To define this set, he defined the transfinite ordinal numbers and transformed the infinite indices into ordinals by replacing ∞ with ω, the first transfinite ordinal number. Cantor called the set of finite ordinals the first number class. The second number class is the set of ordinals whose predecessors form a countably infinite set. The set of all α having countably many predecessors—that is, the set of countable ordinals—is the union of these two number classes. Cantor proved that the cardinality of the second number class is the first uncountable cardinality.[9]
|
1201 |
+
|
1202 |
+
Cantor's second theorem becomes: If P' is countable, then there is a countable ordinal α such that P(α) = ∅. Its proof uses proof by contradiction. Let P' be countable, and assume there is no such α. This assumption produces two cases.
|
1203 |
+
|
1204 |
+
In both cases, P' is uncountable, which contradicts P' being countable. Therefore, there is a countable ordinal α such that P(α) = ∅. Cantor's work with derived sets and ordinal numbers led to the Cantor-Bendixson theorem.[11]
|
1205 |
+
|
1206 |
+
Using successors, limits, and cardinality, Cantor generated an unbounded sequence of ordinal numbers and number classes.[12] The (α + 1)-th number class is the set of ordinals whose predecessors form a set of the same cardinality as the α-th number class. The cardinality of the (α + 1)-th number class is the cardinality immediately following that of the α-th number class.[13] For a limit ordinal α, the α-th number class is the union of the β-th number classes for β < α.[14] Its cardinality is the limit of the cardinalities of these number classes.
|
1207 |
+
|
1208 |
+
If n is finite, the n-th number class has cardinality
|
1209 |
+
|
1210 |
+
|
1211 |
+
|
1212 |
+
|
1213 |
+
ℵ
|
1214 |
+
|
1215 |
+
n
|
1216 |
+
−
|
1217 |
+
1
|
1218 |
+
|
1219 |
+
|
1220 |
+
|
1221 |
+
|
1222 |
+
{\displaystyle \aleph _{n-1}}
|
1223 |
+
|
1224 |
+
. If α ≥ ω, the α-th number class has cardinality
|
1225 |
+
|
1226 |
+
|
1227 |
+
|
1228 |
+
|
1229 |
+
ℵ
|
1230 |
+
|
1231 |
+
α
|
1232 |
+
|
1233 |
+
|
1234 |
+
|
1235 |
+
|
1236 |
+
{\displaystyle \aleph _{\alpha }}
|
1237 |
+
|
1238 |
+
.[15] Therefore, the cardinalities of the number classes correspond one-to-one with the aleph numbers. Also, the α-th number class consists of ordinals different from those in the preceding number classes if and only if α is a non-limit ordinal. Therefore, the non-limit number classes partition the ordinals into pairwise disjoint sets.
|
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1 |
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A computer is a machine that can be instructed to carry out sequences of arithmetic or logical operations automatically via computer programming. Modern computers have the ability to follow generalized sets of operations, called programs. These programs enable computers to perform an extremely wide range of tasks. A "complete" computer including the hardware, the operating system (main software), and peripheral equipment required and used for "full" operation can be referred to as a computer system. This term may as well be used for a group of computers that are connected and work together, in particular a computer network or computer cluster.
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Computers are used as control systems for a wide variety of industrial and consumer devices. This includes simple special purpose devices like microwave ovens and remote controls, factory devices such as industrial robots and computer-aided design, and also general purpose devices like personal computers and mobile devices such as smartphones. The Internet is run on computers and it connects hundreds of millions of other computers and their users.
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Early computers were only conceived as calculating devices. Since ancient times, simple manual devices like the abacus aided people in doing calculations. Early in the Industrial Revolution, some mechanical devices were built to automate long tedious tasks, such as guiding patterns for looms. More sophisticated electrical machines did specialized analog calculations in the early 20th century. The first digital electronic calculating machines were developed during World War II. The first semiconductor transistors in the late 1940s were followed by the silicon-based MOSFET (MOS transistor) and monolithic integrated circuit (IC) chip technologies in the late 1950s, leading to the microprocessor and the microcomputer revolution in the 1970s. The speed, power and versatility of computers have been increasing dramatically ever since then, with MOS transistor counts increasing at a rapid pace (as predicted by Moore's law), leading to the Digital Revolution during the late 20th to early 21st centuries.
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Conventionally, a modern computer consists of at least one processing element, typically a central processing unit (CPU) in the form of a metal-oxide-semiconductor (MOS) microprocessor, along with some type of computer memory, typically MOS semiconductor memory chips. The processing element carries out arithmetic and logical operations, and a sequencing and control unit can change the order of operations in response to stored information. Peripheral devices include input devices (keyboards, mice, joystick, etc.), output devices (monitor screens, printers, etc.), and input/output devices that perform both functions (e.g., the 2000s-era touchscreen). Peripheral devices allow information to be retrieved from an external source and they enable the result of operations to be saved and retrieved.
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According to the Oxford English Dictionary, the first known use of the word "computer" was in 1613 in a book called The Yong Mans Gleanings by English writer Richard Braithwait: "I haue [sic] read the truest computer of Times, and the best Arithmetician that euer [sic] breathed, and he reduceth thy dayes into a short number." This usage of the term referred to a human computer, a person who carried out calculations or computations. The word continued with the same meaning until the middle of the 20th century. During the latter part of this period women were often hired as computers because they could be paid less than their male counterparts.[1] By 1943, most human computers were women.[2]
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The Online Etymology Dictionary gives the first attested use of "computer" in the 1640s, meaning "one who calculates"; this is an "agent noun from compute (v.)". The Online Etymology Dictionary states that the use of the term to mean "'calculating machine' (of any type) is from 1897." The Online Etymology Dictionary indicates that the "modern use" of the term, to mean "programmable digital electronic computer" dates from "1945 under this name; [in a] theoretical [sense] from 1937, as Turing machine".[3]
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Devices have been used to aid computation for thousands of years, mostly using one-to-one correspondence with fingers. The earliest counting device was probably a form of tally stick. Later record keeping aids throughout the Fertile Crescent included calculi (clay spheres, cones, etc.) which represented counts of items, probably livestock or grains, sealed in hollow unbaked clay containers.[4][5] The use of counting rods is one example.
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The abacus was initially used for arithmetic tasks. The Roman abacus was developed from devices used in Babylonia as early as 2400 BC. Since then, many other forms of reckoning boards or tables have been invented. In a medieval European counting house, a checkered cloth would be placed on a table, and markers moved around on it according to certain rules, as an aid to calculating sums of money.[6]
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18 |
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The Antikythera mechanism is believed to be the earliest mechanical analog "computer", according to Derek J. de Solla Price.[7] It was designed to calculate astronomical positions. It was discovered in 1901 in the Antikythera wreck off the Greek island of Antikythera, between Kythera and Crete, and has been dated to c. 100 BC. Devices of a level of complexity comparable to that of the Antikythera mechanism would not reappear until a thousand years later.
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Many mechanical aids to calculation and measurement were constructed for astronomical and navigation use. The planisphere was a star chart invented by Abū Rayhān al-Bīrūnī in the early 11th century.[8] The astrolabe was invented in the Hellenistic world in either the 1st or 2nd centuries BC and is often attributed to Hipparchus. A combination of the planisphere and dioptra, the astrolabe was effectively an analog computer capable of working out several different kinds of problems in spherical astronomy. An astrolabe incorporating a mechanical calendar computer[9][10] and gear-wheels was invented by Abi Bakr of Isfahan, Persia in 1235.[11] Abū Rayhān al-Bīrūnī invented the first mechanical geared lunisolar calendar astrolabe,[12] an early fixed-wired knowledge processing machine[13] with a gear train and gear-wheels,[14] c. 1000 AD.
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The sector, a calculating instrument used for solving problems in proportion, trigonometry, multiplication and division, and for various functions, such as squares and cube roots, was developed in the late 16th century and found application in gunnery, surveying and navigation.
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The planimeter was a manual instrument to calculate the area of a closed figure by tracing over it with a mechanical linkage.
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The slide rule was invented around 1620–1630, shortly after the publication of the concept of the logarithm. It is a hand-operated analog computer for doing multiplication and division. As slide rule development progressed, added scales provided reciprocals, squares and square roots, cubes and cube roots, as well as transcendental functions such as logarithms and exponentials, circular and hyperbolic trigonometry and other functions. Slide rules with special scales are still used for quick performance of routine calculations, such as the E6B circular slide rule used for time and distance calculations on light aircraft.
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In the 1770s, Pierre Jaquet-Droz, a Swiss watchmaker, built a mechanical doll (automaton) that could write holding a quill pen. By switching the number and order of its internal wheels different letters, and hence different messages, could be produced. In effect, it could be mechanically "programmed" to read instructions. Along with two other complex machines, the doll is at the Musée d'Art et d'Histoire of Neuchâtel, Switzerland, and still operates.[15]
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The tide-predicting machine invented by Sir William Thomson in 1872 was of great utility to navigation in shallow waters. It used a system of pulleys and wires to automatically calculate predicted tide levels for a set period at a particular location.
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The differential analyser, a mechanical analog computer designed to solve differential equations by integration, used wheel-and-disc mechanisms to perform the integration. In 1876, Lord Kelvin had already discussed the possible construction of such calculators, but he had been stymied by the limited output torque of the ball-and-disk integrators.[16] In a differential analyzer, the output of one integrator drove the input of the next integrator, or a graphing output. The torque amplifier was the advance that allowed these machines to work. Starting in the 1920s, Vannevar Bush and others developed mechanical differential analyzers.
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Charles Babbage, an English mechanical engineer and polymath, originated the concept of a programmable computer. Considered the "father of the computer",[17] he conceptualized and invented the first mechanical computer in the early 19th century. After working on his revolutionary difference engine, designed to aid in navigational calculations, in 1833 he realized that a much more general design, an Analytical Engine, was possible. The input of programs and data was to be provided to the machine via punched cards, a method being used at the time to direct mechanical looms such as the Jacquard loom. For output, the machine would have a printer, a curve plotter and a bell. The machine would also be able to punch numbers onto cards to be read in later. The Engine incorporated an arithmetic logic unit, control flow in the form of conditional branching and loops, and integrated memory, making it the first design for a general-purpose computer that could be described in modern terms as Turing-complete.[18][19]
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The machine was about a century ahead of its time. All the parts for his machine had to be made by hand – this was a major problem for a device with thousands of parts. Eventually, the project was dissolved with the decision of the British Government to cease funding. Babbage's failure to complete the analytical engine can be chiefly attributed to political and financial difficulties as well as his desire to develop an increasingly sophisticated computer and to move ahead faster than anyone else could follow. Nevertheless, his son, Henry Babbage, completed a simplified version of the analytical engine's computing unit (the mill) in 1888. He gave a successful demonstration of its use in computing tables in 1906.
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During the first half of the 20th century, many scientific computing needs were met by increasingly sophisticated analog computers, which used a direct mechanical or electrical model of the problem as a basis for computation. However, these were not programmable and generally lacked the versatility and accuracy of modern digital computers.[20] The first modern analog computer was a tide-predicting machine, invented by Sir William Thomson in 1872. The differential analyser, a mechanical analog computer designed to solve differential equations by integration using wheel-and-disc mechanisms, was conceptualized in 1876 by James Thomson, the brother of the more famous Lord Kelvin.[16]
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The art of mechanical analog computing reached its zenith with the differential analyzer, built by H. L. Hazen and Vannevar Bush at MIT starting in 1927. This built on the mechanical integrators of James Thomson and the torque amplifiers invented by H. W. Nieman. A dozen of these devices were built before their obsolescence became obvious. By the 1950s, the success of digital electronic computers had spelled the end for most analog computing machines, but analog computers remained in use during the 1950s in some specialized applications such as education (control systems) and aircraft (slide rule).
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By 1938, the United States Navy had developed an electromechanical analog computer small enough to use aboard a submarine. This was the Torpedo Data Computer, which used trigonometry to solve the problem of firing a torpedo at a moving target. During World War II similar devices were developed in other countries as well.
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Early digital computers were electromechanical; electric switches drove mechanical relays to perform the calculation. These devices had a low operating speed and were eventually superseded by much faster all-electric computers, originally using vacuum tubes. The Z2, created by German engineer Konrad Zuse in 1939, was one of the earliest examples of an electromechanical relay computer.[21]
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In 1941, Zuse followed his earlier machine up with the Z3, the world's first working electromechanical programmable, fully automatic digital computer.[22][23] The Z3 was built with 2000 relays, implementing a 22 bit word length that operated at a clock frequency of about 5–10 Hz.[24] Program code was supplied on punched film while data could be stored in 64 words of memory or supplied from the keyboard. It was quite similar to modern machines in some respects, pioneering numerous advances such as floating point numbers. Rather than the harder-to-implement decimal system (used in Charles Babbage's earlier design), using a binary system meant that Zuse's machines were easier to build and potentially more reliable, given the technologies available at that time.[25] The Z3 was not itself a universal computer but could be extended to be Turing complete.[26][27]
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Purely electronic circuit elements soon replaced their mechanical and electromechanical equivalents, at the same time that digital calculation replaced analog. The engineer Tommy Flowers, working at the Post Office Research Station in London in the 1930s, began to explore the possible use of electronics for the telephone exchange. Experimental equipment that he built in 1934 went into operation five years later, converting a portion of the telephone exchange network into an electronic data processing system, using thousands of vacuum tubes.[20] In the US, John Vincent Atanasoff and Clifford E. Berry of Iowa State University developed and tested the Atanasoff–Berry Computer (ABC) in 1942,[28] the first "automatic electronic digital computer".[29] This design was also all-electronic and used about 300 vacuum tubes, with capacitors fixed in a mechanically rotating drum for memory.[30]
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During World War II, the British at Bletchley Park achieved a number of successes at breaking encrypted German military communications. The German encryption machine, Enigma, was first attacked with the help of the electro-mechanical bombes which were often run by women.[31][32] To crack the more sophisticated German Lorenz SZ 40/42 machine, used for high-level Army communications, Max Newman and his colleagues commissioned Flowers to build the Colossus.[30] He spent eleven months from early February 1943 designing and building the first Colossus.[33] After a functional test in December 1943, Colossus was shipped to Bletchley Park, where it was delivered on 18 January 1944[34] and attacked its first message on 5 February.[30]
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Colossus was the world's first electronic digital programmable computer.[20] It used a large number of valves (vacuum tubes). It had paper-tape input and was capable of being configured to perform a variety of boolean logical operations on its data, but it was not Turing-complete. Nine Mk II Colossi were built (The Mk I was converted to a Mk II making ten machines in total). Colossus Mark I contained 1,500 thermionic valves (tubes), but Mark II with 2,400 valves, was both 5 times faster and simpler to operate than Mark I, greatly speeding the decoding process.[35][36]
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The ENIAC[37] (Electronic Numerical Integrator and Computer) was the first electronic programmable computer built in the U.S. Although the ENIAC was similar to the Colossus, it was much faster, more flexible, and it was Turing-complete. Like the Colossus, a "program" on the ENIAC was defined by the states of its patch cables and switches, a far cry from the stored program electronic machines that came later. Once a program was written, it had to be mechanically set into the machine with manual resetting of plugs and switches. The programmers of the ENIAC were six women, often known collectively as the "ENIAC girls".[38][39]
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It combined the high speed of electronics with the ability to be programmed for many complex problems. It could add or subtract 5000 times a second, a thousand times faster than any other machine. It also had modules to multiply, divide, and square root. High speed memory was limited to 20 words (about 80 bytes). Built under the direction of John Mauchly and J. Presper Eckert at the University of Pennsylvania, ENIAC's development and construction lasted from 1943 to full operation at the end of 1945. The machine was huge, weighing 30 tons, using 200 kilowatts of electric power and contained over 18,000 vacuum tubes, 1,500 relays, and hundreds of thousands of resistors, capacitors, and inductors.[40]
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The principle of the modern computer was proposed by Alan Turing in his seminal 1936 paper,[41] On Computable Numbers. Turing proposed a simple device that he called "Universal Computing machine" and that is now known as a universal Turing machine. He proved that such a machine is capable of computing anything that is computable by executing instructions (program) stored on tape, allowing the machine to be programmable. The fundamental concept of Turing's design is the stored program, where all the instructions for computing are stored in memory. Von Neumann acknowledged that the central concept of the modern computer was due to this paper.[42] Turing machines are to this day a central object of study in theory of computation. Except for the limitations imposed by their finite memory stores, modern computers are said to be Turing-complete, which is to say, they have algorithm execution capability equivalent to a universal Turing machine.
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Early computing machines had fixed programs. Changing its function required the re-wiring and re-structuring of the machine.[30] With the proposal of the stored-program computer this changed. A stored-program computer includes by design an instruction set and can store in memory a set of instructions (a program) that details the computation. The theoretical basis for the stored-program computer was laid by Alan Turing in his 1936 paper. In 1945, Turing joined the National Physical Laboratory and began work on developing an electronic stored-program digital computer. His 1945 report "Proposed Electronic Calculator" was the first specification for such a device. John von Neumann at the University of Pennsylvania also circulated his First Draft of a Report on the EDVAC in 1945.[20]
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The Manchester Baby was the world's first stored-program computer. It was built at the Victoria University of Manchester by Frederic C. Williams, Tom Kilburn and Geoff Tootill, and ran its first program on 21 June 1948.[43] It was designed as a testbed for the Williams tube, the first random-access digital storage device.[44] Although the computer was considered "small and primitive" by the standards of its time, it was the first working machine to contain all of the elements essential to a modern electronic computer.[45] As soon as the Baby had demonstrated the feasibility of its design, a project was initiated at the university to develop it into a more usable computer, the Manchester Mark 1. Grace Hopper was the first person to develop a compiler for programming language.[2]
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The Mark 1 in turn quickly became the prototype for the Ferranti Mark 1, the world's first commercially available general-purpose computer.[46] Built by Ferranti, it was delivered to the University of Manchester in February 1951. At least seven of these later machines were delivered between 1953 and 1957, one of them to Shell labs in Amsterdam.[47] In October 1947, the directors of British catering company J. Lyons & Company decided to take an active role in promoting the commercial development of computers. The LEO I computer became operational in April 1951[48] and ran the world's first regular routine office computer job.
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The concept of a field-effect transistor was proposed by Julius Edgar Lilienfeld in 1925. John Bardeen and Walter Brattain, while working under William Shockley at Bell Labs, built the first working transistor, the point-contact transistor, in 1947, which was followed by Shockley's bipolar junction transistor in 1948.[49][50] From 1955 onwards, transistors replaced vacuum tubes in computer designs, giving rise to the "second generation" of computers. Compared to vacuum tubes, transistors have many advantages: they are smaller, and require less power than vacuum tubes, so give off less heat. Junction transistors were much more reliable than vacuum tubes and had longer, indefinite, service life. Transistorized computers could contain tens of thousands of binary logic circuits in a relatively compact space. However, early junction transistors were relatively bulky devices that were difficult to manufacture on a mass-production basis, which limited them to a number of specialised applications.[51]
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At the University of Manchester, a team under the leadership of Tom Kilburn designed and built a machine using the newly developed transistors instead of valves.[52] Their first transistorised computer and the first in the world, was operational by 1953, and a second version was completed there in April 1955. However, the machine did make use of valves to generate its 125 kHz clock waveforms and in the circuitry to read and write on its magnetic drum memory, so it was not the first completely transistorized computer. That distinction goes to the Harwell CADET of 1955,[53] built by the electronics division of the Atomic Energy Research Establishment at Harwell.[53][54]
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The metal–oxide–silicon field-effect transistor (MOSFET), also known as the MOS transistor, was invented by Mohamed M. Atalla and Dawon Kahng at Bell Labs in 1959.[55] It was the first truly compact transistor that could be miniaturised and mass-produced for a wide range of uses.[51] With its high scalability,[56] and much lower power consumption and higher density than bipolar junction transistors,[57] the MOSFET made it possible to build high-density integrated circuits.[58][59] In addition to data processing, it also enabled the practical use of MOS transistors as memory cell storage elements, leading to the development of MOS semiconductor memory, which replaced earlier magnetic-core memory in computers.[60] The MOSFET led to the microcomputer revolution,[61] and became the driving force behind the computer revolution.[62][63] The MOSFET is the most widely used transistor in computers,[64][65] and is the fundamental building block of digital electronics.[66]
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The next great advance in computing power came with the advent of the integrated circuit (IC).
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The idea of the integrated circuit was first conceived by a radar scientist working for the Royal Radar Establishment of the Ministry of Defence, Geoffrey W.A. Dummer. Dummer presented the first public description of an integrated circuit at the Symposium on Progress in Quality Electronic Components in Washington, D.C. on 7 May 1952.[67]
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The first working ICs were invented by Jack Kilby at Texas Instruments and Robert Noyce at Fairchild Semiconductor.[68] Kilby recorded his initial ideas concerning the integrated circuit in July 1958, successfully demonstrating the first working integrated example on 12 September 1958.[69] In his patent application of 6 February 1959, Kilby described his new device as "a body of semiconductor material ... wherein all the components of the electronic circuit are completely integrated".[70][71] However, Kilby's invention was a hybrid integrated circuit (hybrid IC), rather than a monolithic integrated circuit (IC) chip.[72] Kilby's IC had external wire connections, which made it difficult to mass-produce.[73]
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Noyce also came up with his own idea of an integrated circuit half a year later than Kilby.[74] Noyce's invention was the first true monolithic IC chip.[75][73] His chip solved many practical problems that Kilby's had not. Produced at Fairchild Semiconductor, it was made of silicon, whereas Kilby's chip was made of germanium. Noyce's monolithic IC was fabricated using the planar process, developed by his colleague Jean Hoerni in early 1959. In turn, the planar process was based on the silicon surface passivation and thermal oxidation processes developed by Mohamed Atalla at Bell Labs in the late 1950s.[76][77][78]
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Modern monolithic ICs are predominantly MOS (metal-oxide-semiconductor) integrated circuits, built from MOSFETs (MOS transistors).[79] After the first MOSFET was invented by Mohamed Atalla and Dawon Kahng at Bell Labs in 1959,[80] Atalla first proposed the concept of the MOS integrated circuit in 1960, followed by Kahng in 1961, both noting that the MOS transistor's ease of fabrication made it useful for integrated circuits.[51][81] The earliest experimental MOS IC to be fabricated was a 16-transistor chip built by Fred Heiman and Steven Hofstein at RCA in 1962.[82] General Microelectronics later introduced the first commercial MOS IC in 1964,[83] developed by Robert Norman.[82] Following the development of the self-aligned gate (silicon-gate) MOS transistor by Robert Kerwin, Donald Klein and John Sarace at Bell Labs in 1967, the first silicon-gate MOS IC with self-aligned gates was developed by Federico Faggin at Fairchild Semiconductor in 1968.[84] The MOSFET has since become the most critical device component in modern ICs.[85]
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The development of the MOS integrated circuit led to the invention of the microprocessor,[86][87] and heralded an explosion in the commercial and personal use of computers. While the subject of exactly which device was the first microprocessor is contentious, partly due to lack of agreement on the exact definition of the term "microprocessor", it is largely undisputed that the first single-chip microprocessor was the Intel 4004,[88] designed and realized by Federico Faggin with his silicon-gate MOS IC technology,[86] along with Ted Hoff, Masatoshi Shima and Stanley Mazor at Intel.[89][90] In the early 1970s, MOS IC technology enabled the integration of more than 10,000 transistors on a single chip.[59]
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System on a Chip (SoCs) are complete computers on a microchip (or chip) the size of a coin.[91] They may or may not have integrated RAM and flash memory. If not integrated, The RAM is usually placed directly above (known as Package on package) or below (on the opposite side of the circuit board) the SoC, and the flash memory is usually placed right next to the SoC, this all done to improve data transfer speeds, as the data signals don't have to travel long distances. Since ENIAC in 1945, computers have advanced enormously, with modern SoCs (Such as the Snapdragon 865) being the size of a coin while also being hundreds of thousands of times more powerful than ENIAC, integrating billions of transistors, and consuming only a few watts of power.
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The first mobile computers were heavy and ran from mains power. The 50lb IBM 5100 was an early example. Later portables such as the Osborne 1 and Compaq Portable were considerably lighter but still needed to be plugged in. The first laptops, such as the Grid Compass, removed this requirement by incorporating batteries – and with the continued miniaturization of computing resources and advancements in portable battery life, portable computers grew in popularity in the 2000s.[92] The same developments allowed manufacturers to integrate computing resources into cellular mobile phones by the early 2000s.
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These smartphones and tablets run on a variety of operating systems and recently became the dominant computing device on the market.[93] These are powered by System on a Chip (SoCs), which are complete computers on a microchip the size of a coin.[91]
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Computers can be classified in a number of different ways, including:
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The term hardware covers all of those parts of a computer that are tangible physical objects. Circuits, computer chips, graphic cards, sound cards, memory (RAM), motherboard, displays, power supplies, cables, keyboards, printers and "mice" input devices are all hardware.
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A general purpose computer has four main components: the arithmetic logic unit (ALU), the control unit, the memory, and the input and output devices (collectively termed I/O). These parts are interconnected by buses, often made of groups of wires.
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Inside each of these parts are thousands to trillions of small electrical circuits which can be turned off or on by means of an electronic switch. Each circuit represents a bit (binary digit) of information so that when the circuit is on it represents a "1", and when off it represents a "0" (in positive logic representation). The circuits are arranged in logic gates so that one or more of the circuits may control the state of one or more of the other circuits.
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When unprocessed data is sent to the computer with the help of input devices, the data is processed and sent to output devices. The input devices may be hand-operated or automated. The act of processing is mainly regulated by the CPU. Some examples of input devices are:
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The means through which computer gives output are known as output devices. Some examples of output devices are:
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The control unit (often called a control system or central controller) manages the computer's various components; it reads and interprets (decodes) the program instructions, transforming them into control signals that activate other parts of the computer.[95] Control systems in advanced computers may change the order of execution of some instructions to improve performance.
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A key component common to all CPUs is the program counter, a special memory cell (a register) that keeps track of which location in memory the next instruction is to be read from.[96]
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The control system's function is as follows—note that this is a simplified description, and some of these steps may be performed concurrently or in a different order depending on the type of CPU:
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Since the program counter is (conceptually) just another set of memory cells, it can be changed by calculations done in the ALU. Adding 100 to the program counter would cause the next instruction to be read from a place 100 locations further down the program. Instructions that modify the program counter are often known as "jumps" and allow for loops (instructions that are repeated by the computer) and often conditional instruction execution (both examples of control flow).
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The sequence of operations that the control unit goes through to process an instruction is in itself like a short computer program, and indeed, in some more complex CPU designs, there is another yet smaller computer called a microsequencer, which runs a microcode program that causes all of these events to happen.
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The control unit, ALU, and registers are collectively known as a central processing unit (CPU). Early CPUs were composed of many separate components. Since the 1970s, CPUs have typically been constructed on a single MOS integrated circuit chip called a microprocessor.
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The ALU is capable of performing two classes of operations: arithmetic and logic.[97] The set of arithmetic operations that a particular ALU supports may be limited to addition and subtraction, or might include multiplication, division, trigonometry functions such as sine, cosine, etc., and square roots. Some can only operate on whole numbers (integers) while others use floating point to represent real numbers, albeit with limited precision. However, any computer that is capable of performing just the simplest operations can be programmed to break down the more complex operations into simple steps that it can perform. Therefore, any computer can be programmed to perform any arithmetic operation—although it will take more time to do so if its ALU does not directly support the operation. An ALU may also compare numbers and return boolean truth values (true or false) depending on whether one is equal to, greater than or less than the other ("is 64 greater than 65?"). Logic operations involve Boolean logic: AND, OR, XOR, and NOT. These can be useful for creating complicated conditional statements and processing boolean logic.
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Superscalar computers may contain multiple ALUs, allowing them to process several instructions simultaneously.[98] Graphics processors and computers with SIMD and MIMD features often contain ALUs that can perform arithmetic on vectors and matrices.
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A computer's memory can be viewed as a list of cells into which numbers can be placed or read. Each cell has a numbered "address" and can store a single number. The computer can be instructed to "put the number 123 into the cell numbered 1357" or to "add the number that is in cell 1357 to the number that is in cell 2468 and put the answer into cell 1595." The information stored in memory may represent practically anything. Letters, numbers, even computer instructions can be placed into memory with equal ease. Since the CPU does not differentiate between different types of information, it is the software's responsibility to give significance to what the memory sees as nothing but a series of numbers.
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In almost all modern computers, each memory cell is set up to store binary numbers in groups of eight bits (called a byte). Each byte is able to represent 256 different numbers (28 = 256); either from 0 to 255 or −128 to +127. To store larger numbers, several consecutive bytes may be used (typically, two, four or eight). When negative numbers are required, they are usually stored in two's complement notation. Other arrangements are possible, but are usually not seen outside of specialized applications or historical contexts. A computer can store any kind of information in memory if it can be represented numerically. Modern computers have billions or even trillions of bytes of memory.
|
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The CPU contains a special set of memory cells called registers that can be read and written to much more rapidly than the main memory area. There are typically between two and one hundred registers depending on the type of CPU. Registers are used for the most frequently needed data items to avoid having to access main memory every time data is needed. As data is constantly being worked on, reducing the need to access main memory (which is often slow compared to the ALU and control units) greatly increases the computer's speed.
|
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Computer main memory comes in two principal varieties:
|
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RAM can be read and written to anytime the CPU commands it, but ROM is preloaded with data and software that never changes, therefore the CPU can only read from it. ROM is typically used to store the computer's initial start-up instructions. In general, the contents of RAM are erased when the power to the computer is turned off, but ROM retains its data indefinitely. In a PC, the ROM contains a specialized program called the BIOS that orchestrates loading the computer's operating system from the hard disk drive into RAM whenever the computer is turned on or reset. In embedded computers, which frequently do not have disk drives, all of the required software may be stored in ROM. Software stored in ROM is often called firmware, because it is notionally more like hardware than software. Flash memory blurs the distinction between ROM and RAM, as it retains its data when turned off but is also rewritable. It is typically much slower than conventional ROM and RAM however, so its use is restricted to applications where high speed is unnecessary.[99]
|
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+
|
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In more sophisticated computers there may be one or more RAM cache memories, which are slower than registers but faster than main memory. Generally computers with this sort of cache are designed to move frequently needed data into the cache automatically, often without the need for any intervention on the programmer's part.
|
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|
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I/O is the means by which a computer exchanges information with the outside world.[100] Devices that provide input or output to the computer are called peripherals.[101] On a typical personal computer, peripherals include input devices like the keyboard and mouse, and output devices such as the display and printer. Hard disk drives, floppy disk drives and optical disc drives serve as both input and output devices. Computer networking is another form of I/O.
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I/O devices are often complex computers in their own right, with their own CPU and memory. A graphics processing unit might contain fifty or more tiny computers that perform the calculations necessary to display 3D graphics.[citation needed] Modern desktop computers contain many smaller computers that assist the main CPU in performing I/O. A 2016-era flat screen display contains its own computer circuitry.
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While a computer may be viewed as running one gigantic program stored in its main memory, in some systems it is necessary to give the appearance of running several programs simultaneously. This is achieved by multitasking i.e. having the computer switch rapidly between running each program in turn.[102] One means by which this is done is with a special signal called an interrupt, which can periodically cause the computer to stop executing instructions where it was and do something else instead. By remembering where it was executing prior to the interrupt, the computer can return to that task later. If several programs are running "at the same time". then the interrupt generator might be causing several hundred interrupts per second, causing a program switch each time. Since modern computers typically execute instructions several orders of magnitude faster than human perception, it may appear that many programs are running at the same time even though only one is ever executing in any given instant. This method of multitasking is sometimes termed "time-sharing" since each program is allocated a "slice" of time in turn.[103]
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Before the era of inexpensive computers, the principal use for multitasking was to allow many people to share the same computer. Seemingly, multitasking would cause a computer that is switching between several programs to run more slowly, in direct proportion to the number of programs it is running, but most programs spend much of their time waiting for slow input/output devices to complete their tasks. If a program is waiting for the user to click on the mouse or press a key on the keyboard, then it will not take a "time slice" until the event it is waiting for has occurred. This frees up time for other programs to execute so that many programs may be run simultaneously without unacceptable speed loss.
|
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|
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Some computers are designed to distribute their work across several CPUs in a multiprocessing configuration, a technique once employed only in large and powerful machines such as supercomputers, mainframe computers and servers. Multiprocessor and multi-core (multiple CPUs on a single integrated circuit) personal and laptop computers are now widely available, and are being increasingly used in lower-end markets as a result.
|
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Supercomputers in particular often have highly unique architectures that differ significantly from the basic stored-program architecture and from general purpose computers.[104] They often feature thousands of CPUs, customized high-speed interconnects, and specialized computing hardware. Such designs tend to be useful only for specialized tasks due to the large scale of program organization required to successfully utilize most of the available resources at once. Supercomputers usually see usage in large-scale simulation, graphics rendering, and cryptography applications, as well as with other so-called "embarrassingly parallel" tasks.
|
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Software refers to parts of the computer which do not have a material form, such as programs, data, protocols, etc. Software is that part of a computer system that consists of encoded information or computer instructions, in contrast to the physical hardware from which the system is built. Computer software includes computer programs, libraries and related non-executable data, such as online documentation or digital media. It is often divided into system software and application software Computer hardware and software require each other and neither can be realistically used on its own. When software is stored in hardware that cannot easily be modified, such as with BIOS ROM in an IBM PC compatible computer, it is sometimes called "firmware".
|
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|
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There are thousands of different programming languages—some intended to be general purpose, others useful only for highly specialized applications.
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The defining feature of modern computers which distinguishes them from all other machines is that they can be programmed. That is to say that some type of instructions (the program) can be given to the computer, and it will process them. Modern computers based on the von Neumann architecture often have machine code in the form of an imperative programming language. In practical terms, a computer program may be just a few instructions or extend to many millions of instructions, as do the programs for word processors and web browsers for example. A typical modern computer can execute billions of instructions per second (gigaflops) and rarely makes a mistake over many years of operation. Large computer programs consisting of several million instructions may take teams of programmers years to write, and due to the complexity of the task almost certainly contain errors.
|
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|
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This section applies to most common RAM machine–based computers.
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In most cases, computer instructions are simple: add one number to another, move some data from one location to another, send a message to some external device, etc. These instructions are read from the computer's memory and are generally carried out (executed) in the order they were given. However, there are usually specialized instructions to tell the computer to jump ahead or backwards to some other place in the program and to carry on executing from there. These are called "jump" instructions (or branches). Furthermore, jump instructions may be made to happen conditionally so that different sequences of instructions may be used depending on the result of some previous calculation or some external event. Many computers directly support subroutines by providing a type of jump that "remembers" the location it jumped from and another instruction to return to the instruction following that jump instruction.
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Program execution might be likened to reading a book. While a person will normally read each word and line in sequence, they may at times jump back to an earlier place in the text or skip sections that are not of interest. Similarly, a computer may sometimes go back and repeat the instructions in some section of the program over and over again until some internal condition is met. This is called the flow of control within the program and it is what allows the computer to perform tasks repeatedly without human intervention.
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Comparatively, a person using a pocket calculator can perform a basic arithmetic operation such as adding two numbers with just a few button presses. But to add together all of the numbers from 1 to 1,000 would take thousands of button presses and a lot of time, with a near certainty of making a mistake. On the other hand, a computer may be programmed to do this with just a few simple instructions. The following example is written in the MIPS assembly language:
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Once told to run this program, the computer will perform the repetitive addition task without further human intervention. It will almost never make a mistake and a modern PC can complete the task in a fraction of a second.
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In most computers, individual instructions are stored as machine code with each instruction being given a unique number (its operation code or opcode for short). The command to add two numbers together would have one opcode; the command to multiply them would have a different opcode, and so on. The simplest computers are able to perform any of a handful of different instructions; the more complex computers have several hundred to choose from, each with a unique numerical code. Since the computer's memory is able to store numbers, it can also store the instruction codes. This leads to the important fact that entire programs (which are just lists of these instructions) can be represented as lists of numbers and can themselves be manipulated inside the computer in the same way as numeric data. The fundamental concept of storing programs in the computer's memory alongside the data they operate on is the crux of the von Neumann, or stored program[citation needed], architecture. In some cases, a computer might store some or all of its program in memory that is kept separate from the data it operates on. This is called the Harvard architecture after the Harvard Mark I computer. Modern von Neumann computers display some traits of the Harvard architecture in their designs, such as in CPU caches.
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While it is possible to write computer programs as long lists of numbers (machine language) and while this technique was used with many early computers,[105] it is extremely tedious and potentially error-prone to do so in practice, especially for complicated programs. Instead, each basic instruction can be given a short name that is indicative of its function and easy to remember – a mnemonic such as ADD, SUB, MULT or JUMP. These mnemonics are collectively known as a computer's assembly language. Converting programs written in assembly language into something the computer can actually understand (machine language) is usually done by a computer program called an assembler.
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Programming languages provide various ways of specifying programs for computers to run. Unlike natural languages, programming languages are designed to permit no ambiguity and to be concise. They are purely written languages and are often difficult to read aloud. They are generally either translated into machine code by a compiler or an assembler before being run, or translated directly at run time by an interpreter. Sometimes programs are executed by a hybrid method of the two techniques.
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Machine languages and the assembly languages that represent them (collectively termed low-level programming languages) are generally unique to the particular architecture of a computer's central processing unit (CPU). For instance, an ARM architecture CPU (such as may be found in a smartphone or a hand-held videogame) cannot understand the machine language of an x86 CPU that might be in a PC.[106] Historically a significant number of other cpu architectures were created and saw extensive use, notably including the MOS Technology 6502 and 6510 in addition to the Zilog Z80.
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Although considerably easier than in machine language, writing long programs in assembly language is often difficult and is also error prone. Therefore, most practical programs are written in more abstract high-level programming languages that are able to express the needs of the programmer more conveniently (and thereby help reduce programmer error). High level languages are usually "compiled" into machine language (or sometimes into assembly language and then into machine language) using another computer program called a compiler.[107] High level languages are less related to the workings of the target computer than assembly language, and more related to the language and structure of the problem(s) to be solved by the final program. It is therefore often possible to use different compilers to translate the same high level language program into the machine language of many different types of computer. This is part of the means by which software like video games may be made available for different computer architectures such as personal computers and various video game consoles.
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Program design of small programs is relatively simple and involves the analysis of the problem, collection of inputs, using the programming constructs within languages, devising or using established procedures and algorithms, providing data for output devices and solutions to the problem as applicable. As problems become larger and more complex, features such as subprograms, modules, formal documentation, and new paradigms such as object-oriented programming are encountered. Large programs involving thousands of line of code and more require formal software methodologies.
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The task of developing large software systems presents a significant intellectual challenge. Producing software with an acceptably high reliability within a predictable schedule and budget has historically been difficult; the academic and professional discipline of software engineering concentrates specifically on this challenge.
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Errors in computer programs are called "bugs". They may be benign and not affect the usefulness of the program, or have only subtle effects. But in some cases, they may cause the program or the entire system to "hang", becoming unresponsive to input such as mouse clicks or keystrokes, to completely fail, or to crash. Otherwise benign bugs may sometimes be harnessed for malicious intent by an unscrupulous user writing an exploit, code designed to take advantage of a bug and disrupt a computer's proper execution. Bugs are usually not the fault of the computer. Since computers merely execute the instructions they are given, bugs are nearly always the result of programmer error or an oversight made in the program's design.[108]
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Admiral Grace Hopper, an American computer scientist and developer of the first compiler, is credited for having first used the term "bugs" in computing after a dead moth was found shorting a relay in the Harvard Mark II computer in September 1947.[109]
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Computers have been used to coordinate information between multiple locations since the 1950s. The U.S. military's SAGE system was the first large-scale example of such a system, which led to a number of special-purpose commercial systems such as Sabre.[110] In the 1970s, computer engineers at research institutions throughout the United States began to link their computers together using telecommunications technology. The effort was funded by ARPA (now DARPA), and the computer network that resulted was called the ARPANET.[111] The technologies that made the Arpanet possible spread and evolved.
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In time, the network spread beyond academic and military institutions and became known as the Internet. The emergence of networking involved a redefinition of the nature and boundaries of the computer. Computer operating systems and applications were modified to include the ability to define and access the resources of other computers on the network, such as peripheral devices, stored information, and the like, as extensions of the resources of an individual computer. Initially these facilities were available primarily to people working in high-tech environments, but in the 1990s the spread of applications like e-mail and the World Wide Web, combined with the development of cheap, fast networking technologies like Ethernet and ADSL saw computer networking become almost ubiquitous. In fact, the number of computers that are networked is growing phenomenally. A very large proportion of personal computers regularly connect to the Internet to communicate and receive information. "Wireless" networking, often utilizing mobile phone networks, has meant networking is becoming increasingly ubiquitous even in mobile computing environments.
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A computer does not need to be electronic, nor even have a processor, nor RAM, nor even a hard disk. While popular usage of the word "computer" is synonymous with a personal electronic computer, the modern[112] definition of a computer is literally: "A device that computes, especially a programmable [usually] electronic machine that performs high-speed mathematical or logical operations or that assembles, stores, correlates, or otherwise processes information."[113] Any device which processes information qualifies as a computer, especially if the processing is purposeful.[citation needed]
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There is active research to make computers out of many promising new types of technology, such as optical computers, DNA computers, neural computers, and quantum computers. Most computers are universal, and are able to calculate any computable function, and are limited only by their memory capacity and operating speed. However different designs of computers can give very different performance for particular problems; for example quantum computers can potentially break some modern encryption algorithms (by quantum factoring) very quickly.
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There are many types of computer architectures:
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Of all these abstract machines, a quantum computer holds the most promise for revolutionizing computing.[114] Logic gates are a common abstraction which can apply to most of the above digital or analog paradigms. The ability to store and execute lists of instructions called programs makes computers extremely versatile, distinguishing them from calculators. The Church–Turing thesis is a mathematical statement of this versatility: any computer with a minimum capability (being Turing-complete) is, in principle, capable of performing the same tasks that any other computer can perform. Therefore, any type of computer (netbook, supercomputer, cellular automaton, etc.) is able to perform the same computational tasks, given enough time and storage capacity.
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A computer will solve problems in exactly the way it is programmed to, without regard to efficiency, alternative solutions, possible shortcuts, or possible errors in the code. Computer programs that learn and adapt are part of the emerging field of artificial intelligence and machine learning. Artificial intelligence based products generally fall into two major categories: rule based systems and pattern recognition systems. Rule based systems attempt to represent the rules used by human experts and tend to be expensive to develop. Pattern based systems use data about a problem to generate conclusions. Examples of pattern based systems include voice recognition, font recognition, translation and the emerging field of on-line marketing.
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As the use of computers has spread throughout society, there are an increasing number of careers involving computers.
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The need for computers to work well together and to be able to exchange information has spawned the need for many standards organizations, clubs and societies of both a formal and informal nature.
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An alphabet is a standardized set of basic written symbols or graphemes (called letters) that represent the phonemes of certain spoken languages. Not all writing systems represent language in this way; in a syllabary, each character represents a syllable, for instance, and logographic systems use characters to represent words, morphemes, or other semantic units.
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The first fully phonemic script, the Proto-Canaanite script, later known as the Phoenician alphabet, is considered to be the first alphabet, and is the ancestor of most modern alphabets, including Arabic, Cyrillic, Greek, Hebrew, Latin, and possibly Brahmic.[1][2] It was created by Semitic-speaking workers and slaves in the Sinai Peninsula (as the Proto-Sinaitic script), by selecting a small number of hieroglyphs commonly seen in their Egyptian surroundings to describe the sounds, as opposed to the semantic values, of their own Canaanite language.[3][4] Peter T. Daniels, however, distinguishes an abugida or alphasyllabary, a set of graphemes that represent consonantal base letters which diacritics modify to represent vowels (as in Devanagari and other South Asian scripts), an abjad, in which letters predominantly or exclusively represent consonants (as in the original Phoenician, Hebrew or Arabic), and an "alphabet", a set of graphemes that represent both vowels and consonants. In this narrow sense of the word the first "true" alphabet was the Greek alphabet,[5][6] which was developed on the basis of the earlier Phoenician alphabet.
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Of the dozens of alphabets in use today, the most popular is the Latin alphabet,[7] which was derived from the Greek, and which many languages modify by adding letters formed using diacritical marks. While most alphabets have letters composed of lines (linear writing), there are also exceptions such as the alphabets used in Braille. The Khmer alphabet (for Cambodian) is the longest, with 74 letters.[8]
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Alphabets are usually associated with a standard ordering of letters. This makes them useful for purposes of collation, specifically by allowing words to be sorted in alphabetical order. It also means that their letters can be used as an alternative method of "numbering" ordered items, in such contexts as numbered lists and number placements.
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The English word alphabet came into Middle English from the Late Latin word alphabetum, which in turn originated in the Greek ἀλφάβητος (alphabētos). The Greek word was made from the first two letters, alpha(α) and beta(β).[9] The names for the Greek letters came from the first two letters of the Phoenician alphabet; aleph, which also meant ox, and bet, which also meant house.
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Sometimes, like in the alphabet song in English, the term "ABCs" is used instead of the word "alphabet" (Now I know my ABCs...). "Knowing one's ABCs", in general, can be used as a metaphor for knowing the basics about anything.
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The history of the alphabet started in ancient Egypt. Egyptian writing had a set of some 24 hieroglyphs that are called uniliterals,[10] to represent syllables that begin with a single consonant of their language, plus a vowel (or no vowel) to be supplied by the native speaker. These glyphs were used as pronunciation guides for logograms, to write grammatical inflections, and, later, to transcribe loan words and foreign names.[11]
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In the Middle Bronze Age, an apparently "alphabetic" system known as the Proto-Sinaitic script appears in Egyptian turquoise mines in the Sinai peninsula dated to circa the 15th century BC, apparently left by Canaanite workers. In 1999, John and Deborah Darnell discovered an even earlier version of this first alphabet at Wadi el-Hol dated to circa 1800 BC and showing evidence of having been adapted from specific forms of Egyptian hieroglyphs that could be dated to circa 2000 BC, strongly suggesting that the first alphabet had been developed about that time.[12] Based on letter appearances and names, it is believed to be based on Egyptian hieroglyphs.[1] This script had no characters representing vowels, although originally it probably was a syllabary, but unneeded symbols were discarded. An alphabetic cuneiform script with 30 signs including three that indicate the following vowel was invented in Ugarit before the 15th century BC. This script was not used after the destruction of Ugarit.[13]
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The Proto-Sinaitic script eventually developed into the Phoenician alphabet, which is conventionally called "Proto-Canaanite" before ca. 1050 BC.[2] The oldest text in Phoenician script is an inscription on the sarcophagus of King Ahiram. This script is the parent script of all western alphabets. By the tenth century, two other forms can be distinguished, namely Canaanite and Aramaic. The Aramaic gave rise to the Hebrew script.[14] The South Arabian alphabet, a sister script to the Phoenician alphabet, is the script from which the Ge'ez alphabet (an abugida) is descended. Vowelless alphabets are called abjads, currently exemplified in scripts including Arabic, Hebrew, and Syriac. The omission of vowels was not always a satisfactory solution and some "weak" consonants are sometimes used to indicate the vowel quality of a syllable (matres lectionis). These letters have a dual function since they are also used as pure consonants.[15]
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The Proto-Sinaitic or Proto-Canaanite script and the Ugaritic script were the first scripts with a limited number of signs, in contrast to the other widely used writing systems at the time, Cuneiform, Egyptian hieroglyphs, and Linear B. The Phoenician script was probably the first phonemic script[1][2] and it contained only about two dozen distinct letters, making it a script simple enough for common traders to learn. Another advantage of Phoenician was that it could be used to write down many different languages, since it recorded words phonemically.
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25 |
+
The script was spread by the Phoenicians across the Mediterranean.[2] In Greece, the script was modified to add vowels, giving rise to the ancestor of all alphabets in the West. It was the first alphabet in which vowels have independent letter forms separate from those of consonants. The Greeks chose letters representing sounds that did not exist in Greek to represent vowels. Vowels are significant in the Greek language, and the syllabical Linear B script that was used by the Mycenaean Greeks from the 16th century BC had 87 symbols, including 5 vowels. In its early years, there were many variants of the Greek alphabet, a situation that caused many different alphabets to evolve from it.
|
26 |
+
|
27 |
+
The Greek alphabet, in its Euboean form, was carried over by Greek colonists to the Italian peninsula, where it gave rise to a variety of alphabets used to write the Italic languages. One of these became the Latin alphabet, which was spread across Europe as the Romans expanded their empire. Even after the fall of the Roman state, the alphabet survived in intellectual and religious works. It eventually became used for the descendant languages of Latin (the Romance languages) and then for most of the other languages of Europe.
|
28 |
+
|
29 |
+
Some adaptations of the Latin alphabet are augmented with ligatures, such as æ in Danish and Icelandic and Ȣ in Algonquian; by borrowings from other alphabets, such as the thorn þ in Old English and Icelandic, which came from the Futhark runes; and by modifying existing letters, such as the eth ð of Old English and Icelandic, which is a modified d. Other alphabets only use a subset of the Latin alphabet, such as Hawaiian, and Italian, which uses the letters j, k, x, y and w only in foreign words.
|
30 |
+
|
31 |
+
Another notable script is Elder Futhark, which is believed to have evolved out of one of the Old Italic alphabets. Elder Futhark gave rise to a variety of alphabets known collectively as the Runic alphabets. The Runic alphabets were used for Germanic languages from AD 100 to the late Middle Ages. Its usage is mostly restricted to engravings on stone and jewelry, although inscriptions have also been found on bone and wood. These alphabets have since been replaced with the Latin alphabet, except for decorative usage for which the runes remained in use until the 20th century.
|
32 |
+
|
33 |
+
The Old Hungarian script is a contemporary writing system of the Hungarians. It was in use during the entire history of Hungary, albeit not as an official writing system. From the 19th century it once again became more and more popular.
|
34 |
+
|
35 |
+
The Glagolitic alphabet was the initial script of the liturgical language Old Church Slavonic and became, together with the Greek uncial script, the basis of the Cyrillic script. Cyrillic is one of the most widely used modern alphabetic scripts, and is notable for its use in Slavic languages and also for other languages within the former Soviet Union. Cyrillic alphabets include the Serbian, Macedonian, Bulgarian, Russian, Belarusian and Ukrainian. The Glagolitic alphabet is believed to have been created by Saints Cyril and Methodius, while the Cyrillic alphabet was invented by Clement of Ohrid, who was their disciple. They feature many letters that appear to have been borrowed from or influenced by the Greek alphabet and the Hebrew alphabet.
|
36 |
+
|
37 |
+
The longest European alphabet is the Latin-derived Slovak alphabet which has 46 letters.
|
38 |
+
|
39 |
+
Beyond the logographic Chinese writing, many phonetic scripts are in existence in Asia. The Arabic alphabet, Hebrew alphabet, Syriac alphabet, and other abjads of the Middle East are developments of the Aramaic alphabet.
|
40 |
+
|
41 |
+
Most alphabetic scripts of India and Eastern Asia are descended from the Brahmi script, which is often believed to be a descendant of Aramaic.
|
42 |
+
|
43 |
+
In Korea, the Hangul alphabet was created by Sejong the Great.[16] Hangul is a unique alphabet: it is a featural alphabet, where many of the letters are designed from a sound's place of articulation (P to look like the widened mouth, L to look like the tongue pulled in, etc.); its design was planned by the government of the day; and it places individual letters in syllable clusters with equal dimensions, in the same way as Chinese characters, to allow for mixed-script writing[17] (one syllable always takes up one type-space no matter how many letters get stacked into building that one sound-block).
|
44 |
+
|
45 |
+
Zhuyin (sometimes called Bopomofo) is a semi-syllabary used to phonetically transcribe Mandarin Chinese in the Republic of China. After the later establishment of the People's Republic of China and its adoption of Hanyu Pinyin, the use of Zhuyin today is limited, but it is still widely used in Taiwan where the Republic of China still governs. Zhuyin developed out of a form of Chinese shorthand based on Chinese characters in the early 1900s and has elements of both an alphabet and a syllabary. Like an alphabet the phonemes of syllable initials are represented by individual symbols, but like a syllabary the phonemes of the syllable finals are not; rather, each possible final (excluding the medial glide) is represented by its own symbol. For example, luan is represented as ㄌㄨㄢ (l-u-an), where the last symbol ㄢ represents the entire final -an. While Zhuyin is not used as a mainstream writing system, it is still often used in ways similar to a romanization system—that is, for aiding in pronunciation and as an input method for Chinese characters on computers and cellphones.
|
46 |
+
|
47 |
+
European alphabets, especially Latin and Cyrillic, have been adapted for many languages of Asia. Arabic is also widely used, sometimes as an abjad (as with Urdu and Persian) and sometimes as a complete alphabet (as with Kurdish and Uyghur).
|
48 |
+
|
49 |
+
Egyptian hieroglyphs 32 c. BCE
|
50 |
+
|
51 |
+
Hangul 1443
|
52 |
+
|
53 |
+
The term "alphabet" is used by linguists and paleographers in both a wide and a narrow sense. In the wider sense, an alphabet is a script that is segmental at the phoneme level—that is, it has separate glyphs for individual sounds and not for larger units such as syllables or words. In the narrower sense, some scholars distinguish "true" alphabets from two other types of segmental script, abjads and abugidas. These three differ from each other in the way they treat vowels: abjads have letters for consonants and leave most vowels unexpressed; abugidas are also consonant-based, but indicate vowels with diacritics to or a systematic graphic modification of the consonants. In alphabets in the narrow sense, on the other hand, consonants and vowels are written as independent letters.[18] The earliest known alphabet in the wider sense is the Wadi el-Hol script, believed to be an abjad, which through its successor Phoenician is the ancestor of modern alphabets, including Arabic, Greek, Latin (via the Old Italic alphabet), Cyrillic (via the Greek alphabet) and Hebrew (via Aramaic).
|
54 |
+
|
55 |
+
Examples of present-day abjads are the Arabic and Hebrew scripts; true alphabets include Latin, Cyrillic, and Korean hangul; and abugidas are used to write Tigrinya, Amharic, Hindi, and Thai. The Canadian Aboriginal syllabics are also an abugida rather than a syllabary as their name would imply, since each glyph stands for a consonant that is modified by rotation to represent the following vowel. (In a true syllabary, each consonant-vowel combination would be represented by a separate glyph.)
|
56 |
+
|
57 |
+
All three types may be augmented with syllabic glyphs. Ugaritic, for example, is basically an abjad, but has syllabic letters for /ʔa, ʔi, ʔu/. (These are the only time vowels are indicated.) Cyrillic is basically a true alphabet, but has syllabic letters for /ja, je, ju/ (я, е, ю); Coptic has a letter for /ti/. Devanagari is typically an abugida augmented with dedicated letters for initial vowels, though some traditions use अ as a zero consonant as the graphic base for such vowels.
|
58 |
+
|
59 |
+
The boundaries between the three types of segmental scripts are not always clear-cut. For example, Sorani Kurdish is written in the Arabic script, which is normally an abjad. However, in Kurdish, writing the vowels is mandatory, and full letters are used, so the script is a true alphabet. Other languages may use a Semitic abjad with mandatory vowel diacritics, effectively making them abugidas. On the other hand, the Phagspa script of the Mongol Empire was based closely on the Tibetan abugida, but all vowel marks were written after the preceding consonant rather than as diacritic marks. Although short a was not written, as in the Indic abugidas, one could argue that the linear arrangement made this a true alphabet. Conversely, the vowel marks of the Tigrinya abugida and the Amharic abugida (ironically, the original source of the term "abugida") have been so completely assimilated into their consonants that the modifications are no longer systematic and have to be learned as a syllabary rather than as a segmental script. Even more extreme, the Pahlavi abjad eventually became logographic. (See below.)
|
60 |
+
|
61 |
+
Thus the primary classification of alphabets reflects how they treat vowels. For tonal languages, further classification can be based on their treatment of tone, though names do not yet exist to distinguish the various types. Some alphabets disregard tone entirely, especially when it does not carry a heavy functional load, as in Somali and many other languages of Africa and the Americas. Such scripts are to tone what abjads are to vowels. Most commonly, tones are indicated with diacritics, the way vowels are treated in abugidas. This is the case for Vietnamese (a true alphabet) and Thai (an abugida). In Thai, tone is determined primarily by the choice of consonant, with diacritics for disambiguation. In the Pollard script, an abugida, vowels are indicated by diacritics, but the placement of the diacritic relative to the consonant is modified to indicate the tone. More rarely, a script may have separate letters for tones, as is the case for Hmong and Zhuang. For most of these scripts, regardless of whether letters or diacritics are used, the most common tone is not marked, just as the most common vowel is not marked in Indic abugidas; in Zhuyin not only is one of the tones unmarked, but there is a diacritic to indicate lack of tone, like the virama of Indic.
|
62 |
+
|
63 |
+
The number of letters in an alphabet can be quite small. The Book Pahlavi script, an abjad, had only twelve letters at one point, and may have had even fewer later on. Today the Rotokas alphabet has only twelve letters. (The Hawaiian alphabet is sometimes claimed to be as small, but it actually consists of 18 letters, including the ʻokina and five long vowels. However, Hawaiian Braille has only 13 letters.) While Rotokas has a small alphabet because it has few phonemes to represent (just eleven), Book Pahlavi was small because many letters had been conflated—that is, the graphic distinctions had been lost over time, and diacritics were not developed to compensate for this as they were in Arabic, another script that lost many of its distinct letter shapes. For example, a comma-shaped letter represented g, d, y, k, or j. However, such apparent simplifications can perversely make a script more complicated. In later Pahlavi papyri, up to half of the remaining graphic distinctions of these twelve letters were lost, and the script could no longer be read as a sequence of letters at all, but instead each word had to be learned as a whole—that is, they had become logograms as in Egyptian Demotic.
|
64 |
+
|
65 |
+
The largest segmental script is probably an abugida, Devanagari. When written in Devanagari, Vedic Sanskrit has an alphabet of 53 letters, including the visarga mark for final aspiration and special letters for kš and jñ, though one of the letters is theoretical and not actually used. The Hindi alphabet must represent both Sanskrit and modern vocabulary, and so has been expanded to 58 with the khutma letters (letters with a dot added) to represent sounds from Persian and English. Thai has a total of 59 symbols, consisting of 44 consonants, 13 vowels and 2 syllabics, not including 4 diacritics for tone marks and one for vowel length.
|
66 |
+
|
67 |
+
The largest known abjad is Sindhi, with 51 letters. The largest alphabets in the narrow sense include Kabardian and Abkhaz (for Cyrillic), with 58 and 56 letters, respectively, and Slovak (for the Latin script), with 46. However, these scripts either count di- and tri-graphs as separate letters, as Spanish did with ch and ll until recently, or uses diacritics like Slovak č.
|
68 |
+
|
69 |
+
The Georgian alphabet (Georgian: ანბანი Anbani) is an alphabetic writing system. With 33 letters, it is the largest true alphabet where each letter is graphically independent.[citation needed] The original Georgian alphabet had 38 letters but 5 letters were removed in the 19th century by Ilia Chavchavadze. The Georgian alphabet is much closer to Greek than the other Caucasian alphabets. The letter order parallels the Greek, with the consonants without a Greek equivalent organized at the end of the alphabet. The origins of the alphabet are still unknown. Some Armenian and Western scholars believe it was created by Mesrop Mashtots (Armenian: Մեսրոպ Մաշտոց Mesrop Maštoc') also known as Mesrob the Vartabed, who was an early medieval Armenian linguist, theologian, statesman and hymnologist, best known for inventing the Armenian alphabet c. 405 AD;[19][20] other Georgian[21] and Western[22] scholars are against this theory.
|
70 |
+
|
71 |
+
Syllabaries typically contain 50 to 400 glyphs, and the glyphs of logographic systems typically number from the many hundreds into the thousands. Thus a simple count of the number of distinct symbols is an important clue to the nature of an unknown script.
|
72 |
+
|
73 |
+
The Armenian alphabet (Armenian: Հայոց գրեր Hayots grer or Հայոց այբուբեն Hayots aybuben) is a graphically unique alphabetical writing system that has been used to write the Armenian language. It was created in year 405 A.D. originally contained 36 letters. Two more letters, օ (o) and ֆ (f), were added in the Middle Ages. During the 1920s orthography reform, a new letter և (capital ԵՎ) was added, which was a ligature before ե+ւ, while the letter Ւ ւ was discarded and reintroduced as part of a new letter ՈՒ ու (which was a digraph before).
|
74 |
+
|
75 |
+
The Armenian script's directionality is horizontal left-to-right, like the Latin and Greek alphabets.[23] It also uses bicameral script like those. The Armenian word for "alphabet" is այբուբեն aybuben (Armenian pronunciation: [ɑjbubɛn]), named after the first two letters of the Armenian alphabet Ա այբ ayb and Բ բեն ben.
|
76 |
+
|
77 |
+
Alphabets often come to be associated with a standard ordering of their letters, which can then be used for purposes of collation—namely for the listing of words and other items in what is called alphabetical order.
|
78 |
+
|
79 |
+
The basic ordering of the Latin alphabet (A
|
80 |
+
B
|
81 |
+
C
|
82 |
+
D
|
83 |
+
E
|
84 |
+
F
|
85 |
+
G
|
86 |
+
H
|
87 |
+
I
|
88 |
+
J
|
89 |
+
K
|
90 |
+
L
|
91 |
+
M
|
92 |
+
N
|
93 |
+
O
|
94 |
+
P
|
95 |
+
Q
|
96 |
+
R
|
97 |
+
S
|
98 |
+
T
|
99 |
+
U
|
100 |
+
V
|
101 |
+
W
|
102 |
+
X
|
103 |
+
Y
|
104 |
+
Z), which is derived from the Northwest Semitic "Abgad" order,[24] is well established, although languages using this alphabet have different conventions for their treatment of modified letters (such as the French é, à, and ô) and of certain combinations of letters (multigraphs). In French, these are not considered to be additional letters for the purposes of collation. However, in Icelandic, the accented letters such as á, í, and ö are considered distinct letters representing different vowel sounds from the sounds represented by their unaccented counterparts. In Spanish, ñ is considered a separate letter, but accented vowels such as á and é are not. The ll and ch were also considered single letters, but in 1994 the Real Academia Española changed the collating order so that ll is between lk and lm in the dictionary and ch is between cg and ci, and in 2010 the tenth congress of the Association of Spanish Language Academies changed it so they were no longer letters at all.[25][26]
|
105 |
+
|
106 |
+
In German, words starting with sch- (which spells the German phoneme /ʃ/) are inserted between words with initial sca- and sci- (all incidentally loanwords) instead of appearing after initial sz, as though it were a single letter—in contrast to several languages such as Albanian, in which dh-, ë-, gj-, ll-, rr-, th-, xh- and zh- (all representing phonemes and considered separate single letters) would follow the letters d, e, g, l, n, r, t, x and z respectively, as well as Hungarian and Welsh. Further, German words with an umlaut are collated ignoring the umlaut—contrary to Turkish that adopted the graphemes ö and ü, and where a word like tüfek, would come after tuz, in the dictionary. An exception is the German telephone directory where umlauts are sorted like ä = ae since names such as Jäger also appear with the spelling Jaeger, and are not distinguished in the spoken language.
|
107 |
+
|
108 |
+
The Danish and Norwegian alphabets end with æ—ø—å, whereas the Swedish and Finnish ones conventionally put å—ä—ö at the end.
|
109 |
+
|
110 |
+
It is unknown whether the earliest alphabets had a defined sequence. Some alphabets today, such as the Hanuno'o script, are learned one letter at a time, in no particular order, and are not used for collation where a definite order is required. However, a dozen Ugaritic tablets from the fourteenth century BC preserve the alphabet in two sequences. One, the ABCDE order later used in Phoenician, has continued with minor changes in Hebrew, Greek, Armenian, Gothic, Cyrillic, and Latin; the other, HMĦLQ, was used in southern Arabia and is preserved today in Ethiopic.[27] Both orders have therefore been stable for at least 3000 years.
|
111 |
+
|
112 |
+
Runic used an unrelated Futhark sequence, which was later simplified. Arabic uses its own sequence, although Arabic retains the traditional abjadi order for numbering.
|
113 |
+
|
114 |
+
The Brahmic family of alphabets used in India use a unique order based on phonology: The letters are arranged according to how and where they are produced in the mouth. This organization is used in Southeast Asia, Tibet, Korean hangul, and even Japanese kana, which is not an alphabet.
|
115 |
+
|
116 |
+
The Phoenician letter names, in which each letter was associated with a word that begins with that sound (acrophony), continue to be used to varying degrees in Samaritan, Aramaic, Syriac, Hebrew, Greek and Arabic.
|
117 |
+
|
118 |
+
The names were abandoned in Latin, which instead referred to the letters by adding a vowel (usually e) before or after the consonant; the two exceptions were Y and Z, which were borrowed from the Greek alphabet rather than Etruscan, and were known as Y Graeca "Greek Y" (pronounced I Graeca "Greek I") and zeta (from Greek)—this discrepancy was inherited by many European languages, as in the term zed for Z in all forms of English other than American English. Over time names sometimes shifted or were added, as in double U for W ("double V" in French), the English name for Y, and American zee for Z. Comparing names in English and French gives a clear reflection of the Great Vowel Shift: A, B, C and D are pronounced /eɪ, biː, siː, diː/ in today's English, but in contemporary French they are /a, be, se, de/. The French names (from which the English names are derived) preserve the qualities of the English vowels from before the Great Vowel Shift. By contrast, the names of F, L, M, N and S (/ɛf, ɛl, ɛm, ɛn, ɛs/) remain the same in both languages, because "short" vowels were largely unaffected by the Shift.
|
119 |
+
|
120 |
+
In Cyrillic originally the letters were given names based on Slavic words; this was later abandoned as well in favor of a system similar to that used in Latin.
|
121 |
+
|
122 |
+
Letters of Armenian alphabet also have distinct letter names.
|
123 |
+
|
124 |
+
When an alphabet is adopted or developed to represent a given language, an orthography generally comes into being, providing rules for the spelling of words in that language. In accordance with the principle on which alphabets are based, these rules will generally map letters of the alphabet to the phonemes (significant sounds) of the spoken language. In a perfectly phonemic orthography there would be a consistent one-to-one correspondence between the letters and the phonemes, so that a writer could predict the spelling of a word given its pronunciation, and a speaker would always know the pronunciation of a word given its spelling, and vice versa. However this ideal is not usually achieved in practice; some languages (such as Spanish and Finnish) come close to it, while others (such as English) deviate from it to a much larger degree.
|
125 |
+
|
126 |
+
The pronunciation of a language often evolves independently of its writing system, and writing systems have been borrowed for languages they were not designed for, so the degree to which letters of an alphabet correspond to phonemes of a language varies greatly from one language to another and even within a single language.
|
127 |
+
|
128 |
+
Languages may fail to achieve a one-to-one correspondence between letters and sounds in any of several ways:
|
129 |
+
|
130 |
+
National languages sometimes elect to address the problem of dialects by simply associating the alphabet with the national standard. Some national languages like Finnish, Armenian, Turkish, Russian, Serbo-Croatian (Serbian, Croatian and Bosnian) and Bulgarian have a very regular spelling system with a nearly one-to-one correspondence between letters and phonemes. Strictly speaking, these national languages lack a word corresponding to the verb "to spell" (meaning to split a word into its letters), the closest match being a verb meaning to split a word into its syllables. Similarly, the Italian verb corresponding to 'spell (out)', compitare, is unknown to many Italians because spelling is usually trivial, as Italian spelling is highly phonemic. In standard Spanish, one can tell the pronunciation of a word from its spelling, but not vice versa, as certain phonemes can be represented in more than one way, but a given letter is consistently pronounced. French, with its silent letters and its heavy use of nasal vowels and elision, may seem to lack much correspondence between spelling and pronunciation, but its rules on pronunciation, though complex, are actually consistent and predictable with a fair degree of accuracy.
|
131 |
+
|
132 |
+
At the other extreme are languages such as English, where the pronunciations of many words simply have to be memorized as they do not correspond to the spelling in a consistent way. For English, this is partly because the Great Vowel Shift occurred after the orthography was established, and because English has acquired a large number of loanwords at different times, retaining their original spelling at varying levels. Even English has general, albeit complex, rules that predict pronunciation from spelling, and these rules are successful most of the time; rules to predict spelling from the pronunciation have a higher failure rate.
|
133 |
+
|
134 |
+
Sometimes, countries have the written language undergo a spelling reform to realign the writing with the contemporary spoken language. These can range from simple spelling changes and word forms to switching the entire writing system itself, as when Turkey switched from the Arabic alphabet to a Latin-based Turkish alphabet.
|
135 |
+
|
136 |
+
The standard system of symbols used by linguists to represent sounds in any language, independently of orthography, is called the International Phonetic Alphabet.
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Oregon (/ˈɒr(ɪ)ɡən/ (listen) ORR-(ih)-gən)[7] is a state in the Pacific Northwest region on the West Coast of the United States. The Columbia River delineates much of Oregon's northern boundary with Washington, while the Snake River delineates much of its eastern boundary with Idaho. The 42° north parallel delineates the southern boundary with California and Nevada.
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Oregon has been home to many Indigenous nations for thousands of years. The first European traders, explorers, and settlers began exploring what is now Oregon's Pacific coast in the early-mid 1500s. As early as 1565, the Spanish began sending vessels northeast from the Philippines, riding the Kuroshio Current in a sweeping circular route across the northern part of the Pacific. In 1592, Juan de Fuca undertook detailed mapping and studies of ocean currents in the Pacific Northwest, including the Oregon coast as well as the strait now bearing his name. Spanish ships – 250 in as many years – would typically not land before reaching Cape Mendocino in California, but some landed or wrecked in what is now Oregon. Nehalem tales recount strangers and the discovery of items like chunks of beeswax and a lidded silver vase, likely connected to the 1707 wreck of the San Francisco Xavier.[8]
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In 1843, an autonomous government was formed in the Oregon Country, and the Oregon Territory was created in 1848. Oregon became the 33rd state of the U.S. on February 14, 1859. Today, with 4 million people over 98,000 square miles (250,000 km2), Oregon is the ninth largest and 27th most populous U.S. state. The capital, Salem, is the second-most populous city in Oregon, with 169,798 residents. Portland, with 647,805, ranks as the 26th among U.S. cities. The Portland metropolitan area, which also includes the city of Vancouver, Washington, to the north, ranks the 25th largest metro area in the nation, with a population of 2,453,168.
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Oregon is one of the most geographically diverse states in the U.S.,[9] marked by volcanoes, abundant bodies of water, dense evergreen and mixed forests, as well as high deserts and semi-arid shrublands. At 11,249 feet (3,429 m), Mount Hood, a stratovolcano, is the state's highest point. Oregon's only national park, Crater Lake National Park, comprises the caldera surrounding Crater Lake, the deepest lake in the United States. The state is also home to the single largest organism in the world, Armillaria ostoyae, a fungus that runs beneath 2,200 acres (8.9 km2) of the Malheur National Forest.[10]
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Because of its diverse landscapes and waterways, Oregon's economy is largely powered by various forms of agriculture, fishing, and hydroelectric power. Oregon is also the top timber producer of the contiguous United States, and the timber industry dominated the state's economy in the 20th century.[11] Technology is another one of Oregon's major economic forces, beginning in the 1970s with the establishment of the Silicon Forest and the expansion of Tektronix and Intel. Sportswear company Nike, Inc., headquartered in Beaverton, is the state's largest public corporation with an annual revenue of $30.6 billion.[12]
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The earliest evidence of the name Oregon has Spanish origins. The term "orejón" comes from the historical chronicle Relación de la Alta y Baja California (1598)[13] written by the new Spaniard Rodrigo Montezuma and made reference to the Columbia River when the Spanish explorers penetrated into the actual North American territory that became part of the Viceroyalty of New Spain. This chronicle is the first topographical and linguistic source with respect to the place name Oregon. There are also two other sources with Spanish origins, such as the word oregano, referring to a plant which grows in the southern part of the region. It is possible that the American territory was named by the Spaniards, as there is a stream in Spain called the "Arroyo del Oregón" (which is located in the province of Ciudad Real); it is also possible that the "j" in the Spanish phrase "El Orejón" was later corrupted into a "g",[14] and in context might refer to the 'earful' of the massive Columbia River at its mouth.
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Another early use of the name, spelled Ouragon, was by Major Robert Rogers in a 1765 petition to the Kingdom of Great Britain. The term referred to the then-mythical River of the West (the Columbia River). By 1778, the spelling had shifted to Oregon.[15] Rogers wrote:
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... from the Great Lakes towards the Head of the Mississippi, and from thence to the River called by the Indians Ouragon ...[16]
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One theory is that the name comes from the French word ouragan ("windstorm" or "hurricane"), which was applied to the River of the West based on Native American tales of powerful Chinook winds on the lower Columbia River, or perhaps from firsthand French experience with the Chinook winds of the Great Plains. At the time, the River of the West was thought to rise in western Minnesota and flow west through the Great Plains.[17]
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Joaquin Miller explained in Sunset magazine, in 1904, how Oregon's name was derived:
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The name, Oregon, is rounded down phonetically, from Ouve água—Oragua, Or-a-gon, Oregon—given probably by the same Portuguese navigator that named the Farallones after his first officer, and it literally, in a large way, means cascades: "Hear the waters." You should steam up the Columbia and hear and feel the waters falling out of the clouds of Mount Hood to understand entirely the full meaning of the name Ouve a água, Oregon.[18]
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Another account, endorsed as the "most plausible explanation" in the book Oregon Geographic Names, was advanced by George R. Stewart in a 1944 article in American Speech. According to Stewart, the name came from an engraver's error in a French map published in the early 18th century, on which the Ouisiconsink (Wisconsin) River was spelled "Ouaricon-sint", broken on two lines with the -sint below, so there appeared to be a river flowing to the west named "Ouaricon".
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According to the Oregon Tourism Commission, present-day Oregonians /ˌɒrɪˈɡoʊniənz/[19] pronounce the state's name as "or-uh-gun, never or-ee-gone".[20] After being drafted by the Detroit Lions in 2002, former Oregon Ducks quarterback Joey Harrington distributed "Orygun" stickers to members of the media as a reminder of how to pronounce the name of his home state.[21][22] The stickers are sold by the University of Oregon Bookstore.[23]
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Oregon is 295 miles (475 km) north to south at longest distance, and 395 miles (636 km) east to west. With an area of 98,381 square miles (254,810 km2), Oregon is slightly larger than the United Kingdom. It is the ninth largest state in the United States.[24] Oregon's highest point is the summit of Mount Hood, at 11,249 feet (3,429 m), and its lowest point is the sea level of the Pacific Ocean along the Oregon Coast.[25] Oregon's mean elevation is 3,300 feet (1,006 m). Crater Lake National Park, the state's only national park, is the site of the deepest lake in the United States at 1,943 feet (592 m).[26] Oregon claims the D River as the shortest river in the world,[27] though the state of Montana makes the same claim of its Roe River.[28] Oregon is also home to Mill Ends Park (in Portland),[29] the smallest park in the world at 452 square inches (0.29 m2).
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Oregon is split into eight geographical regions. In Western Oregon: Oregon Coast (west of the Coast Range), the Willamette Valley, Rogue Valley, Cascade Range and Klamath Mountains; and in Central and Eastern Oregon: the Columbia Plateau, the High Desert, and the Blue Mountains.
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Oregon lies in two time zones. Most of Malheur County is in the Mountain Time Zone, while the rest of the state lies in the Pacific Time Zone.
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Western Oregon's mountainous regions, home to three of the most prominent mountain peaks of the United States including Mount Hood, were formed by the volcanic activity of the Juan de Fuca Plate, a tectonic plate that poses a continued threat of volcanic activity and earthquakes in the region. The most recent major activity was the 1700 Cascadia earthquake. Washington's Mount St. Helens erupted in 1980, an event visible from northern Oregon and affecting some areas there.[30]
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The Columbia River, which forms much of Oregon's northern border, also played a major role in the region's geological evolution, as well as its economic and cultural development. The Columbia is one of North America's largest rivers, and one of two rivers to cut through the Cascades (the Klamath River in southern Oregon is the other). About 15,000 years ago, the Columbia repeatedly flooded much of Oregon during the Missoula Floods; the modern fertility of the Willamette Valley is largely the result. Plentiful salmon made parts of the river, such as Celilo Falls, hubs of economic activity for thousands of years.
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Today, Oregon's landscape varies from rain forest in the Coast Range to barren desert in the southeast, which still meets the technical definition of a frontier. Oregon's geographical center is further west than any of the other 48 contiguous states (although the westernmost point of the lower 48 states is in Washington). Central Oregon's geographical features range from high desert and volcanic rock formations resulting from lava beds. The Oregon Badlands Wilderness is in this region of the state.[31]
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Typical of a western state, Oregon is home to a unique and diverse array of wildlife. Roughly 60 percent of the state is covered in forest,[32] while the areas west of the Cascades are more densely populated by forest, making up around 80 percent of the landscape. Some 60 percent of Oregon's forests are within federal land.[32] Oregon is the top timber producer of the lower 48 states.[11][33]
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Moose have not always inhabited the state but came to Oregon in the 1960s; the Wallowa Valley herd numbered about 60 as of 2013[update].[39] Gray wolves were extirpated from Oregon around 1930 but have since found their way back; most reside in northeast Oregon, with two packs living in the south-central part.[40] Although their existence in Oregon is unconfirmed, reports of grizzly bears still turn up, and it is probable some still move into eastern Oregon from Idaho.[41]
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Oregon is home to what is considered the largest single organism in the world, an Armillaria solidipes fungus beneath the Malheur National Forest of eastern Oregon.[10]
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Oregon has three national park sites: Crater Lake National Park in the southern part of the Cascades, John Day Fossil Beds National Monument east of the Cascades, and Lewis and Clark National and State Historical Parks on the north coast.[42][43]
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Most of Oregon has a generally mild climate, though there is significant variation given the variety of landscapes across the state.[44] The state's western region (west of the Cascade Range) has an oceanic climate, populated by dense evergreen mixed forests. Western Oregon's climate is heavily influenced by the Pacific Ocean; the western third of Oregon is very wet in the winter, moderately to very wet during the spring and fall, and dry during the summer. The relative humidity of Western Oregon is high except during summer days, which are semi-dry to semi-humid; Eastern Oregon typically sees low humidity year-round.[45]
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The state's southwestern portion, particularly the Rogue Valley, has a Mediterranean climate with drier and sunnier winters and hotter summers, similar to Northern California.[46]
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Oregon's northeastern portion has a steppe climate, and its high terrain regions have a subarctic climate. Like Western Europe, Oregon, and the Pacific Northwest in general, is considered warm for its latitude, and the state has far milder winters at a given elevation than comparable latitudes elsewhere in North America, such as the Upper Midwest, Ontario, Quebec and New England.[45] However, the state ranks fifth for coolest summer temperatures of any state in the country, after Maine, Idaho, Wyoming, and Alaska.[47]
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The eastern two thirds of Oregon, which largely comprise high desert, have cold, snowy winters and very dry summers. Much of the east is semiarid to arid like the rest of the Great Basin, though the Blue Mountains are wet enough to support extensive forests. Most of Oregon receives significant snowfall, but the Willamette Valley, where 60 percent of the population lives,[48] has considerably milder winters for its latitude and typically sees only light snowfall.[45]
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Oregon's highest recorded temperature is 119 °F (48 °C) at Pendleton on August 10, 1898, and the lowest recorded temperature is −54 °F (−48 °C) at Seneca on February 10, 1933.[49]
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Humans have inhabited the area that is now Oregon for at least 15,000 years. In recorded history, mentions of the land date to as early as the 16th century. During the 18th and 19th centuries, European powers—and later the United States—quarreled over possession of the region until 1846, when the U.S. and Great Britain finalized division of the region. Oregon became a state on February 14, 1859, and as of 2015[update] has more than four million residents.[50]
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While there is considerable evidence that Paleo-Indians inhabited the region, the oldest evidence of habitation in Oregon was found at Fort Rock Cave and the Paisley Caves in Lake County. Archaeologist Luther Cressman dated material from Fort Rock to 13,200 years ago,[51] and there is evidence supporting inhabitants in the region at least 15,000 years ago.[52] By 8000 BC there were settlements throughout the state, with populations concentrated along the lower Columbia River, in the western valleys, and around coastal estuaries.
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During the prehistoric period, the Willamette Valley region was flooded after the collapse of glacial dams from then Lake Missoula, located in what would later become Montana. These massive floods occurred during the last glacial period and filled the valley with 300 to 400 feet (91 to 122 m) of water.[53]
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By the 16th century, Oregon was home to many Native American groups, including the Chinook, Coquille (Ko-Kwell), Bannock, Chasta, Kalapuya, Klamath, Klickitat, Molalla, Nez Perce, Takelma, Killamuk, Neah-kah-nie, Umatilla, and Umpqua.[54][55][56][57]
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The first Europeans to visit Oregon were Spanish explorers led by Juan Rodríguez Cabrillo, who sighted southern Oregon off the Pacific coast in 1543.[58] Sailing Golden Hind in search of the Strait of Anian, the English explorer and privateer Sir Francis Drake briefly sought shelter, anchoring at South Cove, Cape Arago, just south of Coos Bay, before sailing for what is now California.[59][60] Exploration continued routinely in 1774, starting with the expedition of the frigate Santiago by Juan José Pérez Hernández (see Spanish expeditions to the Pacific Northwest), and the coast of Oregon became a valuable trade route to Asia. In 1778, British captain James Cook also explored the coast.[61]
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French Canadians, Scots, Métis and other continental natives (e.g. Iroquois) trappers arrived in the late 18th and early 19th centuries, soon to be followed by Catholic clergy. Some travelled as members of the Lewis and Clark and 1811 Astor expedition. Few stayed permanently such as Étienne Lussier, often referred as the first "European" farmer in the state of Oregon. Evidence of the French Canadian presence can be found in numerous names of French origin such as Malheur Lake, Malheur River, Grande Ronde, Deschutes rivers and the city of La Grande. Furthermore, many of the early pioneers first came out West with the North West Company and the Hudson's Bay Company before heading South of the Columbia for better farmland as the fur trade declined. French Prairie by the Willamette River and French Settlement by the Umpqua River are known as early mixed ancestry settlements.
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The Lewis and Clark Expedition travelled through northern Oregon also in search of the Northwest Passage. They built their winter fort in 1805–06 at Fort Clatsop, near the mouth of the Columbia River, staying at the encampment from December until March.[62]
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British explorer David Thompson also conducted overland exploration. In 1811, while working for the North West Company, Thompson became the first European to navigate the entire Columbia River.[63] Stopping on the way, at the junction of the Snake River, he posted a claim to the region for Great Britain and the North West Company. Upon returning to Montreal, he publicized the abundance of fur-bearing animals in the area.[64]
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Also in 1811, New Yorker John Jacob Astor financed the establishment of Fort Astoria at the mouth of the Columbia River as a western outpost to his Pacific Fur Company;[65] this was the first permanent European settlement in Oregon.
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In the War of 1812, the British gained control of all Pacific Fur Company posts. The Treaty of 1818 established joint British and American occupancy of the region west of the Rocky Mountains to the Pacific Ocean. By the 1820s and 1830s, the Hudson's Bay Company dominated the Pacific Northwest from its Columbia District headquarters at Fort Vancouver (built in 1825 by the district's chief factor, John McLoughlin, across the Columbia from present-day Portland).
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In 1841, the expert trapper and entrepreneur Ewing Young died leaving considerable wealth and no apparent heir, and no system to probate his estate. A meeting followed Young's funeral, at which a probate government was proposed.[66] Doctor Ira Babcock of Jason Lee's Methodist Mission was elected supreme judge.[67] Babcock chaired two meetings in 1842 at Champoeg, (halfway between Lee's mission and Oregon City), to discuss wolves and other animals of contemporary concern. These meetings were precursors to an all-citizen meeting in 1843, which instituted a provisional government headed by an executive committee made up of David Hill, Alanson Beers, and Joseph Gale.[68] This government was the first acting public government of the Oregon Country before annexation by the government of the United States. It was succeeded by a Second Executive Committee, made up of Peter G. Stewart, Osborne Russell, and William J. Bailey, and this committee was itself succeeded by George Abernethy, who was the first and only Governor of Oregon under the provisional government.
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Also in 1841, Sir George Simpson, governor of the Hudson's Bay Company, reversed the Hudson's Bay Company's long-standing policy of discouraging settlement because it interfered with the lucrative fur trade.[69] He directed that some 200 Red River Colony settlers be relocated to HBC farms near Fort Vancouver, (the James Sinclair expedition), in an attempt to hold Columbia District.
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Starting in 1842–43, the Oregon Trail brought many new American settlers to the Oregon Country. For some time, it seemed that Britain and the United States would go to war for a third time in 75 years (see Oregon boundary dispute), but the border was defined peacefully in 1846 by the Oregon Treaty. The border between the United States and British North America was set at the 49th parallel.[70] The Oregon Territory was officially organized on August 13, 1848.[71]
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Settlement increased with the Donation Land Claim Act of 1850 and the forced relocation of the native population to Indian reservations in Oregon.
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In December 1844, Oregon passed its Black Exclusion Law, which prohibited African Americans from entering the territory while simultaneously prohibiting slavery. Slave owners who brought their slaves with them were given three years before they were forced to free them. Any African Americans in the region after the law was passed were forced to leave, and those who did not comply were arrested and beaten. They received no less than twenty and no more than thirty-nine stripes across their bare back if they still did not leave. This process could be repeated every six months.[72] Slavery played a major part in Oregon's history and even influenced its path to statehood. The territory's request for statehood was delayed several times, as members of Congress argued among themselves whether the territory should be admitted as a "free" or "slave" state. Eventually politicians from the south agreed to allow Oregon to enter as a "free" state, in exchange for opening slavery to the southwest United States.[73]
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Oregon was admitted to the Union on February 14, 1859. Founded as a refuge from disputes over slavery, Oregon had a "whites only" clause in its original state Constitution.[74][75] At the outbreak of the American Civil War, regular U.S. troops were withdrawn and sent east. Volunteer cavalry recruited in California were sent north to Oregon to keep peace and protect the populace. The First Oregon Cavalry served until June 1865.
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Beginning in the 1880s, the growth of railroads expanded the state's lumber, wheat, and other agricultural markets, and the rapid growth of its cities.[76] Due to the abundance of timber and waterway access via the Willamette River, Portland became a major force in the lumber industry of the Pacific Northwest, and quickly became the state's largest city. It would earn the nickname "Stumptown",[77] and would later become recognized as one of the most dangerous port cities in the United States due to racketeering and illegal activities at the turn of the 20th century.[78] In 1902, Oregon introduced direct legislation by the state's citizens through initiatives and referenda, known as the Oregon System.[79]
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On May 5, 1945, six civilians were killed by a Japanese balloon bomb that exploded on Gearhart Mountain near Bly.[80][81] They remained the only people on American soil whose deaths were attributed to an enemy balloon bomb explosion during World War II. The bombing site is now located in the Mitchell Recreation Area.
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Industrial expansion began in earnest following the 1933–37 construction of the Bonneville Dam on the Columbia River. Hydroelectric power, food, and lumber provided by Oregon helped fuel the development of the West, although the periodic fluctuations in the U.S. building industry have hurt the state's economy on multiple occasions. Portland in particular experienced a population boom between 1900 and 1930, tripling in size; the arrival of World War II also provided the northwest region of the state with an industrial boom, where Liberty ships and aircraft carriers were constructed.[82]
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During the 1970s, the Pacific Northwest was particularly affected by the 1973 oil crisis, with Oregon suffering a substantial shortage.[83]
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In 1994, Oregon became the first U.S. state to legalize physician-assisted suicide through the Oregon Death with Dignity Act. A measure to legalize recreational use of marijuana in Oregon was approved on November 4, 2014, making Oregon only the second state at the time to have legalized gay marriage, physician-assisted suicide, and recreational marijuana.[84]
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Oregon's population is largely concentrated in the Willamette Valley, which stretches from Eugene in the south (home of the University of Oregon) through Corvallis (home of Oregon State University) and Salem (the capital) to Portland (Oregon's largest city).[85]
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Astoria, at the mouth of the Columbia River, was the first permanent English-speaking settlement west of the Rockies in what is now the United States. Oregon City, at the end of the Oregon Trail, was the Oregon Territory's first incorporated city, and was its first capital from 1848 until 1852, when the capital was moved to Salem. Bend, near the geographic center of the state, is one of the ten fastest-growing metropolitan areas in the United States.[86][better source needed] In southern Oregon, Medford is a rapidly growing metro area and is home to the Rogue Valley International-Medford Airport, the state's third-busiest airport. To the south, near the California border, is the city of Ashland. Eastern Oregon is very sparsely populated, but is home to Hermiston, which with a population of 18,000 is the largest and fastest growing city in the region.[87]
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A writer in the Oregon Country book A Pacific Republic, written in 1839, predicted the territory was to become an independent republic. Four years later, in 1843, settlers of the Willamette Valley voted in majority for a republic government.[90] The Oregon Country functioned in this way until August 13, 1848, when Oregon was annexed by the United States and a territorial government was established. Oregon maintained a territorial government until February 14, 1859, when it was granted statehood.[91]
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Oregon state government has a separation of powers similar to the federal government. It has three branches:
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Governors in Oregon serve four-year terms and are limited to two consecutive terms, but an unlimited number of total terms. Oregon has no lieutenant governor; in the event that the office of governor is vacated, Article V, Section 8a of the Oregon Constitution specifies that the Secretary of State is first in line for succession.[92] The other statewide officers are Treasurer, Attorney General, Superintendent, and Labor Commissioner. The biennial Oregon Legislative Assembly consists of a thirty-member Senate and a sixty-member House. The state supreme court has seven elected justices, currently including the only two openly gay state supreme court justices in the nation. They choose one of their own to serve a six-year term as Chief Justice.
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The debate over whether to move to annual sessions is a long-standing battle in Oregon politics, but the voters have resisted the move from citizen legislators to professional lawmakers. Because Oregon's state budget is written in two-year increments and, there being no sales tax, state revenue is based largely on income taxes, it is often significantly over- or under-budget. Recent legislatures have had to be called into special session repeatedly to address revenue shortfalls resulting from economic downturns, bringing to a head the need for more frequent legislative sessions. Oregon Initiative 71, passed in 2010, mandates the legislature to begin meeting every year, for 160 days in odd-numbered years, and 35 days in even-numbered years.
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Oregonians have voted for the Democratic presidential candidate in every election since 1988. In 2004 and 2006, Democrats won control of the state Senate, and then the House. Since the late 1990s, Oregon has been represented by four Democrats and one Republican in the U.S. House of Representatives. Since 2009, the state has had two Democratic U.S. senators, Ron Wyden and Jeff Merkley. Oregon voters have elected Democratic governors in every election since 1986, most recently electing Kate Brown over Republican Bud Pierce in a 2016 special election for a two-year term, and re-electing her for a full four-year term over Republican Knute Buehler in 2018.
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The base of Democratic support is largely concentrated in the urban centers of the Willamette Valley. The eastern two-thirds of the state beyond the Cascade Mountains typically votes Republican; in 2000 and 2004, George W. Bush carried every county east of the Cascades. However, the region's sparse population means the more populous counties in the Willamette Valley usually outweigh the eastern counties in statewide elections.
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In the 2002 general election, Oregon voters approved a ballot measure to increase the state minimum wage automatically each year according to inflationary changes, which are measured by the consumer price index (CPI).[93] In the 2004 general election, Oregon voters passed ballot measures banning same-sex marriage[94] and restricting land use regulation.[95] In the 2006 general election, voters restricted the use of eminent domain and extended the state's discount prescription drug coverage.[96]
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Like all U.S. states, Oregon is represented by two senators. Since the 1980 census, Oregon has had five congressional districts. After Oregon was admitted to the Union, it began with a single member in the House of Representatives (La Fayette Grover, who served in the 35th United States Congress for less than a month). Congressional apportionment increased the size of the delegation following the censuses of 1890, 1910, 1940, and 1980. A detailed list of the past and present Congressional delegations from Oregon is available.
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The United States District Court for the District of Oregon hears federal cases in the state. The court has courthouses in Portland, Eugene, Medford, and Pendleton. Also in Portland is the federal bankruptcy court, with a second branch in Eugene.[97] Oregon (among other western states and territories) is in the 9th Court of Appeals. One of the court's meeting places is at the Pioneer Courthouse in downtown Portland, a National Historic Landmark built in 1869.
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Political opinions in Oregon are geographically split by the Cascade Range, with western Oregon being more liberal and Eastern Oregon being conservative.[99] In a 2008 analysis of the 2004 presidential election, a political analyst found that according to the application of a Likert scale, Oregon boasted both the most liberal Kerry voters and the most conservative Bush voters, making it the most politically polarized state in the country.[100]
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While Republicans typically win more counties by running up huge margins in the east, the Democratic tilt of the more populated west is usually enough to swing the entire state Democratic. In 2008, for instance, Republican Senate incumbent Gordon H. Smith lost his bid for a third term, even though he carried all but eight counties. His Democratic challenger, Jeff Merkley, won Multnomah County by 142,000 votes, more than double the overall margin of victory.
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During Oregon's history, it has adopted many electoral reforms proposed during the Progressive Era, through the efforts of William S. U'Ren and his Direct Legislation League. Under his leadership, the state overwhelmingly approved a ballot measure in 1902 that created the initiative and referendum for citizens to introduce or approve proposed laws or amendments to the state constitution directly, making Oregon the first state to adopt such a system. Today, roughly half of U.S. states do so.[101]
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In following years, the primary election to select party candidates was adopted in 1904, and in 1908 the Oregon Constitution was amended to include recall of public officials. More recent amendments include the nation's first doctor-assisted suicide law,[102] called the Death with Dignity Act (which was challenged, unsuccessfully, in 2005 by the Bush administration in a case heard by the U.S. Supreme Court), legalization of medical cannabis, and among the nation's strongest anti-urban sprawl and pro-environment laws.[citation needed] More recently, 2004's Measure 37 reflects a backlash against such land-use laws. However, a further ballot measure in 2007, Measure 49, curtailed many of the provisions of 37.
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Of the measures placed on the ballot since 1902, the people have passed 99 of the 288 initiatives and 25 of the 61 referendums on the ballot, though not all of them survived challenges in courts (see Pierce v. Society of Sisters, for an example). During the same period, the legislature has referred 363 measures to the people, of which 206 have passed.
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Oregon pioneered the American use of postal voting, beginning with experimentation approved by the Oregon Legislative Assembly in 1981 and culminating with a 1998 ballot measure mandating that all counties conduct elections by mail. It remains one of just two states, the other being Washington, where voting by mail is the only method of voting.
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In 1994, Oregon adopted the Oregon Health Plan, which made health care available to most of its citizens without private health insurance.[103]
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In the U.S. Electoral College, Oregon casts seven votes. Oregon has supported Democratic candidates in the last eight elections. Democratic incumbent Barack Obama won the state by a margin of twelve percentage points, with over 54% of the popular vote in 2012. In the 2016 election, Hillary Clinton won Oregon by 11 percentage points.[104]
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As of 2015[update], Oregon ranks as the 17th highest in median household income at $60,834.[4] The gross domestic product (GDP) of Oregon in 2013 was $219.6 billion, a 2.7% increase from 2012; Oregon is the 25th wealthiest state by GDP. In 2003, Oregon was 28th in the U.S. by GDP. The state's per capita personal income (PCPI) in 2013 was $39,848, a 1.5% increase from 2012. Oregon ranks 33rd in the U.S. by PCPI, compared to 31st in 2003. The national PCPI in 2013 was $44,765.[106]
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Oregon's unemployment rate was 5.5% in September 2016,[107] while the U.S. unemployment rate was 5.0% that month.[108] Oregon has the third largest amount of food stamp users in the nation (21% of the population).[109]
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Oregon's diverse landscapes provide ideal environments for various types of farming. Land in the Willamette Valley owes its fertility to the Missoula Floods, which deposited lake sediment from Glacial Lake Missoula in western Montana onto the valley floor.[110] In 2016, the Willamette Valley region produced over 100 million pounds (45 kt) of blueberries.[111]
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Oregon is also one of four major world hazelnut growing regions, and produces 95% of the domestic hazelnuts in the United States. While the history of the wine production in Oregon can be traced to before Prohibition, it became a significant industry beginning in the 1970s. In 2005, Oregon ranked third among U.S. states with 303 wineries.[112] Due to regional similarities in climate and soil, the grapes planted in Oregon are often the same varieties found in the French regions of Alsace and Burgundy. In 2014, 71 wineries opened in the state. The total is currently 676, which represents growth of 12% over 2013.[113]
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In the southern Oregon coast, commercially cultivated cranberries account for about 7 percent of U.S. production, and the cranberry ranks 23rd among Oregon's top 50 agricultural commodities. Cranberry cultivation in Oregon uses about 27,000 acres (110 square kilometers) in southern Coos and northern Curry counties, centered around the coastal city of Bandon. In the northeastern region of the state, particularly around Pendleton, both irrigated and dry land wheat is grown.[114] Oregon farmers and ranchers also produce cattle, sheep, dairy products, eggs and poultry.
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Vast forests have historically made Oregon one of the nation's major timber-producing and logging states, but forest fires (such as the Tillamook Burn), over-harvesting, and lawsuits over proper management of the extensive federal forest holdings have reduced the timber produced. Between 1989 and 2011, the amount of timber harvested from federal lands in Oregon dropped about 90%, although harvest levels on private land have remained relatively constant.[115]
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Even the shift in recent years towards finished goods such as paper and building materials has not slowed the decline of the timber industry in the state. The effects of this decline have included Weyerhaeuser's acquisition of Portland-based Willamette Industries in January 2002, the relocation of Louisiana-Pacific's corporate headquarters from Portland to Nashville, and the decline of former lumber company towns such as Gilchrist. Despite these changes, Oregon still leads the United States in softwood lumber production; in 2011, 4,134 million board feet (9,760,000 m3) was produced in Oregon, compared with 3,685 million board feet (8,700,000 m3) in Washington, 1,914 million board feet (4,520,000 m3) in Georgia, and 1,708 million board feet (4,030,000 m3) in Mississippi.[116] The slowing of the timber and lumber industry has caused high unemployment rates in rural areas.[117]
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Oregon has one of the largest salmon-fishing industries in the world, although ocean fisheries have reduced the river fisheries in recent years.[118] Because of the abundance of waterways in the state, it is also a major producer of hydroelectric energy.[119]
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Tourism is also a strong industry in the state. Tourism is centered on the state's natural features – mountains, forests, waterfalls, rivers, beaches and lakes, including Crater Lake National Park, Multnomah Falls, the Painted Hills, the Deschutes River, and the Oregon Caves. Mount Hood and Mount Bachelor also draw visitors year-round for skiing and other snow activities.[120]
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Portland is home to the Oregon Museum of Science and Industry, the Portland Art Museum, and the Oregon Zoo, which is the oldest zoo west of the Mississippi River.[121] The International Rose Test Garden is another prominent attraction in the city. Portland has also been named the best city in the world for street food by several publications, including the U.S. News & World Report and CNN.[122][123] Oregon is home to many breweries, and Portland has the largest number of breweries of any city in the world.[124]
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The state's coastal region produces significant tourism as well.[125] The Oregon Coast Aquarium comprises 23 acres (9.3 ha) along Yaquina Bay in Newport, and was also home to Keiko the orca whale.[126] It has been noted as one of the top ten aquariums in North America.[127] Fort Clatsop in Warrenton features a replica of Lewis and Clark's encampment at the mouth of the Columbia River in 1805. The Sea Lion Caves in Florence are the largest system of sea caverns in the United States, and also attract many visitors.[128]
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In Southern Oregon, the Oregon Shakespeare Festival, held in Ashland, is also a tourist draw, as is the Oregon Vortex and the Wolf Creek Inn State Heritage Site, a historic inn where Jack London wrote his 1913 novel Valley of the Moon.[129]
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Oregon has also historically been a popular region for film shoots due to its diverse landscapes, as well as its proximity to Hollywood (see List of films shot in Oregon).[130] Movies filmed in Oregon include: Animal House, Free Willy, The General, The Goonies, Kindergarten Cop, One Flew Over the Cuckoo's Nest, and Stand By Me. Oregon native Matt Groening, creator of The Simpsons, has incorporated many references from his hometown of Portland into the TV series.[131] Additionally, several television shows have been filmed throughout the state including Portlandia, Grimm, Bates Motel, and Leverage.[132] The Oregon Film Museum is located in the old Clatsop County Jail in Astoria.
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High technology industries located in Silicon Forest have been a major employer since the 1970s. Tektronix was the largest private employer in Oregon until the late 1980s. Intel's creation and expansion of several facilities in eastern Washington County continued the growth that Tektronix had started. Intel, the state's largest for-profit private employer,[133][134] operates four large facilities, with Ronler Acres, Jones Farm and Hawthorn Farm all located in Hillsboro.[135]
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The spinoffs and startups that were produced by these two companies led to establishment of the so-called Silicon Forest. The recession and dot-com bust of 2001 hit the region hard; many high technology employers reduced the number of their employees or went out of business. Open Source Development Labs made news in 2004 when they hired Linus Torvalds, developer of the Linux kernel. In 2010, biotechnology giant Genentech opened a $400 million facility in Hillsboro to expand its production capabilities.[136] Oregon is home to several large datacenters that take advantage of cheap power and a climate conducive to reducing cooling costs. Google operates a large datacenter in The Dalles, and Facebook built a large datacenter near Prineville in 2010. Amazon opened a datacenter near Boardman in 2011, and a fulfillment center in Troutdale in 2018.[137][138]
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Oregon is also the home of large corporations in other industries. The world headquarters of Nike are located near Beaverton. Medford is home to Harry and David, which sells gift items under several brands. Medford is also home to the national headquarters of Lithia Motors. Portland is home to one of the West's largest trade book publishing houses, Graphic Arts Center Publishing. Oregon is also home to Mentor Graphics Corporation, a world leader in electronic design automation located in Wilsonville and employs roughly 4,500 people worldwide.
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Adidas Corporations American Headquarters is located in Portland and employs roughly 900 full-time workers at its Portland campus.[140] Nike, located in Beaverton, employs roughly 5,000 full-time employees at its 200-acre (81 ha) campus. Nike's Beaverton campus is continuously ranked as a top employer in the Portland area-along with competitor Adidas.[141] Intel Corporation employs 18,600 in Oregon[134] with the majority of these employees located at the company's Hillsboro campus located about 30 minutes west of Portland. Intel has been a top employer in Oregon since 1974.[142]
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The U.S. Federal Government and Providence Health systems are respective contenders for top employers in Oregon with roughly 12,000 federal workers and 14,000 Providence Health workers.
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In 2015, a total of seven companies headquartered in Oregon landed in the Fortune 1000: Nike, at 106; Precision Castparts Corp. at 302; Lithia Motors at 482; StanCorp Financial Group at 804; Schnitzer Steel Industries at 853; The Greenbrier Companies at 948; and Columbia Sportswear at 982.[143]
|
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Oregon's biennial state budget, $2.6 billion in 2017, comprises General Funds, Federal Funds, Lottery Funds, and Other Funds.[144]
|
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Oregon is one of only five states that have no sales tax.[145] Oregon voters have been resolute in their opposition to a sales tax, voting proposals down each of the nine times they have been presented.[146] The last vote, for 1993's Measure 1, was defeated by a 75–25% margin.[147]
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The state also has a minimum corporate tax of only $150 a year,[148] amounting to 5.6% of the General Fund in the 2005–07 biennium; data about which businesses pay the minimum is not available to the public.[149][better source needed] As a result, the state relies on property and income taxes for its revenue. Oregon has the fifth highest personal income tax in the nation. According to the U.S. Census Bureau, Oregon ranked 41st out of the 50 states in taxes per capita in 2005 with an average amount paid of 1,791.45.[150]
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A few local governments levy sales taxes on services: the city of Ashland, for example, collects a 5% sales tax on prepared food.[151]
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The City of Portland imposes an Arts Education and Access Income Tax on residents over 18—a flat tax of $35 collected from individuals earning $1,000 or more per year and residing in a household with an annual income exceeding the federal poverty level. The tax funds Portland school teachers, and art focused non-profit organizations in Portland.[152]
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The State of Oregon also allows transit district to levy an income tax on employers and the self-employed. The State currently collects the tax for TriMet and the Lane Transit District.[153][154]
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Oregon is one of six states with a revenue limit.[155] The "kicker law" stipulates that when income tax collections exceed state economists' estimates by two percent or more, any excess must be returned to taxpayers.[156] Since the enactment of the law in 1979, refunds have been issued for seven of the eleven biennia.[157] In 2000, Ballot Measure 86 converted the "kicker" law from statute to the Oregon Constitution, and changed some of its provisions.
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Federal payments to county governments that were granted to replace timber revenue when logging in National Forests was restricted in the 1990s, have been under threat of suspension for several years. This issue dominates the future revenue of rural counties, which have come to rely on the payments in providing essential services.[158]
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55% of state revenues are spent on public education, 23% on human services (child protective services, Medicaid, and senior services), 17% on public safety, and 5% on other services.[159]
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For health insurance, as of 2018 Cambia Health Solutions has the highest market share at 21%, followed by Providence Health.[160] In the Portland region, Kaiser Permanente leads.[160] Providence and Kaiser are vertically integrated delivery systems which operate hospitals and offer insurance plans.[161] Aside from Providence and Kaiser, hospital systems which are primarily Oregon-based include Legacy Health mostly covering Portland, Samaritan Health Services with five hospitals in various areas across the state, and Tuality Healthcare in the western Portland metropolitan area. In Southern Oregon, Asante runs several hospitals, including Rogue Regional Medical Center. Some hospitals are operated by multi-state organizations such as PeaceHealth and CommonSpirit Health. Some hospitals such Salem Hospital operate independently of larger systems.
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Oregon Health & Science University is a Portland-based medical school which operates two hospitals and clinics.
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The Oregon Health Plan is the state's Medicaid managed care plan, and it is known for innovations.[162] The Portland area is a mature managed care and two-thirds of Medicare enrollees are in Medicare Advantage plans.[162]
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The United States Census Bureau estimates that the population of Oregon was 4,217,737 on July 1, 2019, a 10.09% increase over the 2010 census.[50]
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Oregon was the nation's "Top Moving Destination" in 2014, with two families moving into the state for every one moving out (66.4% to 33.6%).[166] Oregon was also the top moving destination in 2013,[167] and the second-most popular destination in 2010 through 2012.[168][169]
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As of the 2010 census, the population of Oregon was 3,831,074. The gender makeup of the state was 49.5% male and 50.5% female. 22.6% of the population were under the age of 18; 63.5% were between the ages of 18 and 64; and 12.5% were 65 years of age or older.[170]
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The table below shows the racial composition of Oregon's population as of 2016.
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According to the 2016 American Community Survey, 12.4% of Oregon's population were of Hispanic or Latino origin (of any race): Mexican (10.4%), Puerto Rican (0.3%), Cuban (0.1%), and other Hispanic or Latino origin (1.5%).[171] The five largest ancestry groups were: German (19.1%), Irish (11.7%), English (11.3%), American (5.3%), and Norwegian (3.8%).[174]
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The state's most populous ethnic group, non-Hispanic white, has declined from 95.8% in 1970 to 77.8% in 2012.[175][176]
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As of 2011[update], 38.7% of Oregon's children under one year of age belonged to minority groups, meaning they had at least one parent who was not a non-Hispanic white.[177] Of the state's total population, 22.6% was under the age 18, and 77.4% were 18 or older.
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The center of population of Oregon is located in Linn County, in the city of Lyons.[178] Around 60% of Oregon's population resides within the Portland metropolitan area.[179]
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As of 2009[update], Oregon's population comprised 361,393 foreign-born residents.[180] Of the foreign-born residents, the three largest groups are originally from countries in: Latin America (47.8%), Asia (27.4%), and Europe (16.5%).[180]
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Oregon has frequently been cited by statistical agencies for having a smaller percentage of religious communities than other U.S. states.[182][183] According to a 2009 Gallup poll, Oregon was paired with Vermont as the two "least religious" states in the United States.[184]
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In the same 2009 Gallup poll, 69% of Oregonians identified themselves as being Christian.[185] The largest Christian denominations in Oregon by number of adherents in 2010 were the Roman Catholic Church with 398,738; The Church of Jesus Christ of Latter-day Saints with 147,965; and the Assemblies of God with 45,492.[186] Oregon also contains the largest community of Russian Old Believers to be found in the United States.[187] Judaism is the largest non-Christian religion in Oregon with more than 50,000 adherents, 47,000 of whom live in the Portland area.[188][189] Recently, new kosher food and Jewish educational offerings have led to a rapid increase in Portland's Orthodox Jewish population.[190] The Northwest Tibetan Cultural Association is headquartered in Portland. There are an estimated 6,000 to 10,000 Muslims in Oregon, most of whom live in and around Portland.[191]
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Most of the remainder of the population had no religious affiliation; the 2008 American Religious Identification Survey (ARIS) placed Oregon as tied with Nevada in fifth place of U.S. states having the highest percentage of residents identifying themselves as "non-religious", at 24 percent.[192][193] Secular organizations include the Center for Inquiry (CFI), the Humanists of Greater Portland (HGP), and the United States Atheists (USA).
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During much of the 1990s, a group of conservative Christians formed the Oregon Citizens Alliance, and unsuccessfully tried to pass legislation to prevent "gay sensitivity training" in public schools and legal benefits for homosexual couples.[194]
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Projections from the U.S. Census Bureau show Oregon's population increasing to 4,833,918 by 2030, an increase of 41.3% compared to the state's population of 3,421,399 in 2000.[201] The state's own projections forecast a total population of 5,425,408 in 2040.[202]
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In the 2013–2014 school year, the state had 567,000 students in public schools.[203] There were 197 public school districts, served by 19 education service districts.[203]
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In 2016, the largest school districts in the state were:[204] Portland Public Schools, comprising 47,323 students; Salem-Keizer School District, comprising 40,565 students; Beaverton School District, comprising 39,625 students; Hillsboro School District, comprising 21,118 students; and North Clackamas School District, comprising 17,053 students.
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Approximately 90.5% of Oregon high school students graduate, improving on the national average of 88.3% as measured from the 2010 United States Census.[205]
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Especially since the 1990 passage of Measure 5, which set limits on property tax levels, Oregon has struggled to fund higher education. Since then, Oregon has cut its higher education budget and now ranks 46th in the country in state spending per student. However, 2007 legislation funded the university system far beyond the governor's requested budget though still capping tuition increases at 3% per year.[206] Oregon supports a total of seven public universities and one affiliate. It is home to three public research universities: The University of Oregon (UO) in Eugene and Oregon State University (OSU) in Corvallis, both classified as research universities with very high research activity, and Portland State University which is classified as a research university with high research activity.[207]
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UO is the state's highest nationally-ranked and most selective[208] public university by U.S. News & World Report and Forbes.[209] OSU is the state's only land-grant university, has the state's largest enrollment for fall 2014,[210] and is the state's highest ranking university according to Academic Ranking of World Universities, Washington Monthly, and QS World University Rankings.[211] OSU receives more annual funding for research than all other public higher education institutions in Oregon combined.[212] The state's urban Portland State University has Oregon's second largest enrollment.
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The state has three regional universities: Western Oregon University in Monmouth, Southern Oregon University in Ashland, and Eastern Oregon University in La Grande. The Oregon Institute of Technology has its campus in Klamath Falls. The quasi-public Oregon Health & Science University (OHSU) includes medical, dental, and nursing schools, and graduate programs in biomedical sciences in Portland and a science and engineering school in Hillsboro. The state also supports 17 community colleges.
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Oregon is home to a wide variety of private colleges, the majority of which are located in the Portland area. The University of Portland and Marylhurst University are both Catholic universities located in or near Portland, affiliated with the Congregation of Holy Cross, and the Sisters of the Holy Names of Jesus and Mary, respectively. Reed College, a rigorous liberal arts college in Portland, was ranked by Forbes as the 52nd best college in the country in 2015.[213]
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Other private institutions in Portland include Lewis & Clark College; Multnomah University; Portland Bible College; Warner Pacific College; Cascade College; the National University of Natural Medicine; and Western Seminary, a theological graduate school. Pacific University is in the Portland suburb of Forest Grove. There are also private colleges further south in the Willamette Valley. McMinnville is home to Linfield College, while nearby Newberg is home to George Fox University. Salem is home to two private schools: Willamette University (the state's oldest, established during the provisional period) and Corban University. Also located near Salem is Mount Angel Seminary, one of America's largest Roman Catholic seminaries. The state's second medical school, the College of Osteopathic Medicine of the Pacific, Northwest, is located in Lebanon. Eugene is home to three private colleges: Northwest Christian University, New Hope Christian College, and Gutenberg College.
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Oregon is home to three major professional sports teams: the Portland Trail Blazers of the NBA, the Portland Thorns of the NWSL and the Portland Timbers of MLS.[214]
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Until 2011, the only major professional sports team in Oregon was the Portland Trail Blazers of the National Basketball Association. From the 1970s to the 1990s, the Blazers were one of the most successful teams in the NBA in terms of both win-loss record and attendance.[215] In the early 21st century, the team's popularity declined due to personnel and financial issues, but revived after the departure of controversial players and the acquisition of new players such as Brandon Roy and LaMarcus Aldridge, and still later Damian Lillard.[216][217] The Blazers play in the Moda Center in Portland's Lloyd District, which also is home to the Portland Winterhawks of the junior Western Hockey League.[218]
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The Portland Timbers play at Providence Park, just west of downtown Portland. The Timbers have a strong following, with the team regularly selling out its games.[219] The Timbers repurposed the formerly multi-use stadium into a soccer-specific stadium in fall 2010, increasing the seating in the process.[220] The Timbers operate Portland Thorns FC, a women's soccer team that has played in the National Women's Soccer League since the league's first season in 2013. The Thorns, who also play at Providence Park, have won two league championships, in the inaugural 2013 season and also in 2017, and have been by far the NWSL's attendance leader in each of the league's seasons.
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Eugene, Salem and Hillsboro have minor-league baseball teams: the Eugene Emeralds, the Salem-Keizer Volcanoes, and the Hillsboro Hops all play in the Single-A Northwest League.[221] Portland has had minor-league baseball teams in the past, including the Portland Beavers and Portland Rockies, who played most recently at Providence Park when it was known as PGE Park.
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The Oregon State Beavers and the University of Oregon Ducks football teams of the Pac-12 Conference meet annually in the Civil War. Both schools have had recent success in other sports as well: Oregon State won back-to-back college baseball championships in 2006 and 2007,[222] winning a third in 2018;[223] and the University of Oregon won back-to-back NCAA men's cross country championships in 2007 and 2008.[224]
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Government
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Tourism and recreation
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History and culture
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Maps and geology
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Coordinates: 44°00′N 120°30′W / 44°N 120.5°W / 44; -120.5
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The ear is the organ of hearing and, in mammals, balance. In mammals, the ear is usually described as having three parts—the outer ear, the middle ear and the inner ear. The outer ear consists of the pinna and the ear canal. Since the outer ear is the only visible portion of the ear in most animals, the word "ear" often refers to the external part alone.[1] The middle ear includes the tympanic cavity and the three ossicles. The inner ear sits in the bony labyrinth, and contains structures which are key to several senses: the semicircular canals, which enable balance and eye tracking when moving; the utricle and saccule, which enable balance when stationary; and the cochlea, which enables hearing. The ears of vertebrates are placed somewhat symmetrically on either side of the head, an arrangement that aids sound localisation.
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The ear develops from the first pharyngeal pouch and six small swellings that develop in the early embryo called otic placodes, which are derived from ectoderm.
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The ear may be affected by disease, including infection and traumatic damage. Diseases of the ear may lead to hearing loss, tinnitus and balance disorders such as vertigo, although many of these conditions may also be affected by damage to the brain or neural pathways leading from the ear.
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The ear has been adorned by earrings and other jewelry in numerous cultures for thousands of years, and has been subjected to surgical and cosmetic alterations.
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The human ear consists of three parts—the outer ear, middle ear and inner ear.[2] The ear canal of the outer ear is separated from the air-filled tympanic cavity of the middle ear by the eardrum. The middle ear contains the three small bones—the ossicles—involved in the transmission of sound, and is connected to the throat at the nasopharynx, via the pharyngeal opening of the Eustachian tube. The inner ear contains the otolith organs—the utricle and saccule—and the semicircular canals belonging to the vestibular system, as well as the cochlea of the auditory system.[2]
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The outer ear is the external portion of the ear and includes the fleshy visible pinna (also called the auricle), the ear canal, and the outer layer of the eardrum (also called the tympanic membrane).[2][3]
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The pinna consists of the curving outer rim called the helix, the inner curved rim called the antihelix, and opens into the ear canal. The tragus protrudes and partially obscures the ear canal, as does the facing antitragus. The hollow region in front of the ear canal is called the concha. The ear canal stretches for about 1 inch (2.5 cm). The first part of the canal is surrounded by cartilage, while the second part near the eardrum is surrounded by bone. This bony part is known as the auditory bulla and is formed by the tympanic part of the temporal bone. The skin surrounding the ear canal contains ceruminous and sebaceous glands that produce protective ear wax. The ear canal ends at the external surface of the eardrum.[3]
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Two sets of muscles are associated with the outer ear: the intrinsic and extrinsic muscles. In some mammals, these muscles can adjust the direction of the pinna.[3] In humans, these muscles have little or no effect.[4] The ear muscles are supplied by the facial nerve, which also supplies sensation to the skin of the ear itself, as well as to the external ear cavity. The great auricular nerve, auricular nerve, auriculotemporal nerve, and lesser and greater occipital nerves of the cervical plexus all supply sensation to parts of the outer ear and the surrounding skin.[3]
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The pinna consists of a single piece of elastic cartilage with a complicated relief on its inner surface and a fairly smooth configuration on its posterior surface. A tubercle, known as Darwin's tubercle, is sometimes present, lying in the descending part of the helix and corresponding to the ear-tip of mammals. The earlobe consists of areola and adipose tissue.[5] The symmetrical arrangement of the two ears allows for the localisation of sound. The brain accomplishes this by comparing arrival-times and intensities from each ear, in circuits located in the superior olivary complex and the trapezoid bodies which are connected via pathways to both ears.[6]
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The middle ear lies between the outer ear and the inner ear. It consists of an air-filled cavity called the tympanic cavity and includes the three ossicles and their attaching ligaments; the auditory tube; and the round and oval windows. The ossicles are three small bones that function together to receive, amplify, and transmit the sound from the eardrum to the inner ear. The ossicles are the malleus (hammer), incus (anvil), and the stapes (stirrup). The stapes is the smallest named bone in the body. The middle ear also connects to the upper throat at the nasopharynx via the pharyngeal opening of the Eustachian tube.[3][7]
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The three ossicles transmit sound from the outer ear to the inner ear. The malleus receives vibrations from sound pressure on the eardrum, where it is connected at its longest part (the manubrium or handle) by a ligament. It transmits vibrations to the incus, which in turn transmits the vibrations to the small stapes bone. The wide base of the stapes rests on the oval window. As the stapes vibrates, vibrations are transmitted through the oval window, causing movement of fluid within the cochlea.[3]
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The round window allows for the fluid within the inner ear to move. As the stapes pushes the secondary tympanic membrane, fluid in the inner ear moves and pushes the membrane of the round window out by a corresponding amount into the middle ear. The ossicles help amplify sound waves by nearly 15–20 times.[2]
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The inner ear sits within the temporal bone in a complex cavity called the bony labyrinth. A central area known as the vestibule contains two small fluid-filled recesses, the utricle and saccule. These connect to the semicircular canals and the cochlea. There are three semicircular canals angled at right angles to each other which are responsible for dynamic balance. The cochlea is a spiral shell-shaped organ responsible for the sense of hearing. These structures together create the membranous labyrinth.[8]
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The bony labyrinth refers to the bony compartment which contains the membranous labyrinth, contained within the temporal bone. The inner ear structurally begins at the oval window, which receives vibrations from the incus of the middle ear. Vibrations are transmitted into the inner ear into a fluid called endolymph, which fills the membranous labyrinth. The endolymph is situated in two vestibules, the utricle and saccule, and eventually transmits to the cochlea, a spiral-shaped structure. The cochlea consists of three fluid-filled spaces: the vestibular duct, the cochlear duct, and the tympanic duct.[3] Hair cells responsible for transduction—changing mechanical changes into electrical stimuli are present in the organ of Corti in the cochlea.[8]
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The blood supply of the ear differs according to each part of the ear.
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The outer ear is supplied by a number of arteries. The posterior auricular artery provides the majority of the blood supply. The anterior auricular arteries provide some supply to the outer rim of the ear and scalp behind it. The posterior auricular artery is a direct branch of the external carotid artery, and the anterior auricular arteries are branches from the superficial temporal artery. The occipital artery also plays a role.[8]
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The middle ear is supplied by the mastoid branch of either the occipital or posterior auricular arteries and the deep auricular artery, a branch of the maxillary artery. Other arteries which are present but play a smaller role include branches of the middle meningeal artery, ascending pharyngeal artery, internal carotid artery, and the artery of the pterygoid canal.[8]
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The inner ear is supplied by the anterior tympanic branch of the maxillary artery; the stylomastoid branch of the posterior auricular artery; the petrosal branch of middle meningeal artery; and the labyrinthine artery, arising from either the anterior inferior cerebellar artery or the basilar artery.[8]
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Sound waves travel through the outer ear, are modulated by the middle ear, and are transmitted to the vestibulocochlear nerve in the inner ear. This nerve transmits information to the temporal lobe of the brain, where it is registered as sound.
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Sound that travels through the outer ear impacts on the eardrum, and causes it to vibrate. The three ossicles bones transmit this sound to a second window (the oval window) which protects the fluid-filled inner ear. In detail, the pinna of the outer ear helps to focus a sound, which impacts on the eardrum. The malleus rests on the membrane, and receives the vibration. This vibration is transmitted along the incus and stapes to the oval window. Two small muscles, the tensor tympani and stapedius, also help modulate noise. The two muscles reflexively contract to dampen excessive vibrations. Vibration of the oval window causes vibration of the endolymph within the vestibule and the cochlea.[9]
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The inner ear houses the apparatus necessary to change the vibrations transmitted from the outside world via the middle ear into signals passed along the vestibulocochlear nerve to the brain. The hollow channels of the inner ear are filled with liquid, and contain a sensory epithelium that is studded with hair cells. The microscopic "hairs" of these cells are structural protein filaments that project out into the fluid. The hair cells are mechanoreceptors that release a chemical neurotransmitter when stimulated. Sound waves moving through fluid flows against the receptor cells of the organ of Corti. The fluid pushes the filaments of individual cells; movement of the filaments causes receptor cells to become open to receive the potassium-rich endolymph. This causes the cell to depolarise, and creates an action potential that is transmitted along the spiral ganglion, which sends information through the auditory portion of the vestibulocochlear nerve to the temporal lobe of the brain.[9]
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The human ear can generally hear sounds with frequencies between 20 Hz and 20 kHz (the audio range). Sounds outside this range are considered infrasound (below 20 Hz)[10] or ultrasound (above 20 kHz)[11] Although hearing requires an intact and functioning auditory portion of the central nervous system as well as a working ear, human deafness (extreme insensitivity to sound) most commonly occurs because of abnormalities of the inner ear, rather than in the nerves or tracts of the central auditory system.
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Providing balance, when moving or stationary, is also a central function of the ear. The ear facilitates two types of balance: static balance, which allows a person to feel the effects of gravity, and dynamic balance, which allows a person to sense acceleration.
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Static balance is provided by two ventricles, the utricle and the saccule. Cells lining the walls of these ventricles contain fine filaments, and the cells are covered with a fine gelatinous layer. Each cell has 50–70 small filaments, and one large filament, the kinocilium. Within the gelatinous layer lie otoliths, tiny formations of calcium carbonate. When a person moves, these otoliths shift position. This shift alters the positions of the filaments, which opens ion channels within the cell membranes, creating depolarisation and an action potential that is transmitted to the brain along the vestibulocochlear nerve.[9][12]
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Dynamic balance is provided through the three semicircular canals. These three canals are orthogonal (at right angles) to each other. At the end of each canal is a slight enlargement, known as the ampulla, which contains numerous cells with filaments in a central area called the cupula. The fluid in these canals rotates according to the momentum of the head. When a person changes acceleration, the inertia of the fluid changes. This affects the pressure on the cupula, and results in the opening of ion channels. This causes depolarisation, which is passed as a signal to the brain along the vestibulocochlear nerve.[9] Dynamic balance also helps maintain eye tracking when moving, via the vestibulo–ocular reflex.
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During embryogenesis the ear develops as three distinct structures: the inner ear, the middle ear and the outer ear.[13] Each structure originates from a different germ layer: the ectoderm, endoderm and mesenchyme.[14][15]
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After implantation, around the second to third week the developing embryo consists of three layers: endoderm, mesoderm and ectoderm. The first part of the ear to develop is the inner ear,[15] which begins to form from the ectoderm around the 22nd day of the embryo's development.[14] Specifically, the inner ear derives from two thickenings called otic placodes on either side of the head. Each otic placode recedes below the ectoderm, forms an otic pit and then an otic vesicle.[16] This entire mass will eventually become surrounded by mesenchyme to form the bony labyrinth.[16][17]
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Around the 33rd day of development, the vesicles begin to differentiate. Closer to the back of the embryo, they form what will become the utricle and semicircular canals. Closer to the front of the embryo, the vesicles differentiate into a rudimentary saccule, which will eventually become the saccule and cochlea. Part of the saccule will eventually give rise and connect to the cochlear duct. This duct appears approximately during the sixth week and connects to the saccule through the ductus reuniens.[14]
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As the cochlear duct's mesenchyme begins to differentiate, three cavities are formed: the scala vestibuli, the scala tympani and the scala media.[14][17] Both the scala vestibuli and the scala tympani contain an extracellular fluid called perilymph. The scala media contains endolymph.[17] A set of membranes called the vestibular membrane and the basilar membrane develop to separate the cochlear duct from the vestibular duct and the tympanic duct, respectively.[14]
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Parts of the otic vesicle in turn form the vestibulocochlear nerve.[18] These form bipolar neurons which supply sensation to parts of the inner ear (namely the sensory parts of the semicircular canals, macular of the utricle and saccule, and organ of Corti). The nerve begins to form around the 28th day.[16]
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Most of the genes responsible for the regulation of inner ear formation and its morphogenesis are members of the homeobox gene family such as Pax, Msx and Otx homeobox genes. The development of inner ear structures such as the cochlea is regulated by Dlx5/Dlx6, Otx1/Otx2 and Pax2, which in turn are controlled by the master gene Shh. Shh is secreted by the notochord.[19]
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The middle ear and its components develop from the first and second pharyngeal arches.[16] The tympanic cavity and auditory tube develop from the first part of the pharyngeal pouch between the first two arches in an area which will also go on to develop the pharynx. This develops as a structure called the tubotympanic recess.[16] The ossicles (malleus, incus and stapes) normally appear during the first half of fetal development. The first two (malleus and incus) derive from the first pharyngeal arch and the stapes derives from the second.[14] All three ossicles develop from the neural crest.[16] Eventually cells from the tissue surrounding the ossicles will experience apoptosis and a new layer of endodermal epithelial will constitute the formation of the tympanic cavity wall.[14][15]
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Unlike structures of the inner and middle ear, which develop from pharyngeal pouches, the ear canal originates from the dorsal portion of the first pharyngeal cleft.[14][16] It is fully expanded by the end of the 18th week of development.[17] The eardrum is made up of three layers (ectoderm, endoderm and connective tissue). The pinna originates as a fusion of six hillocks. The first three hillocks are derived from the lower part of the first pharyngeal arch and form the tragus, crus of the helix, and helix, respectively. The final three hillocks are derived from the upper part of the second pharyngeal arch and form the antihelix, antitragus, and earlobe.[14][16][17] The outer ears develop in the lower neck. As the mandible forms they move towards their final position level with the eyes.[13][18]
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Hearing loss may be either partial or total. This may be a result of injury or damage, congenital disease, or physiological causes. When hearing loss is a result of injury or damage to the outer ear or middle ear, it is known as conductive hearing loss. When deafness is a result of injury or damage to the inner ear, vestibulochoclear nerve, or brain, it is known as sensorineural hearing loss.
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Causes of conductive hearing loss include an ear canal blocked by ear wax, ossicles that are fixed together or absent, or holes in the eardrum. Conductive hearing loss may also result from middle ear inflammation causing fluid build-up in the normally air-filled space, such as by otitis media. Tympanoplasty is the general name of the operation to repair the middle ear's eardrum and ossicles. Grafts from muscle fascia are ordinarily used to rebuild an intact eardrum. Sometimes artificial ear bones are placed to substitute for damaged ones, or a disrupted ossicular chain is rebuilt in order to conduct sound effectively.
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Hearing aids or cochlear implants may be used if the hearing loss is severe or prolonged. Hearing aids work by amplifying the sound of the local environment and are best suited to conductive hearing loss.[20] Cochlear implants transmit the sound that is heard as if it were a nervous signal, bypassing the cochlea.
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Anomalies and malformations of the pinna are common. These anomalies include chromosome syndromes such as ring 18. Children may also present cases of abnormal ear canals and low ear implantation.[15] In rare cases no pinna is formed (atresia), or is extremely small (microtia). Small pinnae can develop when the auricular hillocks do not develop properly. The ear canal can fail to develop if it does not channelise properly or if there is an obstruction.[15] Reconstructive surgery to treat hearing loss is considered as an option for children older than five,[21] with a cosmetic surgical procedure to reduce the size or change the shape of the ear is called an otoplasty. The initial medical intervention is aimed at assessing the baby's hearing and the condition of the ear canal, as well as the middle and inner ear. Depending on the results of tests, reconstruction of the outer ear is done in stages, with planning for any possible repairs of the rest of the ear.[22][23][24]
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Approximately one out of one thousand children suffer some type of congenital deafness related to the development of the inner ear.[25] Inner ear congenital anomalies are related to sensorineural hearing loss and are generally diagnosed with a computed tomography (CT) scan or a magnetic resonance imaging (MRI) scan.[21] Hearing loss problems also derive from inner ear anomalies because its development is separate from that of the middle and external ear.[15] Middle ear anomalies can occur because of errors during head and neck development. The first pharyngeal pouch syndrome associates middle ear anomalies to the malleus and incus structures as well as to the non-differentiation of the annular stapedial ligament. Temporal bone and ear canal anomalies are also related to this structure of the ear and are known to be associated with sensorineural hearing loss and conductive hearing loss.[21]
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Vertigo refers to the inappropriate perception of motion. This is due to dysfunction of the vestibular system. One common type of vertigo is benign paroxysmal positional vertigo, when an otolith is displaced from the ventricles to the semicircular canal. The displaced otolith rests on the cupola, causing a sensation of movement when there is none. Ménière's disease, labyrinthitis, strokes, and other infective and congenital diseases may also result in the perception of vertigo.[26]
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Injuries to the external ear occur fairly frequently, and can leave minor to major deformity. Injuries include: laceration, avulsion injuries, burn and repeated twisting or pulling of an ear, for discipline or torture.[27] Chronic damage to the ears can cause cauliflower ear, a common condition in boxers and wrestlers in which the cartilage around the ears becomes lumpy and distorted owing to persistence of a haematoma around the perichondrium, which can impair blood supply and healing.[28] Owing to its exposed position, the external ear is susceptible to frostbite[29] as well as skin cancers, including squamous-cell carcinoma and basal-cell carcinomas.[30]
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The ear drum may become perforated in the event of a large sound or explosion, when diving or flying (called barotrauma), or by objects inserted into the ear. Another common cause of injury is due to an infection such as otitis media.[31] These may cause a discharge from the ear called otorrhea,[32] and are often investigated by otoscopy and audiometry. Treatment may include watchful waiting, antibiotics and possibly surgery, if the injury is prolonged or the position of the ossicles is affected.[33] Skull fractures that go through the part of the skull containing the ear structures (the temporal bone) can also cause damage to the middle ear.[34] A cholesteatoma is a cyst of squamous skin cells that may develop from birth or secondary to other causes such as chronic ear infections. It may impair hearing or cause dizziness or vertigo, and is usually investigated by otoscopy and may require a CT scan. The treatment for cholesteatoma is surgery.[35]
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There are two principal damage mechanisms to the inner ear in industrialised society, and both injure hair cells. The first is exposure to elevated sound levels (noise trauma), and the second is exposure to drugs and other substances (ototoxicity). A large number of people are exposed to sound levels on a daily basis that are likely to lead to significant hearing loss.[36] The National Institute for Occupational Safety and Health has recently published research on the estimated numbers of persons with hearing difficulty (11%) and the percentage of those that can be attributed to occupational noise exposure (24%).[37] Furthermore, according to the National Health and Nutrition Examination Survey (NHANES), approximately twenty-two million (17%) US workers reported exposure to hazardous workplace noise.[38] Workers exposed to hazardous noise further exacerbate the potential for developing noise-induced hearing loss when they do not wear hearing protection.
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Tinnitus is the hearing of sound when no external sound is present.[39] While often described as a ringing, it may also sound like a clicking, hiss or roaring.[40] Rarely, unclear voices or music are heard.[41] The sound may be soft or loud, low pitched or high pitched and appear to be coming from one ear or both.[40] Most of the time, it comes on gradually.[41] In some people, the sound causes depression, anxiety, or concentration difficulties.[40]
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Tinnitus is not a disease but a symptom that can result from a number of underlying causes. One of the most common causes is noise-induced hearing loss. Other causes include: ear infections, disease of the heart or blood vessels, Ménière's disease, brain tumors, emotional stress, exposure to certain medications, a previous head injury, and earwax.[40][42] It is more common in those with depression and anxiety.[41]
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The ears have been ornamented with jewelry for thousands of years, traditionally by piercing of the earlobe. In ancient and modern cultures, ornaments have been placed to stretch and enlarge the earlobes, allowing for larger plugs to be slid into a large fleshy gap in the lobe. Tearing of the earlobe from the weight of heavy earrings, or from traumatic pull of an earring (for example, by snagging on a sweater), is fairly common.[43]
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Injury to the ears has been present since Roman times as a method of reprimand or punishment – "In Roman times, when a dispute arose that could not be settled amicably, the injured party cited the name of the person thought to be responsible before the Praetor; if the offender did not appear within the specified time limit, the complainant summoned witnesses to make statements. If they refused, as often happened, the injured party was allowed to drag them by the ear and to pinch them hard if they resisted. Hence the French expression "se faire tirer l’oreille", of which the literal meaning is "to have one's ear pulled" and the figurative meaning "to take a lot of persuading". We use the expression "to tweak (or pull) someone's ears" to mean "inflict a punishment"."[27]
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The pinnae have an effect on facial appearance. In Western societies, protruding ears (present in about 5% of ethnic Europeans) have been considered unattractive, particularly if asymmetric.[44] The first surgery to reduce the projection of prominent ears was published in the medical literature by Ernst Dieffenbach in 1845, and the first case report in 1881.[45]
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Pointy ears are a characteristic of some creatures in folklore such as the French croquemitaine, Brazilian curupira[46] or Japanese earth spider.[47] It has been a feature of characters on art as old as that of Ancient Greece[48] and medieval Europe.[49] Pointy ears are a common characteristic of many creatures in the fantasy genre,[50] including elves,[51][52][53] faeries,[54][55] pixies,[56] hobbits,[57] or orcs.[58] They are a characteristic of creatures in the horror genre, such as vampires.[59][60] Pointy ears are also found in the science fiction genre; for example among the Vulcan and Romulan races of the Star Trek universe[61] and the Nightcrawler character from the X-Men universe.[62]
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Georg von Békésy was a Hungarian biophysicist born in Budapest, Hungary. In 1961, he was awarded the Nobel Prize in Physiology or Medicine for his research on the function of the cochlea in the mammalian hearing organ.[63]
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The Vacanti mouse was a laboratory mouse that had what looked like a human ear grown on its back. The "ear" was actually an ear-shaped cartilage structure grown by seeding cow cartilage cells into a biodegradable ear-shaped mold and then implanted under the skin of the mouse; then the cartilage naturally grew by itself.[64] It was developed as an alternative to ear repair or grafting procedures and the results met with much publicity and controversy in 1997.[65][66]
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The pinna helps direct sound through the ear canal to the eardrum. The complex geometry of ridges on the inner surface of some mammalian ears helps to sharply focus sounds produced by prey, using echolocation signals. These ridges can be regarded as the acoustic equivalent of a fresnel lens, and may be seen in a wide range of animals, including the bat, aye-aye, lesser galago, bat-eared fox, mouse lemur and others.[67][68][69]
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Some large primates such as gorillas and orang-utans (and also humans) have undeveloped ear muscles that are non-functional vestigial structures, yet are still large enough to be easily identified.[70] An ear muscle that cannot move the ear, for whatever reason, has lost that biological function. This serves as evidence of homology between related species. In humans, there is variability in these muscles, such that some people are able to move their ears in various directions, and it has been said that it may be possible for others to gain such movement by repeated trials.[70] In such primates, the inability to move the ear is compensated for mainly by the ability to easily turn the head on a horizontal plane, an ability which is not common to most monkeys—a function once provided by one structure is now replaced by another.[71]
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In some animals with mobile pinnae (like the horse), each pinna can be aimed independently to better receive the sound. For these animals, the pinnae help localise the direction of the sound source.
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The ear, with its blood vessels close to the surface, is an essential thermoregulator in some land mammals, including the elephant, the fox, and the rabbit.[72] There are five types of ear carriage in domestic rabbits, some of which have been bred for exaggerated ear length[73]—a potential health risk that is controlled in some countries.[74] Abnormalities in the skull of a half-lop rabbit were studied by Charles Darwin in 1868. In marine mammals, Earless seals are one of three groups of Pinnipedia.
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Only vertebrate animals have ears, though many invertebrates detect sound using other kinds of sense organs. In insects, tympanal organs are used to hear distant sounds. They are located either on the head or elsewhere, depending on the insect family.[75] The tympanal organs of some insects are extremely sensitive, offering acute hearing beyond that of most other animals. The female cricket fly Ormia ochracea has tympanal organs on each side of her abdomen. They are connected by a thin bridge of exoskeleton and they function like a tiny pair of eardrums, but, because they are linked, they provide acute directional information. The fly uses her "ears" to detect the call of her host, a male cricket. Depending on where the song of the cricket is coming from, the fly's hearing organs will reverberate at slightly different frequencies. This difference may be as little as 50 billionths of a second, but it is enough to allow the fly to home in directly on a singing male cricket and parasitise it.[76]
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Simpler structures allow other arthropods to detect near-field sounds. Spiders and cockroaches, for example, have hairs on their legs which are used for detecting sound. Caterpillars may also have hairs on their body that perceive vibrations[77] and allow them to respond to sound.
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A pillow is a support of the body at rest for comfort, therapy, or decoration. Pillows are used by many species, including humans. Some types of pillows include throw pillows, body pillows, decorative pillows and many more.[1] Pillows that aid sleeping are a form of bedding that supports the head and neck. Other types of pillows are designed to support the body when lying down or sitting. There are also pillows that consider human body shape for increased comfort during sleep. Decorative pillows used on people, couches or chairs are sometimes referred to as cushions.[2][3]
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In contemporary western culture, pillows consist of a plain or patterned fabric envelope (known as a pillowcase) which contains a soft stuffing, typically synthetic and typically standardized in sizes and shape.[4] Pillows have been historically made of a variety of natural materials and many cultures continue to use pillows made from natural materials in the world.
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The word pillow comes from Middle English pilwe, from Old English pyle (akin to Old High German pfuliwi) and from Latin pulvinus. The first known use of the word pillow was before the 12th century.[5]
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Though the exact origin is unknown, use of pillows evolved in animals well into prehistory, the earliest examples including reptiles and mammals resting their heads on themselves, and one another, to support the head and neck.[1] Animals, including humans, evolved use of inanimate objects in their nests out of wood and stone as pillows.[1] Since domestication, many animals have also learned to make use of human-made pillows and cushions, as well as to rest on members of own and other species, for this purpose.[1]
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Sometime between 5 and 23 million years ago tree-dwelling great apes began building sleeping platforms, including wooden pillows, to improve their sleep.[6] According to studies on chimpanzees that sleep up to eight to nine hours a night using specifically selected ironwood pillows, sturdy pillows enabled great apes to escape being hunted by night predators and not fall out of the trees while asleep.[6] It is likely that this was necessitated by the evolution of large, energy-consuming brains.[6] Though it may also have led to longer periods of REM sleep, that in turn increased their cognitive capacity.[6]
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The earliest recorded use of the modern human device dates back to the civilizations of Mesopotamia around 7,000 BC.[7] During this time, only the wealthy used pillows.[7] The number of pillows symbolized status so the more pillows one owned the more affluence they held.[7] Pillows have long been produced around the world in order to help solve the reoccurring problem of neck, back, and shoulder pain while sleeping.[8] Besides for comfort, the pillow was also used for keeping bugs and insects out of people's hair, mouth, nose, and ears while sleeping.[8]
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Pillow use has been associated with the mummies and tombs of ancient Egypt during the 11th dynasty, dating to 2055–1985 B.C.[9] Ancient Egyptian pillows were wooden or stone headrests.[9] These pillows were mostly used by placing them under the heads of the deceased because the head of a human was considered to be the essence of life and sacred.[9]
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The Romans and Greeks of ancient Europe mastered the creation of the softer type pillow. These pillows were stuffed with reeds, feathers, and straw in order to make them softer and more comfortable.[10] Only upper-class people typically owned these softer pillows; however, all classes of people were allowed to use some type of pillow while sleeping, lying down or sitting in order to give them support.[10] People in ancient Europe started to use pillows when going to church in order to kneel on while praying and to place holy books on.[11] This is a tradition that still lives on today. In addition, the Romans and Greeks used their pillows by placing them under the head of those deceased just like the ancient Egyptians did.[10]
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Chinese pillows were traditionally solid, though sometimes used with a softer fabric over them. Over many Chinese dynasties, pillows were made from a wide range of materials including bamboo, jade, porcelain, wood, and bronze.[12] Ceramic pillows became the most popular.[12] The use of the ceramic pillow first appeared in the Sui Dynasty between 581 and 618 while mass production appeared in the Tang Dynasty between 618 and 907.[12] The Chinese decorated their pillows by making them different shapes and by painting pictures of animals, humans, and plants on them.[12] One common type of pottery used was Cizhou ware. Chinese ceramic pillows reached their peak in terms of production and use during the Song, Jin, and Yuan dynasties between the 10th and 14th century, but slowly phased out during the Ming and Qing dynasties between 1368 and 1911 with the emergence of better pillow making materials.[12]
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Pillows consist of a filler material enclosed in a fabric cover or shell. Covers are made of cloth, such as silk, known as the pillow case or pillow slip. Some pillows have a fancier cover called a sham which is closed on all sides and usually has a slit in the back through which the pillow is placed. Rectangular standard bed pillow cases usually do not have zippers, but instead, have one side open all the time. Often, a zippered pillow protector is often placed around standard pillows with the case in turn covering the protector.
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Fillers are chosen on the basis of comfort, resilience, thermal properties, cost and also for medical and ethical reasons. The most common synthetic fillers are materials derived from polymer fibers, such as polyester and memory foam. Synthetic fillers in general are inexpensive and in the case of memory foam can retain their form longer. Natural fillers have been used since antiquity. The most common are feathers, down, wool, latex, cotton (particularly in India), and buckwheat. More exotic materials have also been used including straw, wood, or stone. Down is usually the softer filler and offer good insulation, but is more expensive due to its relative scarcity. Down has been known to be plucked from live geese,[13] but there are already cruelty-free certifications for down products.[14] In India, traditional pillows are made with kapok, the fluffy, glossy fruit-fibres of the trees Ceiba pentandra and Bombax ceiba.
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The normal lifespan of a pillow is two to four years.[15] Replacement is recommended for sanitary reasons. All types of pillow covers should be laundered periodically since they are the part that is in contact with a person's body. Pillows accumulate dust and microbes among the fill, even when washable pillows are washed. Manufacturers recommend tumble-drying for fifteen minutes every week to freshen them up, and for the heat to kill dust mites. Charities in most countries will not accept used pillows due to hygiene regulations. While some animal shelters accept forms of bedding, most reject donation of used pillows due to the mess they can cause.
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Recycling of pillows, like most textile and bedding items, is expensive and has poor yield.[16] As such, few are recycled and most end up in landfills.[15] Their light weight means that they make up a low proportion of household waste by mass.[16] Most of the few pillows collected for recycling are sent to India and Pakistan and used as low-cost bedding, or in South East Asia, co-mingled with other textiles to manufacture cheap bedding.[16]
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A pillow is designed to provide support and comfort to the body and head. There are three main types of pillows; bed pillows, orthopedic pillows and decorative pillows, with some overlapping of use between these. The appropriate size of a bed pillow depends on the size of the bed. Larger pillows than standard are available for queen- and king-sized beds.
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The choice of bed pillow depends to some extent upon sleeping positions: one manufacturer recommends a thinner and softer pillow for sleeping face down, medium support for sleeping on one's back, and a thicker and firmer pillow for sleeping on the side.[15]
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The classic bed pillow shape is usually a square or rectangle. In the US, they are common in these three sizes (in inches): Standard (20×26 inches), Queen (20×30 inches), and King (20×36 inches). In the US, a less common size is Jumbo (20×28 inches), which is larger than the Standard Size but smaller than the Queen Size.
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Pillows are generally covered with a removable pillow case, which facilitates laundering. Apart from the color and from the material of which they are made, pillowcases have three contrasting characteristics:
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Size conforms to the pillow the case is to contain. They are typically described as:
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Square is also called continental in the UK. German pillow sizes are 80×80 cm (older) or 40×80 cm (newer). When considered as a subset of bed pillows, Euro pillows finish 26×26 in and older style travel pillows commonly finished 12×16 in.[citation needed]
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The main distinguishing feature is whether the pillow case is plain or with a valance around the edge. In the former case this is described as 'plain style' and in the latter as 'Oxford style'.[This quote needs a citation]
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"Oxford ... has a 5cm-10cm valance round all four sides. With a hemstitched or corded decoration around the inner edge of the valance."
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The opening/closure of pillow cases ranges from the straightforward "bag style" common in the United States to the "housewife style" more common in Europe, with a pocket inside the open end to fully contain the pillow.[This quote needs a citation]
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"Housewife is ... essentially a bag, with a flap in the open end to tuck the pillowcase behind to keep it in..."
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Other methods of closure are ties or buttons/buttonholes.
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Body pillows are as long as a full adult body, providing support to the head and neck at the top and to the knees and legs lower down. This type of pillow can be especially useful in providing support for those who sleep on their sides and for pregnant women. Size is 40×140 cm. (See also: Dutch wife)
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Neck pillows support the neck by providing a deep area for the head to rest and a supportive area to keep the neck in alignment with the spine while sleeping. These can also be known as cervical pillows. Cervical pillows help patients to maintain comfortable positioning after therapeutic, orthopedic and surgical measures.[17]
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Travel pillows provide support for the neck in a sitting position. Their "U" shape fits around the back of the neck and keeps the head from slipping into an uncomfortable and possibly harmful position during sleep. However, U-shaped pillows can sometimes force the head forward, creating neck stiffness.
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Doughnut pillows are firm pillows shaped like a torus, with a space in the middle to alleviate pressure on the tailbone area while sitting. These pillows are used primarily by individuals who have suffered an injury to the tailbone area, or who suffer pain from hemorrhoids or another ailment of the colon.
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Lumbar pillows are designed to support the inward curve of the lower back, filling the space created between the lower back and the back of the chair when in a sitting position. These pillows are generally used to support the lower back while driving or sitting, such as in an office chair. Orthopedic pillows are similar to memory foam pillows.
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Decorative pillows serve a dual purpose. They often have fancy cover material which serves to decorate the room where they are found. Since decorative textiles are commonly 54 inches in width, many decorative pillows finish about 17x17 inches. (54/3 = 18, less seam allowance) When used to decorate a fully made up bed, decorative pillows are likely thrown aside at bedtime, since they are not covered with a washable pillow case, thus, while found on the bed, they are primarily there for decoration, hence they fall under this category. These pillows may be custom made, as well as made by freelancers.
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Decorative pillows are also found on furnishings in more public parts of the home, such as sofas, chairs and window seats. Here, their common use may overlap both orthopedic and bed pillows. For example, unless a person has some particular medical condition, they will likely use a handy decorative pillow for lumbar support, as needed, while seated on a sofa. Likewise, for the occasional nap, decorative pillows are handy for supporting the head or neck, even though they may not be covered with a pillow case, as are bed pillows.
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There are five common synonyms for decorative pillows which are descriptive of their use in the home. "Accent" pillows emphasize or accent some other part of the home decor. The terms "sofa pillow" and "couch pillow" refer to the place these decorative pillows are likely found. The terms "toss pillow" and "throw pillow" may refer to the way they generally arrive in their places.[citation needed]
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Novelty pillows are shaped like humorous objects (a banana, tweety bird, a human leg, a chainsaw, a dill pickle, a former president) and are meant to brighten up and add humor to a room or lounge area.
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Tent-flap pillows are placed at the front of a stack of pillows decorating a bed. This pillows have a separate flap of fabric that is attached at the top of the pillow and folds down over the face. The tent flap can be loose or tacked down; if the flap is loose a decorative tassel or bead is usually used to weigh the flap down so it hangs properly.[18]
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Floor pillows are another subset of decorative pillows. These pillows often finish 26×26 inches (one half of the width of the textile, less seam allowance)
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A combination pillow is a pillow which is for those sleepers which are sleeping in a combination of positions and would like to use their pillow to hug or play. [19]
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In many parts of the world, pillows have cultural significance, and references have extended to a wide variety of other uses, forms and activities.
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In some cultures, pillows have forms for hugging. An example is the Dakimakura, a kind of "hugging pillow" originating in Japan that have been endowed with anthropomorphic and zoomorphic qualities and pop culture references for additional psychological comfort. Other types of hugging pillows are more practical, such as the Guling, long hugging pillow originating from Indonesia and the Abrazador, a long hugging pillow originating from the Philippines. A husband pillow (also known as a boyfriend pillow) is a large, high-backed pillow with two "arms." It is used to prop the user upright while in bed or on the floor, as for reading or watching television. Because of this common use, a husband pillow is also called a reading pillow.[citation needed]
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A pillow fight is a common game mostly played by young children (but also by teens and adults) in which they engage in mock physical conflict, using pillows as weapons.
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Pillow fights are known to occur during children's sleepovers. Since pillows are usually soft, injuries rarely occur. The heft of a pillow can still knock a young person off balance, especially on a soft surface such as a bed, which is a common venue. In earlier eras, pillows would often break, shedding feathers throughout a room. Modern pillows tend to be stronger and are often filled with a solid block of artificial filling, so breakage occurs far less frequently.
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Pillow talk refers to the relaxed, intimate conversation that often occurs between two sexual partners after sexual activity, usually accompanied by cuddling, caresses, and other physical intimacy. It is associated with honesty, sexual afterglow, and bonding.[20] In the western world there are many cultural references to Pillow talk.
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Chinese rock pillows (simplified Chinese: 石头枕头; traditional Chinese: 石頭枕頭; pinyin: shítou zhěntou) played an important role in ancient China. Made from jade, they were believed to translate the energy from the stone to the human brain.[citation needed] Originating in the Ming dynasty, this piece of material was trusted to cure headaches or depression, or simply to better the intelligence of those who use it. It was more common among royalty because it was expensive and rare. Families often married off children based on the quality and intricacy of these pillows. Today, rock pillows are still commonly used during the hot summer months in China.[citation needed]
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Cushioning designed for specific parts of the body are sometimes called pillows. An example is Eye pillows which are designed to comfort the eyes. Mouse pads wrist rests though not called pillows, are another examples.
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A sex pillow is a specially-designed and typically firm pillow used to enhance sexual intercourse. An ordinary firm pillow, however, may be used in place of a special one. Some contain a high-density urethane core to balance firm support with softness. In addition to more common pillow shapes, there are wedge-shaped, ramp-shaped, prism-shaped, etc. pillows which lend themselves to various sexual positions, some of which might be difficult or uncomfortable without them. A very common use of a sex pillow is to place it under the woman's buttocks or hips before assuming the missionary position; this pivots the pelvis and increases depth of penetration.[21]
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A pillow menu is a list of available pillows provided by a hotel to guests, usually free of charge.[22] It allows guests to make an alternative pillow choice. Some common pillow alternatives are memory foam, buckwheat hull, and hypoallergenic. Some hotels offer pillows to treat specific conditions such as headaches or stress.
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Adele Laurie Blue Adkins MBE (/əˈdɛl/; born 5 May 1988) is an English singer-songwriter. After graduating from the BRIT School in 2006, Adele signed a recording contract with XL Recordings. In 2007, she received the Brit Awards Critics' Choice Award and won the BBC Sound of 2008 poll. Her debut album, 19, was released in 2008. It is certified eight times platinum in the UK, and three times platinum in the US. The album contains her first song, "Hometown Glory", written when she was 16, which is based on her home suburb of West Norwood in London. An appearance she made on Saturday Night Live in late 2008 boosted her career in the US. At the 51st Grammy Awards in 2009, Adele won the awards for Best New Artist and Best Female Pop Vocal Performance.
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Adele released her second studio album, 21, in 2011. The album was critically well-received and surpassed the success of her debut, earning numerous awards in 2012, among them a record-tying six Grammy Awards, including Album of the Year; two Brit Awards, including British Album of the Year; and three American Music Awards, including Favorite Pop/Rock Album. The album has been certified 17 times platinum in the UK, and is overall the fourth best-selling album in the nation. In the US, it has held the top position longer than any album since 1985, and is certified diamond. The best-selling album worldwide of 2011 and 2012, 21 has sold over 31 million copies. The success of 21 earned Adele numerous mentions in the Guinness Book of Records. She was the first woman in the history of the Billboard Hot 100 to have three simultaneous top 10 singles as a lead artist, with "Rolling in the Deep", "Someone Like You", and "Set Fire to the Rain", all of which also topped the chart.
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In 2012, Adele released "Skyfall", which she co-wrote and recorded for the James Bond film of the same name. The song won an Academy Award, a Grammy Award, and a Golden Globe for Best Original Song, as well as the Brit Award for British Single of the Year. After taking a three-year break, Adele released her third studio album, 25, in 2015. It became the year's best-selling album and broke first-week sales records in the UK and US. 25 was her second album to be certified diamond in the US and earned her five Grammy Awards, including Album of the Year, and four Brit Awards, including British Album of the Year. The lead single, "Hello", became the first song in the US to sell over one million digital copies within a week of its release. Her third concert tour, Adele Live 2016, visited Europe, North America and Oceania, and concluded with finale concerts at Wembley Stadium in late June 2017.[7]
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In 2011, 2012, and 2016, Adele was named Artist of the Year by Billboard. At the 2012 and 2016 Ivor Novello Awards, Adele was named Songwriter of the Year by the British Academy of Songwriters, Composers, and Authors. In 2012, she was listed at number five on VH1's 100 Greatest Women in Music. Time magazine named her one of the most influential people in the world in 2012 and 2016. Her 2016–2017 tour, saw her break attendance records in a number of countries, including the UK, Australia, and the US. With sales of more than 120 million records, Adele is one of the world's best-selling music artists.[8]
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Adele Laurie Blue Adkins was born on 5 May 1988 in Tottenham, London, to an English mother, Penny Adkins, and a Welsh father, Marc Evans.[9] Evans left when Adele was two, leaving her mother to raise her.[10][11] She began singing at the age of four and asserts that she became obsessed with voices.[12][13] In 1997, at the age of nine, Adele and her mother, who by then had found work as a furniture maker and an adult-learning activities organiser, relocated to Brighton on the south coast of England.[14]
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In 1999, she and her mother moved back to London; first to Brixton, then to the neighbouring district of West Norwood in south London, which is the subject of her first song “Hometown Glory”.[15] She spent much of her youth in Brockwell Park where she would play the guitar and sing to friends, which she recalled in her 2015 song “Million Years Ago”. She stated, “It has quite monumental moments of my life that I’ve spent there, and I drove past it [in 2015] and I just literally burst into tears. I really missed it.”[16] Adele graduated from the BRIT School for Performing Arts & Technology in Croydon in May 2006,[17] where she was a classmate of Leona Lewis and Jessie J.[1][18] Adele credits the school with nurturing her talent[19] even though, at the time, she was more interested in going into A&R and hoped to launch other people's careers.[1]
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Four months after graduation, she published two songs on the fourth issue of the online arts publication PlatformsMagazine.com.[20] She had recorded a three-song demo for a class project and given it to a friend.[1] The friend posted the demo on Myspace, where it became very successful and led to a phone call from Richard Russell, boss of the music label XL Recordings. She doubted if the offer was real because the only record company she knew was Virgin Records, and she took a friend with her to the meeting.[18][21]
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Nick Huggett, at XL, recommended Adele to manager Jonathan Dickins at September Management, and in June 2006, Dickins became her official representative.[22] September was managing Jamie T at the time and this proved a major draw for Adele, a big fan of the British singer-songwriter. Huggett then signed Adele to XL in September 2006.[22] Adele provided vocals for Jack Peñate's song, "My Yvonne," for his debut album, and it was during this session she first met producer Jim Abbiss, who would go on to produce both the majority of her debut album, 19, and tracks on 21.[23] In June 2007, Adele made her television debut, performing "Daydreamer" on the BBC's Later... with Jools Holland.[24] Adele's breakthrough song, "Hometown Glory", written when she was 16, was released in October 2007.[22]
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By 2008, Adele had become the headliner and performed an acoustic set, in which she was supported by Damien Rice.[25] She became the first recipient of the Brit Awards Critics' Choice and was named the number-one predicted breakthrough act of 2008 in an annual BBC poll of music critics, Sound of 2008.[26][27] The album 19, named for her age at the time she wrote and composed many of its songs, entered the British charts at number one. The Times Encyclopedia of Modern Music named 19 an "essential" blue-eyed soul recording.[28] She released her second single, "Chasing Pavements", on 14 January 2008, two weeks ahead of her debut album, 19. The song reached number two on the UK Chart, and stayed there for four weeks.[29] Adele was nominated for a 2008 Mercury Prize award for 19.[30] She also won an Urban Music Award for "Best Jazz Act."[31] She also received a Music of Black Origin (MOBO) nomination in the category of Best UK Female.[32] In March 2008, Adele signed a deal with Columbia Records and XL Recordings for her foray into the United States.[33] She embarked on a short North American tour in the same month,[33] and 19 was released in the US in June.[19] Billboard magazine stated of it: "Adele truly has potential to become among the most respected and inspiring international artists of her generation."[34] The An Evening with Adele world tour began in May 2008 and ended in June 2009.[35]
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She later cancelled the 2008 US tour dates to be with a former boyfriend.[36] She said in Nylon magazine in June 2009, "I'm like, 'I can't believe I did that.' It seems so ungrateful.... I was drinking far too much and that was kind of the basis of my relationship with this boy. I couldn't bear to be without him, so I was like, 'Well, I'll just cancel my stuff then.'"[36] She referred to this period as her "early life crisis".[36] She is also known for her dislike of flying and bouts of homesickness when away from her native London.[37] By the middle of October 2008, Adele's attempt to break in America appeared to have failed.[38] But then she was booked as the musical guest on 18 October 2008 episode of NBC's Saturday Night Live. The episode, which included an expected appearance by then US vice-presidential candidate Sarah Palin, earned the program its best ratings in 14 years with 17 million viewers. Adele performed "Chasing Pavements" and "Cold Shoulder,"[39] and the following day, 19 topped the iTunes charts and ranked at number five at Amazon.com while "Chasing Pavements" rose into the top 25.[40] The album reached number 11 on the Billboard 200 as a result, a jump of 35 places over the previous week.[41]
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In November 2008 Adele moved to Notting Hill, London after leaving her mother's house, a move that prompted her to give up drinking.[42] The album was certified gold in early 2009, by the RIAA.[43] By July 2009, the album had sold 2.2 million copies worldwide.[44]
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At the 51st Annual Grammy Awards in February 2009, Adele won the award for Best New Artist, in addition to the award for Best Female Pop Vocal Performance for "Chasing Pavements", which was also nominated for Record of the Year and Song of the Year.[45] Adele performed "Chasing Pavements" at the ceremony in a duet with Jennifer Nettles. In 2010, Adele received a Grammy nomination for Best Female Pop Vocal Performance for "Hometown Glory."[46] In April her song "My Same" entered the German Singles Chart after it had been performed by Lena Meyer-Landrut in the talent show contest Unser Star für Oslo, or Our Star for Oslo, in which the German entry to the Eurovision Song Contest 2010 was determined.[47][48]
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In late September, after being featured on The X Factor, Adele's version of Bob Dylan's "Make You Feel My Love" re-entered the UK singles chart at number 4.[49] During the 2010 CMT Artists of the Year special, Adele performed a widely publicised duet of Lady Antebellum's "Need You Now" with Darius Rucker.[50] This performance was later nominated for a CMT Music Award.[51]
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Adele released her second studio album, 21, on 24 January 2011 in the UK and 22 February in the US.[52][53] She said that the album was inspired by the break-up with her former partner.[11] The album's sound is described as classic and contemporary country and roots music. The change in sound from her first album was the result of her bus driver playing contemporary music from Nashville when she was touring the American South, and the title reflected the growth she had experienced in the prior two years.[53] Adele told Spin Magazine: "It was really exciting for me because I never grew up around [that music]."[54] 21 hit number 1 in 30 countries, including the UK and the US.[55][56][57]
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An emotional performance of "Someone Like You" at the 2011 Brit Awards on 15 February propelled the song to number one in the UK.[58] Her first album, 19, re-entered the UK album chart alongside 21, while first and second singles "Rolling in the Deep" and "Someone Like You" were in the top 5 of the UK singles chart, making Adele the first living artist to achieve the feat of two top-five hits in both the Official Singles Chart and the Official Albums Chart simultaneously since the Beatles in 1964.[59] Both songs topped the charts in multiple markets and broke numerous sales performance records. Following her performance of "Someone Like You" at the 2011 MTV Video Music Awards, it became Adele's second number-one single on the Billboard Hot 100.[60] By December 2011, 21 sold over 3.4 million copies in the UK, and became the biggest-selling album of the 21st century, overtaking Amy Winehouse's Back to Black,[61][62] with Adele becoming the first artist ever to sell three million albums in the UK in one calendar year.[63][64] "Set Fire to the Rain" became Adele's third number-one single on the Billboard Hot 100, as Adele became the first artist ever to have an album, 21, hold the number-one position on the Billboard 200 concurrently with three number-one singles.[65] Moreover, 21 had the most weeks on the Billboard 200 chart of any album by a woman.[66]
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To promote the album, Adele embarked upon the "Adele Live" tour, which sold out its North American leg.[67] In October 2011, Adele was forced to cancel two tours because of a vocal-cord haemorrhage. She released a statement saying she needed an extended period of rest to avoid permanent damage to her voice.[68] In the first week of November 2011 Steven M. Zeitels, director of the Center for Laryngeal Surgery and Voice Rehabilitation at the Massachusetts General Hospital in Boston, performed laser microsurgery on Adele's vocal cords to remove a benign polyp.[69][70][71] A recording of her tour, Live at the Royal Albert Hall, was released in November 2011, debuting at number one in the US with 96,000 copies sold, the highest one-week tally for a music DVD in four years, becoming the best-selling music DVD of 2011.[72] Adele is the first artist in Nielsen SoundScan history to have the year's number-one album (21), number-one single ("Rolling in the Deep"), and number-one music video (Live at the Royal Albert Hall).[73] At the 2011 American Music Awards on 20 November, Adele won three awards; Favorite Pop/Rock Female Artist, Favorite Adult Contemporary Artist, and Favorite Pop/Rock Album for 21.[74] On 9 December, Billboard named Adele Artist of the Year, Billboard 200 Album of the Year (21), and the Billboard Hot 100 Song of the Year ("Rolling in the Deep"), becoming the first woman ever to top all three categories.[75][76]
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Following the throat microsurgery, she made her live comeback at the 2012 Grammy Awards in February.[77] She won in all six categories for which she was nominated, including Album of the Year, Record of the Year, and Song of the Year, making her the second female artist in Grammy history, after Beyoncé, to win that many categories in a single night.[78] Following that success, 21 achieved the biggest weekly sales increase following a Grammy win since Nielsen SoundScan began tracking data in 1991.[79][80] Adele received the Brit Award for British Female Solo Artist, and British Album of the Year presented to her by George Michael.[81][82] Following the Brit Awards, 21 reached number one for the 21st non-consecutive week in the UK.[83] The album has sold over 4.5 million copies in the UK where it is the fourth best-selling album.[84] In October, the album's sales surpassed 4.5 million in the UK, and in November it surpassed 10 million sales in the US.[85][86][87] The best-selling album worldwide of 2011 and 2012, as of 2016[update], the album has sold over 31 million copies.[88][89][90] By the end of 2014, she had sold an estimated 40 million albums and 50 million singles worldwide.[91] Adele is the only artist or band in the last decade in the US to earn an RIAA diamond certification for a one disc album in less than two years.[86]
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In October 2012, Adele confirmed that she had been writing, composing and recording the theme song for Skyfall, the twenty-third James Bond film.[92][93] The song "Skyfall," written and composed in collaboration with producer Paul Epworth, was recorded at Abbey Road Studios, and features orchestrations by J. A. C. Redford.[94] Adele stated recording "Skyfall" was "one of the proudest moments of my life." On 14 October, "Skyfall" rose to number 2 on the UK Singles Chart with sales of 92,000 copies bringing its overall sales to 176,000, and "Skyfall" entered the Billboard Hot 100 at number 8, selling 261,000 copies in the US in its first three days.[95] This tied "Skyfall" with Duran Duran's "A View to a Kill" as the highest-charting James Bond theme song on the UK Singles Chart;[96] a record surpassed in 2015 by Sam Smith's "Writing's on the Wall".[97]
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"Skyfall" has sold more than five million copies worldwide[98] and earned Adele the Golden Globe Award for Best Original Song[99] and the Academy Award for Best Original Song.[100] In December 2012, Adele was named Billboard Artist of the Year, and 21 was named Album of the Year, making her the first artist to receive both accolades two years in a row.[101][102] Adele was also named top female artist.[102] The Associated Press named Adele Entertainer of the Year for 2012.[103] The 2013 Grammy Awards saw Adele's live version of "Set Fire to the Rain" win the Grammy Award for Best Pop Solo Performance, bringing her total wins to nine.[104]
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On 3 April 2012, Adele confirmed that her third album would likely be at least two years away, stating, "I have to take time and live a little bit. There were a good two years between my first and second albums, so it'll be the same this time." She stated that she would continue writing and composing her own material.[105]
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At the 2013 Grammy Awards, she confirmed that she was in the very early stages of her third album.[106][107] She also stated that she will most likely work with Paul Epworth again.[106]
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In September 2013, Wiz Khalifa confirmed that he and Adele had collaborated on a song for his fifth studio album, Blacc Hollywood, though the collaboration did not make the final track listing.[108] In January 2014, Adele received her tenth Grammy Award with "Skyfall" winning Best Song Written for Visual Media at the 56th Annual Grammy Awards.[109]
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On the eve of her 26th birthday in May 2014, Adele posted a cryptic message via her Twitter account which prompted media discussion about her next album. The message, "Bye bye 25... See you again later in the year," was interpreted by some in the media, including Capital FM, as meaning that her next album would be titled 25 and released later in the year.[110] In 2014, Adele was nominated for nine World Music Awards. In early August, Paul Moss suggested that an album would be released in 2014 or 2015.[111] However, in the October 2014 accounts filed with Companies House by XL Recordings, they ruled out a 2014 release.[112]
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On 27 August 2015, Billboard reported that Adele's label, XL Recordings, had intentions of releasing her third studio album sometime in November 2015.[113] Danger Mouse was revealed to have contributed a song, while Tobias Jesso Jr. had written a track, and Ryan Tedder was "back in the mix after producing and co-writing "Rumour Has It" on 21."[113] At the 72nd Venice International Film Festival in early September 2015, Sia announced that her new single "Alive" was co-written by Adele, and had originally been intended for Adele's third album.[114] On 18 October, a 30-second clip of new material from Adele was shown on UK television during a commercial break on The X Factor. The commercial teases a snippet from a new song from her third album, with viewers hearing a voice singing accompanied by lyrics on a black screen.[115]
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In a statement released three days later she confirmed that the album is titled 25, with Adele stating, "My last record was a break-up record, and if I had to label this one, I would call it a make-up record. Making up for lost time. Making up for everything I ever did and never did. 25 is about getting to know who I've become without realising. And I'm sorry it took so long but, you know, life happened."[116] Adele also believes 25 will be her last album with her age as its title, believing that 25 would be the end to a trilogy.[117] On 22 October, Adele confirmed that 25 would be released on 20 November, while the lead single from the album, "Hello" would be released on 23 October.[118] The song was first played on Nick Grimshaw's Radio 1 Breakfast Show on the BBC on the morning of 23 October with Adele interviewed live.[119] The video of "Hello", released on 22 October, was viewed over 27.7 million times on YouTube in its first 24 hours, breaking the Vevo record for the most views in a day, surpassing the 20.1 million views for "Bad Blood" by Taylor Swift.[120] On 28 October, BBC News reported that "Hello" was being viewed on YouTube an average one million times an hour.[121] "Hello" went on to become the fastest video to hit one billion views on YouTube, which it achieved after 88 days.[122] The video for "Hello" captured iconic British elements such as a red telephone box and a cup of tea.[123] The song debuted at number one on the UK Singles Chart on 30 October, with first week sales of 330,000 copies, making it the biggest-selling number one single in three years.[124] "Hello" also debuted at number one in many countries around the world, including Australia, France, Canada, New Zealand, Ireland and Germany, and on 2 November, the song debuted at number one on the Billboard Hot 100, becoming the first song in the US to sell at least one million downloads in a week, setting the record at 1.11 million.[125] By the end of 2015, it had sold 12.3 million units globally and was the year's 7th best-selling single despite being released in late October.[126]
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On 27 October, BBC One announced plans for Adele at the BBC, a one-hour special presented by Graham Norton, in which Adele talks about her new album and performs new songs.[128] This was her first television appearance since performing at the 2013 Academy Awards ceremony, and the show was recorded before a live audience on 2 November for broadcast on 20 November, coinciding with the release of 25.[129] On 27 October it was also announced that Adele would appear on the US entertainment series Saturday Night Live on 21 November.[128][130] On 30 October, Adele confirmed that she would be performing a one-night-only concert titled Adele Live in New York City at the Radio City Music Hall on 17 November. Subsequently, NBC aired the concert special on 14 December.[131][132]
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On 27 November, 25 debuted at number one on the UK Albums Chart and became the fastest selling album in UK chart history with over 800,000 copies sold in its first week.[133] The album debuted at number one in the US where it sold a record-breaking 3.38 million copies in its first week, the largest single sales week for an album since Nielsen began monitoring sales in 1991.[134] 25 also broke first week sales records in Canada and New Zealand.[135][136] 25 became the best-selling album of 2015 in a number of countries, including Australia, the UK and the US, spending seven consecutive weeks at number one in each country, before being displaced by David Bowie's Blackstar.[137][138][139] It was the best-selling album worldwide of 2015 with 17.4 million copies sold.[126] 25 has since sold 20 million copies globally.[140] Adele's seven weeks at the top of the UK Albums Chart took her total to 31 weeks at number one in the UK with her three albums, surpassing Madonna's previous record of most weeks at number one for a female act.[141] As the best-selling artist worldwide for 2015 the IFPI named Adele the Global Recording Artist of the Year.[142]
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In November 2015, Adele's 2016 tour was announced, her first tour since 2011.[143] Beginning in Europe, Adele Live 2016 included four dates at the Manchester Arena in March 2016, six dates at the O2 Arena, London, with further dates in Ireland, Spain, Germany, Italy and the Netherlands among others.[144] Her North American Tour began on 5 July in St. Paul, Minnesota.[145] The leg included six nights at Madison Square Garden in New York City, eight nights at Staples Center in Los Angeles, and four nights at Air Canada Centre in Toronto.[146] Adele broke Taylor Swift's five-show record for most consecutive sold-out shows at the Staples Center.[147]
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At the 2016 Brit Awards in London on 24 February, Adele received the awards for British Female Solo Artist, British Album of the Year for 25, British Single of the Year for "Hello", and British Global Success, bringing her Brit Award wins to eight.[149] She closed the ceremony by performing "When We Were Young", the second single from 25.[149] Two more singles from 25 were released in 2016: "Send My Love (To Your New Lover)" and "Water Under the Bridge". While on stage at London's O2 Arena on 17 March, Adele announced that she would be headlining on the Pyramid Stage at the 2016 Glastonbury Festival, which was later confirmed by the festival's organisers.[150] She appeared for a 90-minute fifteen song set at the festival on 25 June in front of 150,000 people, and described the experience as "by far, the best moment of my life so far".[151][152] In an interview with Jo Whiley on BBC Radio 2 around 30-minutes before going on stage, Adele had said she had been going to Glastonbury since she was a child and that the festival had meant a lot to her, before she broke down. Whiley recalls, "She was really scared, really, really scared. We were doing the interview and at one point she had to stop as she was in tears. It was amazing to see somebody like that, then to witness her walking out on stage and doing the most incredible set. To know that half an hour before she’d been in tears at the thought of walking out there."[148]
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As part of her world tour, in February and March 2017, Adele performed in Australia for the first time, playing outdoor stadiums around the country.[153] Her first two shows in New Zealand sold out in a record-breaking 23 minutes, and a third show was announced, with all tickets sold in under 30 minutes.[154] Adele sold over 600,000 tickets for her record-breaking eight date Australian tour, setting stadium records throughout the country; her Sydney show at ANZ Stadium on 10 March was seen by 95,000 people, the biggest single concert in Australian history, a record she broke the following night with more than 100,000 fans.[155]
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Adele completed her world tour with two concerts, dubbed "The Finale", at Wembley Stadium, London on 28 and 29 June.[7] She announced the shows at "the home of football" by singing the England football team's "Three Lions" anthem and also the theme song to the BBC's weekly Premier League football show Match of the Day.[7] Adele had added another two concerts at Wembley after the first two dates sold out,[156] however she cancelled the last two dates of the tour after damaging her vocal cords.[157] As a show of support, fans instead gathered outside Wembley Stadium to perform renditions of her songs, in an event titled "Sing for Adele".[158]
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At the end of 2016, Billboard named Adele Artist of the Year for the third time,[159] and also received the Top Billboard 200 album.[160] 25 was the best-selling album for a second consecutive year in the US.[161] With 213 million views, Adele's Carpool Karaoke through the streets of London with James Corden, a sketch which featured on Corden's talk show The Late Late Show with James Corden in January 2016, was the biggest YouTube viral video of 2016.[162] At the 59th Annual Grammy Awards in February 2017, Adele won all five of her nominations, bringing her number of awards to fifteen. She won Album of the Year and Best Pop Vocal Album for 25, and Record of the Year, Song of the Year and Best Pop Solo Performance for "Hello".[163] She also performed a tribute to the late George Michael singing the rendition of his song "Fastlove"; due to technical difficulties which occurred during the performance, Adele decided to stop and restart, explaining "I can't mess this up for him".[164]
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Adele was reportedly working on her fourth studio album by 2018.[165] On 5 May 2019, the date of her 31st birthday, Adele posted several black-and-white pictures of herself on her Instagram account celebrating her birthday along with a message reflecting on the preceding year. The message ended with "30 will be a drum n bass record to spite you". Media outlets, including NME, took the post as an indication that a new album was on the way.[166][167] On 15 February 2020, Adele announced at a friend's wedding that her fourth studio album would be out in September 2020.[168] However, she later stated the album release has been delayed due to the COVID-19 pandemic.[169]
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Adele has cited the Spice Girls as a major influence in regard to her love and passion for music, stating that "they made me what I am today".[170] Adele impersonated the Spice Girls at dinner parties as a young girl.[171] She stated she was left "heartbroken" when her favourite Spice Girl, Geri Halliwell aka "Ginger Spice", left the group.[172][173] Growing up she also listened to Sinéad O'Connor,[174] the Cranberries,[175] Bob Marley,[176] the Cure,[177] Dusty Springfield,[178] Celine Dion,[179] and Annie Lennox.[180]
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One of Adele's earliest influences was Gabrielle, who Adele has admired since the age of five. During Adele's school years, her mother made her an eye patch with sequins which she used to perform as the Hackney born star in a school talent contest.[181] After moving to south London, she became interested in R&B artists such as Aaliyah, Destiny's Child, and Mary J. Blige.[182] Adele says that one of the most defining moments in her life was when she watched Pink perform at Brixton Academy in London. She states: "It was the Missundaztood record, so I was about 13 or 14. I had never heard, being in the room, someone sing like that live [...] I remember sort of feeling like I was in a wind tunnel, her voice just hitting me. It was incredible."[183][184] Adele also cites Jeff Buckley's album Grace as an influence, saying "I remember falling out with my best friend when I was like seven and listening to Jeff Buckley, because my mum was a huge fan. Grace has always been around me".[185]
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In 2002, 14-year-old Adele discovered Etta James and Ella Fitzgerald as she stumbled on the artists' CDs in the jazz section of her local music store. She was struck by their appearance on the album covers.[186] Adele states she then "started listening to Etta James every night for an hour," and in the process was getting "to know my own voice."[186] Adele credits Amy Winehouse and her 2003 album Frank for inspiring her to take up the guitar, stating, "If it wasn't for Amy and Frank, one hundred per cent I wouldn't have picked up a guitar, I wouldn't have written "Daydreamer" or "Hometown [Glory]" and I wrote "Someone Like You" on the guitar too."[187] She also states that her mother, who is very close to her, exposed her to the music of Aaliyah, Lauryn Hill, Mary J. Blige, and Alicia Keys, all of whom inspired her as well.[174] On the rock band Queen, she states, "I love them. I'm the biggest Queen fan ever. They're the kind of band that's just in your DNA, really. Everyone just knows who they are."[188] She is also a fan of Lana Del Rey, Grimes, Chvrches, FKA Twigs, Alabama Shakes, Kanye West, Rihanna, Britney Spears, Frank Ocean, and Stevie Nicks.[189][190][191][192] In 2017, she described Beyoncé as a particular inspiration, calling her "my artist of my life" and added "the other artists who mean that much to me are all dead."[193] Adele cited Madonna's 1998 album Ray of Light as a "chief inspiration" behind her album 25.[190] Adele mentioned that Taylor Swift is the inspiration behind her song "Send My Love (To Your New Lover)", with Adele stating "I was in New York, writing "Remedy" with Ryan Tedder. We were having lunch, and "I Knew You Were Trouble" came on the radio — Taylor's song that she did with Max Martin and Shellback. I was like, 'Who did this?' I knew it was Taylor, and I've always loved her, but this is a totally other side — like, 'I want to know who brought that out in her.' And he said Max Martin. I was unaware that I knew who Max Martin was. I Googled him, and I was like, 'He's literally written every massive soundtrack of my life.' So I got my management to reach out. They came to London, and I took my guitar along and was like, 'I've got this riff,' and then "Send My Love" happened really quickly. Max Martin, I just could hang out with him forever. He's so beautiful and lovely and funny and generous and warm and caring. He's a really amazing man."[194]
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Adele's debut album, 19, is of the soul genre, with lyrics describing heartbreak and relationship.[19] Her success occurred simultaneously with several other British female soul singers, with the British press dubbing her a new Amy Winehouse.[1] This was described as a third British Musical Invasion of the US.[18] However, Adele called the comparisons between her and other female soul singers lazy, noting "we're a gender, not a genre".[19][195][196] AllMusic wrote that "Adele is simply too magical to compare her to anyone."[186]
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Her second album, 21, shares the folk and soul influences of her debut album, but was further inspired by the American country and Southern blues music to which she had been exposed during her 2008–09 North American tour An Evening with Adele.[197][198] Composed in the aftermath of Adele's separation from her partner, the album typifies the near dormant tradition of the confessional singer-songwriter in its exploration of heartbreak, self-examination, and forgiveness. Having referred to 21 as a "break-up record", Adele labelled her third studio album, 25, a "make-up record", adding it was about "Making up for lost time. Making up for everything I ever did and never did."[116] Her yearning for her old self, her nostalgia, and melancholy about the passage of time, is a feature of 25, with Adele stating, "I've had a lot of regrets since I turned 25. And sadness hits me in different ways than it used to. There's a lot of things I don't think I'll ever get 'round to doing."[199]
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—Dave Simpson of The Guardian on her voice and down to earth manner.[200]
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Adele is a mezzo-soprano with a range spanning from B2 to C6. However Classic FM state she is often mistaken for a contralto due to the application of a tense chest mix to achieve her lower notes, while also noting that her voice becomes its clearest as she ascends the register, particularly from C4 to C5.[201][202][203][204] Rolling Stone reported that following throat surgery her voice had become "palpably bigger and purer-toned", and that she had added a further four notes to the top of her range.[199] Initially, critics suggested that her vocals were more developed and intriguing than her songwriting, a sentiment with which Adele agreed.[205] She has stated: "I taught myself how to sing by listening to Ella Fitzgerald for acrobatics and scales, Etta James for passion and Roberta Flack for control."[206]
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Adele's singing voice has been acclaimed by critics. In a review of 19, The Observer said, "The way she stretched the vowels, her wonderful soulful phrasing, the sheer unadulterated pleasure of her voice, stood out all the more; little doubt that she's a rare singer".[207] BBC Music wrote, "Her melodies exude warmth, her singing is occasionally stunning and, ...she has tracks that make Lily Allen and Kate Nash sound every bit as ordinary as they are."[208] Also in 2008, Sylvia Patterson of The Guardian wrote, "Of all the gobby new girls, only Adele's bewitching singing voice has the enigmatic quality which causes tears of involuntary emotion to splash down your face in the way Eva Cassidy’s did before her."[209] For their reviews of 21, The New York Times' chief music critic Jon Pareles commended Adele's emotive timbre, comparing her to Dusty Springfield, Petula Clark, and Annie Lennox: "[Adele] can seethe, sob, rasp, swoop, lilt and belt, in ways that draw more attention to the song than to the singer".[210] Ryan Reed of Paste magazine regarded her voice as "a raspy, aged-beyond-its-years thing of full-blooded beauty",[211] while MSN Music's Tom Townshend declared her "the finest singer of [our] generation".[212]
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Adele began dating charity entrepreneur and Old Etonian Simon Konecki in the summer of 2011.[213] In June 2012, Adele announced that she and Konecki were expecting a child.[214][215] Their son Angelo James was born on 19 October 2012.[216] On the topic of becoming a parent, Adele observed that she "felt like I was truly living. I had a purpose, where before I didn't".[217] Adele and Konecki brought a privacy case against a UK-based photo agency that published intrusive paparazzi images of their son taken during family outings in 2013.[218] Lawyers working on their behalf accepted damages from the company in July 2014.[219] Adele has also stated that she has suffered from postnatal depression, generalized anxiety disorder, and panic attacks.[220][221]
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In early 2017, tabloids started speculating that Adele and Konecki had secretly married when they were spotted wearing matching rings on their ring fingers.[222] During her acceptance speech at the 59th Annual Grammy Awards for Album of the Year, Adele confirmed these reports of their marriage by calling Konecki her husband when thanking him.[223] She repeated this in March 2017, telling the audience at a concert in Brisbane, Australia, "I'm married now".[224] Adele became a stay-at-home mother.[225] In April 2019, Adele's representatives announced to the Associated Press that she and Konecki had separated after more than seven years together, but that they would continue to raise their son together.[226][227] On 13 September 2019, it was reported that Adele had filed for divorce from Konecki in the US.[228]
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Politically, she is a supporter of the Labour Party, stating in 2011 that she was a "Labour girl through and through", and in the same interview was critical of the Conservative Party.[229] Despite this declared political affiliation, Adele received backlash after comments about paying taxes during a 2011 interview with Q magazine. She stated, "I use the NHS, I can't use public transport any more, doing what I do, I went to state school, I'm mortified to have to pay 50 percent! Trains are always late, most state schools are shit and I've gotta give you like four million quid, are you having a laugh? When I got my tax bill in from 19 I was ready to go and buy a gun and randomly open fire."[230][231] In 2015, Adele stated "I'm a feminist, I believe that everyone should be treated the same, including race and sexuality".[190]
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Born in the North London district of Tottenham, Adele supports local football club Tottenham Hotspur.[232] In 2017, Adele was ranked the richest musician under 30 years old in the UK and Ireland in the Sunday Times Rich List, which valued her wealth at £125 million. She was ranked the 19th richest musician overall.[233] On the 2019 list, she was valued at £150 million as the 22nd richest musician in the UK.[234]
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Supportive of the LGBT community, on 12 June 2016, an emotional Adele dedicated her show in Antwerp, Belgium to the victims of the mass shooting at a gay nightclub in Orlando, Florida, United States earlier that day, adding "The LGBTQ community, they're like my soul mates since I was really young, so I'm very moved by it."[235][236] In April 2018, it was widely reported that Adele had become an ordained minister in order to officiate at close friend comedian Alan Carr's wedding to Paul Drayton, something which Adele herself subsequently confirmed. The wedding, held in January 2018, took place in the garden of her house in Los Angeles, California.[237]
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Adele has performed in numerous charity concerts throughout her career. In 2007 and 2008 she performed at the Little Noise Sessions held at London's Union Chapel, with proceeds from the concerts donated to Mencap which works with people with learning disabilities.[37] In July and November 2008, Adele performed at the Keep a Child Alive Black Ball in London and New York City respectively.[238][239][240] On 17 September 2009 she performed at the Brooklyn Academy of Music, for the VH1 Divas event, a concert to raise money for the Save The Music Foundation charity.[241][242] On 6 December, Adele opened with a 40-minute set at John Mayer's 2nd Annual Holiday Charity Revue held at the Nokia Theatre in Los Angeles.[243] In 2011, Adele gave a free concert for Pride London, a registered charity which arranges LGBT events in London.[244] The same year, Adele took part in the UK charity telethon Comic Relief for Red Nose Day 2011, performing "Someone like You".[245]
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Adele has been a major contributor to MusiCares, a charity organisation founded by the National Academy of Recording Arts and Sciences for musicians in need. In February 2009, Adele performed at the 2009 MusiCares charity concert in Los Angeles. In 2011 and 2012, Adele donated autographed items for auctions to support MusiCares.[246][247][248] Adele required all backstage visitors to the North American leg of her Adele Live tour to donate a minimum charitable contribution of US$20 for the UK charity SANDS, an organisation dedicated to "supporting anyone affected by the death of a baby and promoting research to reduce the loss of babies' lives".[249]
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On 15 June 2017, Adele attended a vigil in west London for the victims of the Grenfell Tower fire where, keeping a low profile, she was only spotted by a small handful of fans.[250] Four days later she appeared at Chelsea fire station and brought cakes for the firefighters.[251] Station manager Ben King stated "She came in, came up to the mess and had a cup of tea with the watch and then she joined us for the minute's silence."[251] Paying tribute to the victims at her first Wembley show on 28 June, Adele encouraged fans to donate money to help the victims of the blaze rather than waste the money on "overpriced wine".[252]
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At the 51st Annual Grammy Awards in 2009, Adele won awards in the categories of Best New Artist and Best Female Pop Vocal Performance.[253] She was also nominated in the categories of Record of the Year and Song of the Year.[254] The success of her debut album 19 saw Adele nominated for three Brit Awards in the categories of British Female Solo Artist, British Single of the Year and British Breakthrough Act.[255] Then British Prime Minister Gordon Brown sent a thank-you letter to Adele that stated "with the troubles that the country's in financially, you're a light at the end of the tunnel".[256]
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Adele's second album, 21, earned her a record-tying six Grammy Awards, including Album of the Year; two Brit Awards, including British Album of the Year. Adele was the second artist and first female, preceded by Christopher Cross, to have won all four of the general field awards throughout her career.[257] The success of the album saw her receive numerous mentions in the Guinness Book of World Records.[258] With 21 non-consecutive weeks at number 1 in the US, Adele broke the record for the longest number-1 album by a woman in Billboard history, beating the record formerly held by Whitney Houston's soundtrack The Bodyguard.[80] 21 spent its 23rd week at number one in March 2012, making it the longest-running album at number one since 1985,[259] and it became the fourth best-selling album of the past 10 years in the US.[260] The best selling album in the UK of the 21st century, and the best selling album by a female in UK chart history, 21 is also the fourth best-selling album in the UK of all time.[261][262] 21 was her first album certified diamond in the US.[263] On 6 March, 21 reached 30 non-consecutive weeks at number one on the Australian ARIA Chart, making it the longest-running number one album in Australia in the 21st century, and the second longest-running number one ever.[264]
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In February 2012, Adele was listed at number five on VH1's 100 Greatest Women in Music.[265] In April 2012, Time magazine named Adele one of the 100 most influential people in the world.[266][267] People named her one of 2012 Most Beautiful at Every Age.[268] On 30 April 2012, a tribute to Adele was held at New York City's (Le) Poisson Rouge called Broadway Sings Adele, starring various Broadway actors such as Matt Doyle.[269] In July 2012, Adele was listed at number six in Forbes list of the world's highest-paid celebrities under the age of 30, having earned £23 million between May 2011 and May 2012.[270]
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On the week ending 3 March 2012, Adele became the first solo female artist to have three singles in the top 10 of the Billboard Hot 100 at the same time with "Rolling in the Deep", "Someone Like You", and "Set Fire to the Rain" as well as the first female artist to have two albums in the top 5 of the Billboard 200 and two singles in the top 5 of the Billboard Hot 100 simultaneously.[271] Adele topped the 2012 Sunday Times Rich List of musicians in the UK under 30,[272] and made the Top 10 of Billboard magazine's "Top 40 Money Makers".[273] Billboard also announced the same day that Adele's "Rolling in the Deep" is the biggest crossover hit of the past 25 years, topping pop, adult pop and adult contemporary charts and that Adele is one of four female artists to have an album chart at number one for more than 13 weeks (the other three artists being Judy Garland, Carole King, and Whitney Houston).[273]
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At the 2012 Ivor Novello Awards in May, Adele was named Songwriter of the Year, and "Rolling in the Deep" won the award for Most Performed Work of 2011.[274] At the 2012 BMI Awards held in London in October, Adele won Song of the Year (for "Rolling in the Deep") in recognition of the song being the most played on US television and radio in 2011.[275] In 2013, Adele won the Academy Award for Best Original Song for the James Bond theme "Skyfall". This is the first James Bond song to win and the fifth to be nominated—after "For Your Eyes Only" (1981), "Nobody Does It Better" (1977), "Live and Let Die" (1973), and "The Look of Love" (1967).[276][277] "Skyfall" won the Brit Award for Best British Single at the 2013 Brit Awards.[278]
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In June 2013, Adele was appointed a MBE in the Queen's Birthday Honours list for services to music, and she received the award from Prince Charles at Buckingham Palace on 19 December 2013.[279][280] In February 2013 she was named one of the 100 most powerful women in the UK by Woman's Hour on BBC Radio 4.[281] In April 2016, Adele appeared for the second time on the Time 100 list of most influential people.[282] In 2014, Adele was already being regarded as a British cultural icon, with young adults from abroad naming her among a group of people that they most associated with UK culture, which included William Shakespeare, Queen Elizabeth II, David Beckham, J. K. Rowling, The Beatles, Charlie Chaplin and Elton John.[283][284]
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Released in 2015, Adele's third album, 25, became the year's best-selling album and broke first week sales records in a number of markets, including the UK and the US.[285] 25 was her second album to be certified diamond in the US and earned her five Grammy Awards, including her second Grammy Award for Album of the Year, and four Brit Awards for British Female Solo Artist, British Album of the Year, British Single of the Year for "Hello", and British Global Success.[149] Adele became the only artist in history to, on two separate occasions, win the three general categories Grammys in the same ceremony.[286] With 15 awards from 18 nominations, Adele won more Grammys than any other female who was born outside the U.S.[287] Adele's seven weeks at the top of the UK Albums Chart took her total to 31 weeks at number one in the UK with her three albums, surpassing Madonna's previous record of most weeks at number one for a female act in the UK.[141] The lead single, "Hello", became the first song in the US to sell over one million digital copies within a week of its release.[125] At the 2016 Ivor Novello Awards Adele was named Songwriter of the Year by the British Academy of Songwriters, Composers, and Authors.[288] As of 6 August 2019, despite releasing just two albums in the decade (21 and 25), at 36 weeks she has the second most weeks at number one in the UK Album Charts in the 2010s, five weeks behind Ed Sheeran (who has released four albums).[289][290] In December 2019, Israel's largest TV and Radio stations named her singer of the 2010s.[291]
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Athens (/ˈæθɪnz/ ATH-inz;[3] Greek: Αθήνα, romanized: Athína [aˈθina] (listen); Ancient Greek: Ἀθῆναι, romanized: Athênai (pl.) [atʰɛ̂ːnai̯]), is the capital and largest city of Greece. Athens dominates the Attica region and is one of the world's oldest cities, with its recorded history spanning over 3,400 years[citation needed] and its earliest human presence started somewhere between the 11th and 7th millennium BC.[4]
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Classical Athens was a powerful city-state that emerged in conjunction with the seagoing development of the port of Piraeus. A center for the arts, learning and philosophy, home of Plato's Academy and Aristotle's Lyceum,[5][6] it is widely referred to as the cradle of Western civilization and the birthplace of democracy,[7][8] largely because of its cultural and political impact on the European continent, and in particular the Romans.[9] In modern times, Athens is a large cosmopolitan metropolis and central to economic, financial, industrial, maritime, political and cultural life in Greece.
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Athens is a global city and one of the biggest economic centers in southeastern Europe. It has a large financial sector, and its port Piraeus is both the largest passenger port in Europe,[10][11] and the second largest in the world.[12]
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The Municipality of Athens (also City of Athens), which actually constitutes a small administrative unit of the entire city, had a population of 664,046 (in 2011)[2] within its official limits, and a land area of 38.96 km2 (15.04 sq mi).[13][14] The Athens Urban Area (Greater Athens and Greater Piraeus) extends beyond its administrative municipal city limits, with a population of 3,090,508 (in 2011)[15] over an area of 412 km2 (159 sq mi).[14] According to Eurostat[16] in 2011, the functional urban area (FUA) of Athens was the 9th most populous FUA in the European Union (the 6th most populous capital city of the EU), with a population of 3.8 million people. Athens is also the southernmost capital on the European mainland and the warmest major city in Europe.
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The heritage of the classical era is still evident in the city, represented by ancient monuments and works of art, the most famous of all being the Parthenon, considered a key landmark of early Western civilization. The city also retains Roman and Byzantine monuments, as well as a smaller number of Ottoman monuments. Athens is home to two UNESCO World Heritage Sites, the Acropolis of Athens and the medieval Daphni Monastery. Landmarks of the modern era, dating back to the establishment of Athens as the capital of the independent Greek state in 1834, include the Hellenic Parliament and the so-called "architectural trilogy of Athens", consisting of the National Library of Greece, the National and Kapodistrian University of Athens and the Academy of Athens. Athens is also home to several museums and cultural institutions, such as the National Archeological Museum, featuring the world's largest collection of ancient Greek antiquities, the Acropolis Museum, the Museum of Cycladic Art, the Benaki Museum and the Byzantine and Christian Museum. Athens was the host city of the first modern-day Olympic Games in 1896, and 108 years later it hosted the 2004 Summer Olympics, making it one of the few cities to have hosted the Olympics more than once.[17]
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In Ancient Greek, the name of the city was Ἀθῆναι (Athênai, pronounced [atʰɛ̂ːnai̯] in Classical Attic) a plural. In earlier Greek, such as Homeric Greek, the name had been current in the singular form though, as Ἀθήνη (Athḗnē).[18] It was possibly rendered in the plural later on, like those of Θῆβαι (Thêbai) and Μυκῆναι (Μukênai). The root of the word is probably not of Greek or Indo-European origin,[19] and is possibly a remnant of the Pre-Greek substrate of Attica.[19] In antiquity, it was debated whether Athens took its name from its patron goddess Athena (Attic Ἀθηνᾶ, Athēnâ, Ionic Ἀθήνη, Athḗnē, and Doric Ἀθάνα, Athā́nā) or Athena took her name from the city.[20] Modern scholars now generally agree that the goddess takes her name from the city,[20] because the ending -ene is common in names of locations, but rare for personal names.[20] During the medieval period, the name of the city was rendered once again in the singular as Ἀθήνα. However, after the establishment of the modern Greek state, and partly due to the conservatism of the written language, Ἀθῆναι [aˈθine] became again the official name of the city and remained so until the abandonment of Katharevousa in the 1970s, when Ἀθήνα, Athína, became the official name.
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According to the ancient Athenian founding myth, Athena, the goddess of wisdom, competed against Poseidon, the God of the Seas, for patronage of the yet-unnamed city;[21] they agreed that whoever gave the Athenians the better gift would become their patron[21] and appointed Cecrops, the king of Athens, as the judge.[21] According to the account given by Pseudo-Apollodorus, Poseidon struck the ground with his trident and a salt water spring welled up.[21] In an alternative version of the myth from Vergil's Georgics, Poseidon instead gave the Athenians the first horse.[21] In both versions, Athena offered the Athenians the first domesticated olive tree.[21][22] Cecrops accepted this gift[21] and declared Athena the patron goddess of Athens.[21][22] Eight different etymologies, now commonly rejected, were proposed during the 17th century.[citation needed] Christian Lobeck proposed as the root of the name the word ἄθος (áthos) or ἄνθος (ánthos) meaning "flower", to denote Athens as the "flowering city". Ludwig von Döderlein proposed the stem of the verb θάω, stem θη- (tháō, thē-, "to suck") to denote Athens as having fertile soil.[23]
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In classical literature, the city was sometimes referred to as the City of the Violet Crown, first documented in Pindar's ἰοστέφανοι Ἀθᾶναι (iostéphanoi Athânai), or as τὸ κλεινὸν ἄστυ (tò kleinòn ásty, "the glorious city"). In medieval texts, variant names include Setines, Satine, and Astines, all derivations involving false splitting of prepositional phrases.[24] Today the caption η πρωτεύουσα (ī protévousa), "the capital", has become somewhat common. King Alphonse X of Castile in his 'General Storia' (Global History) cites etymology of Athens as: 'Athinai', 'the one without death', death here meaning ignorance.
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Athenians were called cicada-wearers (Ancient Greek: Τεττιγοφόροι) because they used to wear pins of golden cicadas. A symbol of being autochthon (earth-born), because the legendary founder of Athens, Erechtheus was an autochthon or of being musicians, because the cicada is a "musician" insect.[25]
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The oldest known human presence in Athens is the Cave of Schist, which has been dated to between the 11th and 7th millennia BC.[4] Athens has been continuously inhabited for at least 5,000 years.[26][27] By 1400 BC the settlement had become an important centre of the Mycenaean civilization and the Acropolis was the site of a major Mycenaean fortress, whose remains can be recognised from sections of the characteristic Cyclopean walls.[28] Unlike other Mycenaean centers, such as Mycenae and Pylos, it is not known whether Athens suffered destruction in about 1200 BC, an event often attributed to a Dorian invasion, and the Athenians always maintained that they were pure Ionians with no Dorian element. However, Athens, like many other Bronze Age settlements, went into economic decline for around 150 years afterwards.
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Iron Age burials, in the Kerameikos and other locations, are often richly provided for and demonstrate that from 900 BC onwards Athens was one of the leading centres of trade and prosperity in the region.[29] The leading position of Athens may well have resulted from its central location in the Greek world, its secure stronghold on the Acropolis and its access to the sea, which gave it a natural advantage over inland rivals such as Thebes and Sparta.
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By the 6th century BC, widespread social unrest led to the reforms of Solon. These would pave the way for the eventual introduction of democracy by Cleisthenes in 508 BC. Athens had by this time become a significant naval power with a large fleet, and helped the rebellion of the Ionian cities against Persian rule. In the ensuing Greco-Persian Wars Athens, together with Sparta, led the coalition of Greek states that would eventually repel the Persians, defeating them decisively at Marathon in 490 BC, and crucially at Salamis in 480 BC. However, this did not prevent Athens from being captured and sacked twice by the Persians within one year, after a heroic but ultimately failed resistance at Thermopylae by Spartans and other Greeks led by King Leonidas,[30] after both Boeotia and Attica fell to the Persians.
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The decades that followed became known as the Golden Age of Athenian democracy, during which time Athens became the leading city of Ancient Greece, with its cultural achievements laying the foundations for Western civilization. The playwrights Aeschylus, Sophocles and Euripides flourished in Athens during this time, as did the historians Herodotus and Thucydides, the physician Hippocrates, and the philosopher Socrates. Guided by Pericles, who promoted the arts and fostered democracy, Athens embarked on an ambitious building program that saw the construction of the Acropolis of Athens (including the Parthenon), as well as empire-building via the Delian League. Originally intended as an association of Greek city-states to continue the fight against the Persians, the league soon turned into a vehicle for Athens's own imperial ambitions. The resulting tensions brought about the Peloponnesian War (431–404 BC), in which Athens was defeated by its rival Sparta.
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By the mid-4th century BC, the northern Greek kingdom of Macedon was becoming dominant in Athenian affairs. In 338 BC the armies of Philip II defeated an alliance of some of the Greek city-states including Athens and Thebes at the Battle of Chaeronea, effectively ending Athenian independence. Later, under Rome, Athens was given the status of a free city because of its widely admired schools. The Roman emperor Hadrian, in the 2nd century CE, ordered the construction of a library, a gymnasium, an aqueduct which is still in use, several temples and sanctuaries, a bridge and financed the completion of the Temple of Olympian Zeus.
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By the end of Late Antiquity, Athens had shrunk due to sacks by the Herulians, Visigoths, and Early Slavs which caused massive destruction in the city. In this era, the first Christian churches were built in Athens, and the Parthenon and other temples were converted into churches. Athens expanded its settlement in the second half of the Middle Byzantine Period, in the 9th to 10th centuries CE, and was relatively prosperous during the Crusades, benefiting from Italian trade. After the Fourth Crusade the Duchy of Athens was established. In 1458 it was conquered by the Ottoman Empire and entered a long period of decline.
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Following the Greek War of Independence and the establishment of the Greek Kingdom, Athens was chosen as the capital of the newly independent Greek state in 1834, largely because of historical and sentimental reasons. At the time, it was reduced to a town of about 4,000 people in a loose swarm of houses along the foot of the Acropolis. The first King of Greece, Otto of Bavaria, commissioned the architects Stamatios Kleanthis and Eduard Schaubert to design a modern city plan fit for the capital of a state.
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The first modern city plan consisted of a triangle defined by the Acropolis, the ancient cemetery of Kerameikos and the new palace of the Bavarian king (now housing the Greek Parliament), so as to highlight the continuity between modern and ancient Athens. Neoclassicism, the international style of this epoch, was the architectural style through which Bavarian, French and Greek architects such as Hansen, Klenze, Boulanger or Kaftantzoglou designed the first important public buildings of the new capital. In 1896, Athens hosted the first modern Olympic Games. During the 1920s a number of Greek refugees, expelled from Asia Minor after the Greco-Turkish War, swelled Athens's population; nevertheless it was most particularly following World War II, and from the 1950s and 1960s, that the population of the city exploded, and Athens experienced a gradual expansion.
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In the 1980s it became evident that smog from factories and an ever-increasing fleet of automobiles, as well as a lack of adequate free space due to congestion, had evolved into the city's most important challenge. A series of anti-pollution measures taken by the city's authorities in the 1990s, combined with a substantial improvement of the city's infrastructure (including the Attiki Odos motorway, the expansion of the Athens Metro, and the new Athens International Airport), considerably alleviated pollution and transformed Athens into a much more functional city. In 2004 Athens hosted the 2004 Summer Olympics.
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Tondo of the Aison Cup, showing the victory of Theseus over the Minotaur in the presence of Athena. Theseus was responsible, according to the myth, for the synoikismos ("dwelling together")—the political unification of Attica under Athens.
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The earliest coinage of Athens, c. 545–525/15 BC
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Coat of Arms of the Duchy of Athens during the rule of the de la Roche family (13th century)
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The Roman Agora and the Gate of Athena in Plaka district.
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The Temple of Olympian Zeus with river Ilisos by Edward Dodwell, 1821
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The Entry of King Otto in Athens, Peter von Hess, 1839.
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The Stadiou Street in Central Athens in 1908.
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Athens sprawls across the central plain of Attica that is often referred to as the Athens Basin or the Attica Basin (Greek: Λεκανοπέδιο Αθηνών/Αττικής). The basin is bounded by four large mountains: Mount Aigaleo to the west, Mount Parnitha to the north, Mount Pentelicus to the northeast and Mount Hymettus to the east.[31] Beyond Mount Aegaleo lies the Thriasian plain, which forms an extension of the central plain to the west. The Saronic Gulf lies to the southwest. Mount Parnitha is the tallest of the four mountains (1,413 m (4,636 ft)),[32] and has been declared a national park.
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Athens is built around a number of hills. Lycabettus is one of the tallest hills of the city proper and provides a view of the entire Attica Basin. The meteorology of Athens is deemed to be one of the most complex in the world because its mountains cause a temperature inversion phenomenon which, along with the Greek Government's difficulties controlling industrial pollution, was responsible for the air pollution problems the city has faced.[27] This issue is not unique to Athens; for instance, Los Angeles and Mexico City also suffer from similar atmospheric inversion problems.[27]
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The Cephissus river, the Ilisos and the Eridanos stream are the historical rivers of Athens.
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By the late 1970s, the pollution of Athens had become so destructive that according to the then Greek Minister of Culture, Constantine Trypanis, "...the carved details on the five the caryatids of the Erechtheum had seriously degenerated, while the face of the horseman on the Parthenon's west side was all but obliterated."[33] A series of measures taken by the authorities of the city throughout the 1990s resulted in the improvement of air quality; the appearance of smog (or nefos as the Athenians used to call it) has become less common.
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Measures taken by the Greek authorities throughout the 1990s have improved the quality of air over the Attica Basin. Nevertheless, air pollution still remains an issue for Athens, particularly during the hottest summer days. In late June 2007,[34] the Attica region experienced a number of brush fires,[34] including a blaze that burned a significant portion of a large forested national park in Mount Parnitha,[35] considered critical to maintaining a better air quality in Athens all year round.[34] Damage to the park has led to worries over a stalling in the improvement of air quality in the city.[34]
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The major waste management efforts undertaken in the last decade (particularly the plant built on the small island of Psytalia) have greatly improved water quality in the Saronic Gulf, and the coastal waters of Athens are now accessible again to swimmers. In January 2007, Athens faced a waste management problem when its landfill near Ano Liosia, an Athenian suburb, reached capacity.[36] The crisis eased by mid-January when authorities began taking the garbage to a temporary landfill.[36]
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Athens ranks in the lowest percentage for the risk on frequency and severity of terrorist attacks according to the EU Global Terrorism Database (EIU 2007-2016 calculations). The city also ranked 35th in Digital Security, 21st on Health Security, 29th on Infrastructure Security and 41st on Personal Security globally in a 2017 The Economist Intelligence Unit report.[37] It also ranks as a very safe city (39th globally out of 162 cities overall) on the ranking of the safest and most dangerous countries.[38] A 2019 crime index from Numbeo places Athens at 130th position, rating safer than Tampa, Florida or Dublin, Ireland.[39] According to a Mercer 2019 Quality of Living Survey, Athens ranks 89th on the Mercer Quality of Living Survey ranking.[40]
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Athens has a hot-summer Mediterranean climate (Köppen climate classification: Csa). The dominant feature of Athens' climate is alternation between prolonged hot and dry summers and mild winters with moderate rainfall.[41] With an average of 455 millimetres (17.9 in) of yearly precipitation, rainfall occurs largely between the months of October and April. July and August are the driest months when thunderstorms occur sparsely. Furthermore, some coastal areas of Athens, known as the Athens Riviera have a hot semi-arid climate (BSh) according to the climate atlas published by the Hellenic National Meteorological Service (H.N.M.S).[42]
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Owing to the rain shadow of the Pindus Mountains, annual precipitation of Athens is lower than most other parts of Greece, especially western Greece. As an example, Ioannina receives around 1,300 mm (51 in) per year, and Agrinio around 800 mm (31 in) per year. Daily average highs for July have been measured around 34 °C or 93 °F in downtown Athens, but some parts of the city may be even hotter for the higher density of buildings, such as the center,[43] in particular, western areas due to a combination of industrialization and a number of natural factors, knowledge of which has existed since the mid-19th century.[44][45][46]
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Athens is affected by the urban heat island effect in some areas which is caused by human activity,[47][48] altering its temperatures compared to the surrounding rural areas,[49][50][51][52] and leaving detrimental effects on energy usage, expenditure for cooling,[53][54] and health.[48] The urban heat island of the city has also been found to be partially responsible for alterations of the climatological temperature time-series of specific Athens meteorological stations, because of its impact on the temperatures and the temperature trends recorded by some meteorological stations.[55][56][57][58][59] On the other hand, specific meteorological stations, such as the National Garden station and Thiseio meteorological station, are less affected or do not experience the urban heat island.[49][60]
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Athens holds the World Meteorological Organization record for the highest temperature ever recorded in Europe, at 48 °C (118.4 °F), which was recorded in the Elefsina and Tatoi suburbs of Athens on 10 July 1977.[61]
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The Municipality of Athens, the City Centre of the Athens Urban Area, is divided into several districts: Omonoia, Syntagma, Exarcheia, Agios Nikolaos, Neapolis, Lykavittos, Lofos Strefi, Lofos Finopoulou, Lofos Filopappou, Pedion Areos, Metaxourgeio, Aghios Kostantinos, Larissa Station, Kerameikos, Psiri, Monastiraki, Gazi, Thission, Kapnikarea, Aghia Irini, Aerides, Anafiotika, Plaka, Acropolis, Pnyka, Makrygianni, Lofos Ardittou, Zappeion, Aghios Spyridon, Pangrati, Kolonaki, Dexameni, Evaggelismos, Gouva, Aghios Ioannis, Neos Kosmos, Koukaki, Kynosargous, Fix, Ano Petralona, Kato Petralona, Rouf, Votanikos, Profitis Daniil, Akadimia Platonos, Kolonos, Kolokynthou, Attikis Square, Lofos Skouze, Sepolia, Kypseli, Aghios Meletios, Nea Kypseli, Gyzi, Polygono, Ampelokipoi, Panormou-Gerokomeio, Pentagono, Ellinorosson, Nea Filothei, Ano Kypseli, Tourkovounia-Lofos Patatsou, Lofos Elikonos, Koliatsou, Thymarakia, Kato Patisia, Treis Gefyres, Aghios Eleftherios, Ano Patisia, Kypriadou, Menidi, Prompona, Aghios Panteleimonas, Pangrati, Goudi, Vyronas and Ilisia.
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Parnitha National Park is punctuated by well-marked paths, gorges, springs, torrents and caves dotting the protected area. Hiking and mountain-biking in all four mountains are popular outdoor activities for residents of the city. The National Garden of Athens was completed in 1840 and is a green refuge of 15.5 hectares in the centre of the Greek capital. It is to be found between the Parliament and Zappeion buildings, the latter of which maintains its own garden of seven hectares.
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Parts of the City Centre have been redeveloped under a masterplan called the Unification of Archeological Sites of Athens, which has also gathered funding from the EU to help enhance the project.[68][69] The landmark Dionysiou Areopagitou Street has been pedestrianised, forming a scenic route. The route starts from the Temple of Olympian Zeus at Vasilissis Olgas Avenue, continues under the southern slopes of the Acropolis near Plaka, and finishes just beyond the Temple of Hephaestus in Thiseio. The route in its entirety provides visitors with views of the Parthenon and the Agora (the meeting point of ancient Athenians), away from the busy City Centre.
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The hills of Athens also provide green space. Lycabettus, Philopappos hill and the area around it, including Pnyx and Ardettos hill, are planted with pines and other trees, with the character of a small forest rather than typical metropolitan parkland. Also to be found is the Pedion tou Areos (Field of Mars) of 27.7 hectares, near the National Archaeological Museum.
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Athens' largest zoo is the Attica Zoological Park, a 20-hectare (49-acre) private zoo located in the suburb of Spata. The zoo is home to around 2000 animals representing 400 species, and is open 365 days a year. Smaller zoos exist within public gardens or parks, such as the zoo within the National Garden of Athens.
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The Athens Metropolitan Area consists of 58[15] densely populated municipalities, sprawling around the Municipality of Athens (the City Centre) in virtually all directions. For the Athenians, all the urban municipalities surrounding the City Centre are called suburbs. According to their geographic location in relation to the City of Athens, the suburbs are divided into four zones; the northern suburbs (including Agios Stefanos, Dionysos, Ekali, Nea Erythraia, Kifissia, Kryoneri, Maroussi, Pefki, Lykovrysi, Metamorfosi, Nea Ionia, Nea Filadelfeia, Irakleio, Vrilissia, Melissia, Penteli, Chalandri, Agia Paraskevi, Gerakas, Pallini, Galatsi, Psychiko and Filothei); the southern suburbs (including Alimos, Nea Smyrni, Moschato, Tavros, Agios Ioannis Rentis, Kallithea, Piraeus, Agios Dimitrios, Palaio Faliro, Elliniko, Glyfada, Lagonisi, Saronida, Argyroupoli, Ilioupoli, Varkiza, Voula, Vari and Vouliagmeni); the eastern suburbs (including Zografou, Dafni, Vyronas, Kaisariani, Cholargos and Papagou); and the western suburbs (including Peristeri, Ilion, Egaleo, Koridallos, Agia Varvara, Keratsini, Perama, Nikaia, Drapetsona, Chaidari, Petroupoli, Agioi Anargyroi, Ano Liosia, Aspropyrgos, Eleusina, Acharnes and Kamatero).
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The Athens city coastline, extending from the major commercial port of Piraeus to the southernmost suburb of Varkiza for some 25 km (20 mi),[70] is also connected to the City Centre by tram.
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In the northern suburb of Maroussi, the upgraded main Olympic Complex (known by its Greek acronym OAKA) dominates the skyline. The area has been redeveloped according to a design by the Spanish architect Santiago Calatrava, with steel arches, landscaped gardens, fountains, futuristic glass, and a landmark new blue glass roof which was added to the main stadium. A second Olympic complex, next to the sea at the beach of Palaio Faliro, also features modern stadia, shops and an elevated esplanade. Work is underway to transform the grounds of the old Athens Airport – named Elliniko – in the southern suburbs, into one of the largest landscaped parks in Europe, to be named the Hellenikon Metropolitan Park.[71]
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Many of the southern suburbs (such as Alimos, Palaio Faliro, Elliniko, Glyfada, Voula, Vouliagmeni and Varkiza) known as the Athens Riviera, host a number of sandy beaches, most of which are operated by the Greek National Tourism Organisation and require an entrance fee. Casinos operate on both Mount Parnitha, some 25 km (16 mi)[72] from downtown Athens (accessible by car or cable car), and the nearby town of Loutraki (accessible by car via the Athens – Corinth National Highway, or the suburban rail service Proastiakos).
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The large City Centre (Greek: Κέντρο της Αθήνας) of the Greek capital falls directly within the Municipality of Athens or Athens Municipality (Greek: Δήμος Αθηναίων)—also City of Athens. Athens Municipality is the largest in population size in Greece. Piraeus also forms a significant city centre on its own,[73] within the Athens Urban Area and it is the second largest in population size within it, with Peristeri, Kallithea and Kypseli following.
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The Athens Urban Area (Greek: Πολεοδομικό Συγκρότημα Αθηνών), also known as Urban Area of the Capital (Greek: Πολεοδομικό Συγκρότημα Πρωτεύουσας), today consists of 40 municipalities, 35 of which make up what is referred to as the Greater Athens municipalities, located within 4 regional units (North Athens, West Athens, Central Athens, South Athens); and a further 5, which make up the Greater Piraeus municipalities, located within the regional unit of Piraeus as mentioned above. The densely built up urban area of the Greek capital sprawls across 412 km2 (159 sq mi)[14] throughout the Attica Basin and has a total population of 3,074,160 (in 2011).
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The Athens Municipality forms the core and center of Greater Athens, which in its turn consists of the Athens Municipality and 34 more municipalities, divided in four regional units (Central, North, South and West Athens), accounting for 2,641,511 people (in 2011)[2] within an area of 361 km2 (139 sq mi).[14] Until 2010, these four regional units made up the abolished Athens Prefecture. The municipality of Piraeus, the historic Athenian port, with its 4 suburban municipalities make up the regional unit of Piraeus, which in turn forms Greater Piraeus.
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Greater Athens and Greater Piraeus with part of East[74] and West Attica[75] regional units combined make up the continuous Athens Urban Area,[75][76][77] also called the "Urban Area of the Capital" or simply "Athens" (the most common use of the term), spanning over 412 km2 (159 sq mi),[78] with a population of 3,090,508 people as of 2011. The Athens Urban Area is considered to form the city of Athens as a whole, despite its administrative divisions, which is the largest in Greece and one of the most populated urban areas in Europe.
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The Athens Metropolitan Area (Greek: Μητροπολιτική Περιοχή της Αθήνας), also commonly known in English as Greater Athens,[79] spans 2,928.717 km2 (1,131 sq mi) within the Attica region and includes a total of 58 municipalities, which are organized in seven regional units (those outlined above, along with East Attica and West Attica), having reached a population of 3,737,550 based on the preliminary results of the 2011 census. Athens and Piraeus municipalities serve as the two metropolitan centres of the Athens Metropolitan Area.[80] There are also some inter-municipal centres serving specific areas. For example, Kifissia and Glyfada serve as inter-municipal centres for northern and southern suburbs respectively.
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The Municipality of Athens has an official population of 664,046 people.[2] The four regional units that make up what is referred to as Greater Athens have a combined population of 2,640,701. They together with the regional unit of Piraeus (Greater Piraeus) make up the dense Athens Urban Area which reaches a total population of 3,090,508 inhabitants (in 2011).[15] As Eurostat the FUA of Athens had in 2013 3,828,434 inhabitants, being apparently decreasing compared with the pre-economic crisis date of 2009 (4,164,175)[16]
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The municipality (Center) of Athens is the most populous in Greece, with a population of 664,046 people (in 2011)[2] and an area of 38.96 km2 (15.04 sq mi),[13] forming the core of the Athens Urban Area within the Attica Basin. The incumbent Mayor of Athens is Kostas Bakoyannis of New Democracy. The municipality is divided into seven municipal districts which are mainly used for administrative purposes.
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As of the 2011 census, the population for each of the seven municipal districts of Athens is as follows:[81]
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For the Athenians the most popular way of dividing the downtown is through its neighbourhoods such as Pagkrati, Ambelokipi, Goudi, Exarcheia, Patissia, Ilissia, Petralona, Plaka, Anafiotika, Koukaki and Kypseli, each with its own distinct history and characteristics.
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The Athens Metropolitan Area, with an area of 2,928.717 km2 (1,131 sq mi) and inhabited by 3,753,783 people in 2011,[2] consists of the Athens Urban Area with the addition of the towns and villages of East and West Attica, which surround the dense urban area of the Greek capital. It actually sprawls over the whole peninsula of Attica, which is the best part of the region of Attica, excluding the islands.
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Mycenean Athens in 1600–1100 BC could have reached the size of Tiryns; that would put the population at the range of 10,000–15,000.[82] During the Greek Dark Ages the population of Athens was around 4,000 people. In 700 BC the population grew to 10,000. In 500 BC the area probably contained 200,000 people. During the classical period the city's population is estimated from 150,000–350,000 and up to 610,000 according to Thucydides. When Demetrius of Phalerum conducted a population census in 317 BC the population was 21,000 free citizens, plus 10,000 resident aliens and 400,000 slaves. This suggests a total population of 431,000.[83][84][85][86][87][88][89][90][91] This figure is highly suspect because of the lopsided number of slaves and does not include free women and children and resident foreigners: an estimated based on Thucydides is: 40,000 male citizens, 100,000 family members, 70,000 metics (resident foreigners) and 150,000-400,000 slaves. However the numbers would include all of Athenian territory which consisted almost entirely of modern Attica region expect the territory of the city-state of Megaris and the island section.
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The ancient site of the main city is centred on the rocky hill of the acropolis. In the whole of Athenian territory they existed many towns. Acharnae, Afidnes, Cytherus, Colonus, Corydallus, Cropia, Decelea, Euonymos, Vravron among others was important towns in Athenian countryside. The new port of Piraeus was a prototype harbour with the infrastructure and housing located in the site between modern passenger section of the port (named Kantharos in ancient times) and the Pasalimani harbour (named Zea in ancient times). The old one Phaliro was in the site of modern Palaio Faliro and gradually declined after the construction of the new prototype port but remained as a minor port and important settlement with historic significance in late Classical times. The rapid expansion of the modern city, which continues to this day, was initiated in the 1950s and 1960s, because of Greece's transition from an agricultural to an industrial nation.[92] The expansion is now particularly toward the East and North East (a tendency greatly related to the new Eleftherios Venizelos International Airport and the Attiki Odos, the freeway that cuts across Attica). By this process Athens has engulfed many former suburbs and villages in Attica, and continues to do so. The table below shows the historical population of Athens in recent times.
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Athens became the capital of Greece in 1834, following Nafplion, which was the provisional capital from 1829. The municipality (City) of Athens is also the capital of the Attica region. The term Athens can refer either to the Municipality of Athens, to Greater Athens, or to the entire Athens Urban Area.
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The Hellenic Parliament
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The Presidential Mansion, formerly the Crown Prince Palace, in Herodou Attikou Street.
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The Maximos Mansion, official office of the Prime Minister of the Hellenic Republic, in Herodou Attikou Street.
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The Athens City Hall in Kotzia Square was designed by Panagiotis Kolkas and completed in 1874.[98]
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The Embassy of France in Vasilissis Sofias Avenue.
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The Italian Embassy in Vasilissis Sofias Avenue.
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Athens is twinned with:[99]
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Athens is the financial capital of Greece. According to data from 2014, Athens as a metropolitan economic area produced 130 billion US-dollars as GDP in PPP, which consists nearly a half of the production for the whole country. In the list with the strongest economic metropoles of the world Athens was ranked that year 102nd, while the GDP per capita for the same year was 32,000 US-dollars.[112]
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Athens is one of the major economic centres in south-eastern Europe and is considered as a regional economic power in Europe generally. The port of Piraeus, where big investments by COSCO have already been delivered during the recent decade, the completion of the new Cargo Centre in Thriasion,[113] the expansion of the Athens Metro and the Athens Tram, as well as the projected metropolitan park in Elliniko and other economic projects are the economic landmarks of the upcoming years.
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Important Greek companies such as Hellenic Aerospace Industry, Mytilineos Holdings, Titan Cement, Hellenic Petroleum, Papadopoulos E.J., Folli Follie, Jumbo S.A., OPAP, and COSMOTE have their headquarters in the metropolitan area of Athens. Multinational companies such as Ericsson, Sony, Siemens, Motorola, Samsung, Microsoft, Novartis, Mondelez, Coca-Cola, etc. have their regional research and development headquarters also there.
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The banking sector is represented by National Bank of Greece, Alpha Bank, Eurobank, and Piraeus Bank, while the Bank of Greece is also situated in the City Centre. The Athens Stock Exchange, the only in Greece, has been severely hit by the Greek government-debt crisis and the decision of the government to proceed into capital controls during summer 2015. As a whole the economy of Athens and Greece has been severely hit with today's data showing a change from long recession to growth of 1.4% in 2017.[114]
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Tourism is also a great contributor for the economy of the city, which is considered as one of the top destinations in Europe for city-break tourism and is also the gateway for excursions to the islands or the mainland. Greece attracted 26.5 million visitors in 2015, 30.1 million visitors in 2017 and over 33 million in 2018, making Greece one of the most visited countries in Europe and the world, and contributing 18% to the nation's Gross Domestic Product. Athens welcomed more than 5 million tourists in 2018 and 1,4 million of them were "city-breakers" (in 2013 the city-breakers were only 220.000).[115]
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Athens is the major transportation hub of Greece. They combines the largest Airport in Greece as also the largest port in Greece which also the largest port in Mediterranean in containers transport and the largest passenger port in Europe. Is also major national hub for Intercity (Ktel) and international buses as well as for domestic and international rail transport. Public transport is serviced by a variety of transportation means, forming the largest mass transit system of Greece. The Athens Mass Transit System consists of a large bus fleet, a trolleybus fleet that mainly serves Athens's city center, the city's Metro, a Proastiakos service[116] and a tram network, connecting the southern suburbs to the City Centre.[117]
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OSY (Greek: ΟΣΥ) (Odikes Sygkoinonies S.A.) which is subsidiary company of OASA (Athens urban transport organisation), is the main operator of buses and trolleybusses in Athens. Its network consists of about 300 bus lines and 22 trolleybus lines which span the Athens Metropolitan Area,[118] with a fleet of 1,839 buses and 366 trolleybuses.[119] Of those 1,839 buses 416 run on compressed natural gas,[119] making up the largest fleet of natural gas-powered buses in Europe[120] and all trolleybusses are equipped to enable them to run on diesel in case of power failure.[121]
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International links are provided by a number of private companies. National and regional bus links are provided by KTEL from two InterCity Bus Terminals, Kifissos Bus Terminal A and Liosion Bus Terminal B, both located in the north-western part of the city. Kifissos provides connections towards Peloponnese, North Greece, West Greece and some Ionian Islands, whereas Liosion is used for most of Central Greece.
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The Athens Metro is operated by STASY S.A (Greek: ΣΤΑΣΥ) (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation) and provides public transport throughout the Athens Urban Area. While its main purpose is transport, it also houses Greek artifacts found during construction of the system.[122] The Athens Metro has an operating staff of 387 and runs three metro lines, namely the line 1 (Green Line), line 2 (Red Line) and line 3 (Blue Line) lines, of which the first was constructed
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in 1869, and the other two largely during the 1990s, with the initial new sections opened in January 2000. Line 1 mostly runs at ground level and the other two (Line 2 & 3) routes run entirely underground. A fleet of 42 trains, using 252 carriages, operates on the network,[123] with a daily occupancy of 1,353,000 passengers.[124]
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The line 1 (Green Line) serves 24 stations, and is the oldest line of the Athens metro network. It runs from Piraeus station to Kifissia station and covers a distance of 25.6-kilometre (15.9 mi). There are transfer connections with the Blue Line 3 at Monastiraki station and with the Red Line 2 at Omonoia and Attiki stations.
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The line 2 (Red Line) runs from Anthoupoli station to Elliniko station and covers a distance of 17.5 km (10.9 mi).[123] The line connects the western suburbs of Athens with the southeast suburbs, passing through the center of Athens. The Red Line has transfer connections with the Green Line 1 at Attiki and Omonoia stations. There are also transfer connections with the Blue Line 3 at Syntagma station and with the tram at Syntagma, Syngrou Fix and Neos Kosmos stations.
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The line 3 (Blue Line) runs from Agia Marina and Egaleo stations, through the central Monastiraki and Syntagma stations to Doukissis Plakentias avenue in the northeastern suburb of Halandri, covering a distance of 16 km (10 mi).[123] It then ascends to ground level and continues to Eleftherios Venizelos International Airport, using the Suburban Railway infrastructure, extending its total length to 39 km (24 mi).[123] The spring 2007 extension from Monastiraki westwards to Egaleo connected some of the main night life hubs of the city, namely those of Gazi (Kerameikos station) with Psirri (Monastiraki station) and the City Centre (Syntagma station). Extensions are under construction to the west southwest suburbs of Athens, as far as the port of Piraeus. The new stations will be Agia Barvara, Koridallos, Nikaia, Maniatika, Piraeus and Dimotiko Theatro. The completed extension will be ready in 2022 (the first three stations will open in summer 2020), connecting the biggest port of Greece, Piraeus Port, with the biggest airport of Greece the Athens International Airport.
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The Athens commuter rail service, referred to as the "Proastiakós", connects Eleftherios Venizelos International Airport to the city of Kiato, 106 km (66 mi)[125] west of Athens, via Larissa station, the city's central rail station and the port of Piraeus. The length of Athens's commuter rail network extends to 120 km (75 mi),[125] and is expected to stretch to 281 km (175 mi) by 2010.[125]
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Athens tram is operated by STASY S.A (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation). It has a fleet of 35 Sirio type vehicles[126] which serve 48 stations,[126] employ 345 people with an average daily occupancy of 65,000 passengers.[126] The tram network spans a total length of 27 km (17 mi) and covers ten Athenian suburbs.[126] The network runs from Syntagma Square to the southwestern suburb of Palaio Faliro, where the line splits in two branches; the first runs along the Athens coastline toward the southern suburb of Voula, while the other heads toward Neo Faliro. The network covers the majority of the Athens coastline.[127] Further extensions are planned towards the major commercial port of Piraeus.[126] The expansion to Piraeus will include 12 new stations, increase the overall length of tram route by 5.4 km (3 mi), and increase the overall transportation network.[128]
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Athens is served by the Athens International Airport (ATH), located near the town of Spata, in the eastern Messoghia plain, some 35 km (22 mi) east of center of Athens.[129] The airport, awarded the "European Airport of the Year 2004" Award,[130] is intended as an expandable hub for air travel in southeastern Europe and was constructed in 51 months, costing 2.2 billion euros. It employs a staff of 14,000.[130]
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The airport is served by the Metro, the suburban rail, buses to Piraeus port, Athens' City Centre, Liosion and Kifisos Intercity bus stations and Elliniko metro's line 2 southern terminal, and also taxis. The airport accommodates 65 landings and take-offs per hour,[129] with its 24-passenger boarding bridges,[129] 144 check-in counters and broader 150,000 m2 (1,614,587 sq ft) main terminal;[129] and a commercial area of 7,000 m2 (75,347 sq ft) which includes cafés, duty-free shops,[130] and a small museum.
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In 2018, the airport handled 24,135,736 a huge increase over the last 4 years. In 2014, the airport handled 15,196,369 passengers, an increase of 21.2% over the previous year of 2013.[131] Of those 15,196,369 passengers, 5,267,593 passed through the airport for domestic flights,[132] and 9,970,006 passengers travelled through for international flights.[132] Beyond the dimensions of its passenger capacity, ATH handled 205,294 total flights in 2007, or approximately 562 flights per day.[133]
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Athens is the hub of the country's national railway system (OSE), connecting the capital with major cities across Greece and abroad (Istanbul, Sofia, Belgrade and Bucharest). The Port of Piraeus is the largest port in Greece and one of the largest in Europe. It is the largest container port in East Mediterranean Sea Basin. It is also the busiest passenger port in Europe and one of the largest passenger ports in the world. Connects Athens to the numerous Greek islands of the Aegean Sea, with ferries departing, while also serving the cruise ships that arrive.[134][135][136] Rafina and Lavrio act as alternative ports of Athens, connects the city with numerous Greek islands of the Aegean Sea, Evia and Cesme in Turkey,[137][138] while also serving the cruise ships that arrive.
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Two main motorways of Greece begin in Athens, namely the A1/E75, heading north towards Greece's second largest city, Thessaloniki; and the border crossing of Evzones and the A8/E94 heading west, towards Greece's third largest city, Patras, which incorporated the GR-8A. Before their completion much of the road traffic used the GR-1 and the GR-8.
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Athens' Metropolitan Area is served by the motorway network of the Attiki Odos toll-motorway (code: A6). Its main section extends from the western industrial suburb of Elefsina to Athens International Airport; while two beltways, namely the Aigaleo Beltway (A65) and the Hymettus Beltway (A64) serve parts of western and eastern Athens respectively. The span of the Attiki Odos in all its length is 65 km (40 mi),[139] making it the largest metropolitan motorway network in all of Greece.
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Located on Panepistimiou Street, the old campus of the University of Athens, the National Library, and the Athens Academy form the "Athens Trilogy" built in the mid-19th century. Most of the university's workings have been moved to a much larger, modern campus located in the eastern suburb of Zografou. The second higher education institution in the city is the Athens Polytechnic School, found in Patission Street. This was the location where on 17 November 1973, more than 13 students were killed and hundreds injured inside the university during the Athens Polytechnic uprising,[140] against the military junta that ruled the nation from 21 April 1967 until 23 July 1974.
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Other universities that lie within Athens are the Athens University of Economics and Business, the Panteion University, the Agricultural University of Athens and the University of Piraeus. There are overall ten state-supported Institutions of Higher (or Tertiary) education located in the Athens Urban Area, these are by chronological order: Athens School of Fine Arts (1837), National Technical University of Athens (1837), National and Kapodistrian University of Athens (1837), Agricultural University of Athens (1920), Athens University of Economics and Business (1920), Panteion University of Social and Political Sciences (1927), University of Piraeus (1938), Harokopio University of Athens (1990), School of Pedagogical and Technological Education (2002), University of West Attica (2018). There are also several other private colleges, as they called formally in Greece, as the establishment of private universities is prohibited by the constitution. Many of them are accredited by a foreign state or university such as the American College of Greece and the Athens Campus of the University of Indianapolis.[141]
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The city is a world centre of archaeological research. Along with national institutions, such as the Athens University and the Archaeological Society, there are multiple archaeological Museums including the National Archaeological Museum, the Cycladic Museum, the Epigraphic Museum, the Byzantine & Christian Museum, as well as museums at the ancient Agora, Acropolis, Kerameikos, and the Kerameikos Archaeological Museum. The city is also home to the Demokritos laboratory for Archaeometry, alongside regional and national archaeological authorities that form part of the Greek Department of Culture.
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Athens hosts 17 Foreign Archaeological Institutes which promote and facilitate research by scholars from their home countries. As a result, Athens has more than a dozen archaeological libraries and three specialized archaeological laboratories, and is the venue of several hundred specialized lectures, conferences and seminars, as well as dozens of archaeological exhibitions, each year. At any given time, hundreds of international scholars and researchers in all disciplines of archaeology are to be found in the city.
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Athens incorporates architectural styles ranging from Greco-Roman and Neoclassical to modern times. They are often to be found in the same areas, as Athens is not marked by a uniformity of architectural style. A visitor will quickly notice the absence of tall buildings: Athens has very strict height restriction laws in order to ensure the Acropolis hill is visible throughout the city. There is evidence of continuity in elements of the architectural environment through the city's history, although this is still a contested subject [142].
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For the greatest part of the 19th century Neoclassicism dominated Athens, as well as some deviations from it such as Eclecticism, especially in the early 20th century. Thus, the Old Royal Palace was the first important public building to be built, between 1836 and 1843. Later in the mid and late 19th century, Theophil Freiherr von Hansen and Ernst Ziller took part in the construction of many neoclassical buildings such as the Athens Academy and the Zappeion Hall. Ziller also designed many private mansions in the centre of Athens which gradually became public, usually through donations, such as Schliemann's Iliou Melathron.
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Beginning in the 1920s, modern architecture including Bauhaus and Art Deco began to exert an influence on almost all Greek architects, and buildings both public and private were constructed in accordance with these styles. Localities with a great number of such buildings include Kolonaki, and some areas of the centre of the city; neighbourhoods developed in this period include Kypseli.[143]
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In the 1950s and 1960s during the extension and development of Athens, other modern movements such as the International style played an important role. The centre of Athens was largely rebuilt, leading to the demolition of a number of neoclassical buildings. The architects of this era employed materials such as glass, marble and aluminium, and many blended modern and classical elements.[144] After World War II, internationally known architects to have designed and built in the city included Walter Gropius, with his design for the US Embassy, and, among others, Eero Saarinen, in his postwar design for the east terminal of the Ellinikon Airport.
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All over the city can be found several statues or busts. Apart from the neoclassicals by Leonidas Drosis at the Academy of Athens (Plato, Socrates, Apollo, Athena), other notable include the statue of Theseus by Georgios Fytalis at Thiseion, of philhellenes like Lord Byron, George Canning and William Gladstone, the equestrian statue of Theodoros Kolokotronis by Lazaros Sochos in front of the Old Parliament, statues of Ioannis Kapodistrias, Rigas Feraios and Adamantios Korais at the University, of Evangelos Zappas and Konstantinos Zappas at Zappeion, of Ioannis Varvakis at the National Garden, the "woodbreaker" by Dimitrios Filippotis, the equestrian statue of Alexandros Papagos at Papagou district and various busts of fighters of Greek independence at the Pedion tou Areos. A significant landmark is also the Tomb of the Unknown Soldier in Syntagma.
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Athens' most important museums include:
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Athens has been a destination for travellers since antiquity. Over the past decade, the city's infrastructure and social amenities have improved, in part because of its successful bid to stage the 2004 Olympic Games. The Greek Government, aided by the EU, has funded major infrastructure projects such as the state-of-the-art Eleftherios Venizelos International Airport,[145] the expansion of the Athens Metro system,[68] and the new Attiki Odos Motorway.[68]
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Athens was voted as the third best European city to visit in 2015 by European Best Destination. More than 240,000 people voted.
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Athens is home to 148 theatrical stages, more than any other city in the world, including the ancient Odeon of Herodes Atticus, home to the Athens Festival, which runs from May to October each year.[146][147] In addition to a large number of multiplexes, Athens plays host to open air garden cinemas. The city also supports music venues, including the Athens Concert Hall (Megaro Moussikis), which attracts world class artists.[148] The Athens Planetarium,[149] located in Andrea Syngrou Avenue, in Palaio Faliro[150] is one of the largest and best equipped digital planetaria in the world.[151] The Stavros Niarchos Foundation Cultural Center, inaugurated in 2016, will house the National Library of Greece and the Greek National Opera.[152]
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The most successful songs during the period 1870–1930 were the so-called Athenian serenades (Αθηναϊκές καντάδες), based on the Heptanesean kantádhes (καντάδες 'serenades'; sing.: καντάδα) and the songs performed on stage (επιθεωρησιακά τραγούδια 'theatrical revue songs') in revues, musical comedies, operettas and nocturnes that were dominating Athens' theatre scene.
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Notable composers of operettas or nocturnes were Kostas Giannidis, Dionysios Lavrangas, Nikos Hatziapostolou, while Theophrastos Sakellaridis' The Godson remains probably the most popular operetta. Despite the fact that the Athenian songs were not autonomous artistic creations (in contrast with the serenades) and despite their original connection with mainly dramatic forms of Art, they eventually became hits as independent songs. Notable actors of Greek operettas, who made also a series of melodies and songs popular at that time, include Orestis Makris, Kalouta sisters, Vasilis Avlonitis, Afroditi Laoutari, Eleni Papadaki, Marika Nezer, Marika Krevata and others. After 1930, wavering among American and European musical influences as well as the Greek musical tradition. Greek composers begin to write music using the tunes of the tango, waltz, swing, foxtrot, some times combined with melodies in the style of Athenian serenades' repertory. Nikos Gounaris was probably the most renowned composer and singer of the time.
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In 1923, after the population exchange between Greece and Turkey, many ethnic Greeks from Asia Minor fled to Athens as a result of the Greco-Turkish War. They settled in poor neighborhoods and brought with them Rebetiko music, making it popular also in Greece, which became later the base for the Laïko music. Other forms of song popular today in Greece are elafrolaika, entechno, dimotika, and skyladika.[153] Greece's most notable, and internationally famous, composers of Greek song, mainly of the entechno form, are Manos Hadjidakis and Mikis Theodorakis. Both composers have achieved fame abroad for their composition of film scores.[153]
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Athens has a long tradition in sports and sporting events, serving as home to the most important clubs in Greek sport and housing a large number of sports facilities. The city has also been host to sports events of international importance.
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Athens has hosted the Summer Olympic Games twice, in 1896 and 2004. The 2004 Summer Olympics required the development of the Athens Olympic Stadium, which has since gained a reputation as one of the most beautiful stadiums in the world, and one of its most interesting modern monuments.[154] The biggest stadium in the country, it hosted two finals of the UEFA Champions League, in 1994 and 2007. Athens' other major stadium, located in the Piraeus area, is the Karaiskakis Stadium, a sports and entertainment complex, host of the 1971 UEFA Cup Winners' Cup Final.
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Athens has hosted the EuroLeague final three times, the first in 1985 and second in 1993, both at the Peace and Friendship Stadium, most known as SEF, a large indoor arena,[155] and the third time in 2007 at the Olympic Indoor Hall. Events in other sports such as athletics, volleyball, water polo etc., have been hosted in the capital's venues.
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Athens is home to three European multi-sport clubs: Olympiacos, Panathinaikos, AEK Athens. In football, Olympiacos have dominated the domestic competitions, Panathinaikos made it to the 1971 European Cup Final, while AEK Athens is the other member of the big three. These clubs also have basketball teams; Panathinaikos and Olympiacos are among the top powers in European basketball, having won the Euroleague six times and three respectively, whilst AEK Athens was the first Greek team to win a European trophy in any team sport.
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Other notable clubs within Athens are Athinaikos, Panionios, Atromitos, Apollon, Panellinios, Ethnikos Piraeus, Maroussi BCE and Peristeri B.C.. Athenian clubs have also had domestic and international success in other sports.
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The Athens area encompasses a variety of terrain, notably hills and mountains rising around the city, and the capital is the only major city in Europe to be bisected by a mountain range. Four mountain ranges extend into city boundaries and thousands of miles of trails criss-cross the city and neighbouring areas, providing exercise and wilderness access on foot and bike.
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Beyond Athens and across the prefecture of Attica, outdoor activities include skiing, rock climbing, hang gliding and windsurfing. Numerous outdoor clubs serve these sports, including the Athens Chapter of the Sierra Club, which leads over 4,000 outings annually in the area.
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Beside the above clubs, inside the boundaries of Athens Municipality there are some more clubs with presence in national divisions or notable action for short periods. Some of them are PAO Rouf (Rouf) with earlier presence in Gamma Ethniki, Petralona F.C.(el) (Petralona), football club founded in 1963, with earlier presence in Beta Ethniki, Attikos F.C.(el) (Kolonos), football club founded in 1919 with short presence in Gamma Ethniki, Athinais Kypselis [es] (Kypseli), football club founded in 1938 with short presence in Gamma Ethniki, Gyziakos (Gyzi), basketball club founded in 1937 with short presence in Beta Ethniki basketball and Aetos B.C.(el) (Agios Panteleimonas), basketball club founded in 1992 with earlier presence in A2 Ethniki Basketball. Another important Athenian sport club is the Athens Tennis Club founded in 1895 with important offer for the Greek tennis.[156]
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The revival of the modern Olympic Games was brought forth in 1896, by Frenchman Pierre de Coubertin. Thanks to his efforts, Athens was awarded the first modern Olympic Games. In 1896, the city had a population of 123,000[93] and the event helped boost the city's international profile. Of the venues used for these Olympics, the Kallimarmaro Stadium, and Zappeion were most crucial. The Kallimarmaro is a replica of the ancient Athenian stadiums, and the only major stadium (in its capacity of 60,000) to be made entirely of white marble from Mount Penteli, the same material used for construction of the Parthenon.
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Fencing before the king of Greece at the 1896 Summer Olympics.
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The Panathenaic Stadium of Athens (Kallimarmaron) dates back to the 4th century BC and has hosted the first modern Olympic Games in 1896.
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The 1906 Summer Olympics, or the 1906 Intercalated games, were held in Athens. The intercalated competitions were intermediate games to the internationally organized Olympics, and were meant to be organized in Greece every four years, between the main Olympics. This idea later lost support from the IOC and these games were discontinued.
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Athens was awarded the 2004 Summer Olympics on 5 September 1997 in Lausanne, Switzerland, after having lost a previous bid to host the 1996 Summer Olympics, to Atlanta, United States.[17] It was to be the second time Athens would host the games, following the inaugural event of 1896. After an unsuccessful bid in 1990, the 1997 bid was radically improved, including an appeal to Greece's Olympic history. In the last round of voting, Athens defeated Rome with 66 votes to 41.[17] Prior to this round, the cities of Buenos Aires, Stockholm and Cape Town had been eliminated from competition, having received fewer votes.[17]
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During the first three years of preparations, the International Olympic Committee had expressed concern over the speed of construction progress for some of the new Olympic venues. In 2000 the Organising Committee's president was replaced by Gianna Angelopoulos-Daskalaki, who was the president of the original Bidding Committee in 1997. From that point forward, preparations continued at a highly accelerated, almost frenzied pace.
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Although the heavy cost was criticized, estimated at $1.5 billion, Athens was transformed into a more functional city that enjoys modern technology both in transportation and in modern urban development.[157] Some of the finest sporting venues in the world were created in the city, all of which were fully ready for the games. The games welcomed over 10,000 athletes from all 202 countries.[157]
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The 2004 Games were judged a success, as both security and organization worked well, and only a few visitors reported minor problems mainly concerning accommodation issues. The 2004 Olympic Games were described as Unforgettable, dream Games, by IOC President Jacques Rogge for their return to the birthplace of the Olympics, and for meeting the challenges of holding the Olympic Games.[157] The only observable problem was a somewhat sparse attendance of some early events. Eventually, however, a total of more than 3.5 million tickets were sold, which was higher than any other Olympics with the exception of Sydney (more than 5 million tickets were sold there in 2000).[158]
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In 2008 it was reported that most of the Olympic venues had fallen into disrepair: according to those reports, 21 of the 22 facilities built for the games had either been left abandoned or are in a state of dereliction, with several squatter camps having sprung up around certain facilities, and a number of venues afflicted by vandalism, graffiti or strewn with rubbish.[159][160] These claims, however, are disputed and likely to be inaccurate, as most of the facilities used for the Athens Olympics are either in use or in the process of being converted for post-Olympics use. The Greek Government has created a corporation, Olympic Properties SA, which is overseeing the post-Olympics management, development and conversion of these facilities, some of which will be sold off (or have already been sold off) to the private sector,[161] while other facilities are still in use just as during the Olympics, or have been converted for commercial use or modified for other sports.[162] Concerts and theatrical shows, such as those by the troupe Cirque du Soleil, have recently been held in the complex.[153]
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An organ is a group of tissues with similar functions. Plant life and animal life rely on many organs that coexist in organ systems.[1]
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A given organ's tissues can be broadly categorized as parenchyma, the tissue peculiar to (or at least archetypal of) the organ and that does the organ's specialized job, and stroma, the tissues with supportive, structural, connective, or ancillary functions. For example, in a gland, the tissue that makes the hormones is the parenchyma, whereas the stroma includes the nerves that innervate the parenchyma, the blood vessels that oxygenate and nourish it and carry away its metabolic wastes, and the connective tissues that provide a suitable place for it to be situated and anchored. The main tissues that make up an organ tend to have common embryologic origins, such as arising from the same germ layer. Functionally related organs often cooperate to form whole organ systems. Organs exist in most multicellular organisms. In single-celled organisms such as bacteria, the functional analogue of an organ is known as an organelle. In plants, there are three main organs.[2] A hollow organ is an internal organ that forms a hollow tube, or pouch such as the stomach, intestine, or bladder.
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In the study of anatomy, the term viscus refers to an internal organ. Viscera is the plural form.[3][4]
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The number of organs in any organism depends on which precise definition of the term one uses. By one widely used definition, 79 organs have been identified in the human body.[5]
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In biology, tissue is a cellular organizational level between cells and complete organs. A tissue is an ensemble of similar cells and their extracellular matrix from the same origin that together carry out a specific function. Organs are then formed by the functional grouping together of multiple tissues.
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The study of human and animal tissues is known as histology or, in connection with disease, histopathology. For plants, the discipline is called plant morphology. Classical tools for studying tissues include the paraffin block (agarose gel is often used with plants[6][7]) in which tissue is embedded and then sectioned, the histological stain, and the optical microscope. In the last couple of decades, developments in electron microscopy, immunofluorescence, and the use of frozen tissue sections have enhanced the detail that can be observed in tissues. With these tools, the classical appearances of tissues can be examined in health and disease, enabling considerable refinement of medical diagnosis and prognosis.
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Two or more organs working together in the execution of a specific body function form an organ system, also called a biological system or body system. The functions of organ systems often share significant overlap. For instance, the nervous and endocrine system both operate via a shared organ, the hypothalamus. For this reason, the two systems are combined and studied as the neuroendocrine system. The same is true for the musculoskeletal system because of the relationship between the muscular and skeletal systems.
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Common organ system designations in plants includes the differentiation of shoot and root. All parts of the plant above ground (in non-epiphytes), including the functionally distinct leaf and flower organs, may be classified together as the shoot organ system.[8]
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Non-placozoan animals such as humans have a variety of organ systems. These specific systems are also widely studied in human anatomy.
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The study of plant organs is referred to as plant morphology, rather than anatomy – as in animal systems. Organs of plants can be divided into vegetative and reproductive. Vegetative plant organs include roots, stems, and leaves. The reproductive organs are variable. In flowering plants, they are represented by the flower, seed and fruit.[9] In conifers, the organ that bears the reproductive structures is called a cone. In other divisions (phyla) of plants, the reproductive organs are called strobili, in Lycopodiophyta, or simply gametophores in mosses.
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The vegetative organs are essential for maintaining the life of a plant. While there can be 11 organ systems in animals, there are far fewer in plants, where some perform the vital functions, such as photosynthesis, while the reproductive organs are essential in reproduction. However, if there is asexual vegetative reproduction, the vegetative organs are those that create the new generation of plants (see clonal colony).
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Many societies have a system for organ donation, in which a living or deceased donor's organ is transplanted into a person with a failing organ. The transplantation of larger solid organs often requires immunosuppression to prevent organ rejection or graft-versus-host disease.
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There is considerable interest throughout the world in creating laboratory-grown or artificial organs.[citation needed]
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The English word "organ" dates back to the twelfth century, in reference to any musical instrument. By the late 14th century, the musical term's meaning had narrowed to refer specifically to the keyboard-based instrument. At the same time, a second meaning arose, in reference to a "body part adapted to a certain function".[10]
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Plant organs are made from tissue composed of different types of tissue. The three tissue types are ground, vascular, and dermal.[11] When three or more organs are present, it is called an organ system.[12].
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The adjective visceral, also splanchnic, is used for anything pertaining to the internal organs. Historically, viscera of animals were examined by Roman pagan priests like the haruspices or the augurs in order to divine the future by their shape, dimensions or other factors.[13] This practice remains an important ritual in some remote, tribal societies.
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The term "visceral" is contrasted with the term "parietal", meaning "of or relating to the wall of a body part, organ or cavity"[14] The two terms are often used in describing a membrane or piece of connective tissue, referring to the opposing sides.[15]
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Aristotle used the word frequently in his philosophy, both to describe the organs of plants or animals (e.g. the roots of a tree, the heart or liver of an animal), and to describe more abstract "parts" of an interconnected whole (e.g. his logical works, taken as a whole, are referred to as the "organon").[16]
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Some alchemists (e.g. Paracelsus) adopted the Hermetic Qabalah assignment between the seven vital organs and the seven classical planets as follows:[17]
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The variations in natural language definitions of what constitutes an organ, their degree of precision, and the variations in how they map to ontologies and taxonomies in information science (for example, to count how many organs exist in a typical human body) are topics explored by writer Carl Engelking of Discover magazine in 2017 as he analyzed the science journalism coverage of the evolving scientific understanding of the mesentery.[18] He explored a challenge now faced by anatomists: as human understanding of ontology generally (that is, how things are defined, and how the relationship of one thing to another is defined) meets applied ontology and ontology engineering, unification of varying views is in higher demand.[18] However, such unification always faces epistemologic frontiers, as humans can only declare computer ontologies with certainty and finality to the extent that their own cognitive taxonomy (that is, science's understanding of the universe) is certain and final. For example, the fact that the tissues of the mesentery are continuous was something that was simply not known for sure until it was demonstrated with microscopy.[19] Because humans cannot predict all future scientific discoveries, they cannot build a unified ontology that is totally certain and will never again change. However, one of the points made by an anatomist interviewed by Engelking is that, finality aside, much more could be done even now to represent existing human knowledge more clearly for computing purposes.
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Beginning in the 20th century[20] transplants began to occur as scientists knew more about the anatomy of organs. These came later in time as procedures were often dangerous and difficult.[21] Both the source and method of obtaining the organ to transplant are major ethical issues to consider, and because organs as resources for transplant are always more limited than demand for them, various notions of justice, including distributive justice, are developed in the ethical analysis. This situation continues as long as transplantation relies upon organ donors rather than technological innovation, testing, and industrial manufacturing.
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The organ level of organisation in animals can be first detected in flatworms and the more derived phyla. The less-advanced taxa (like Placozoa, Sponges and Radiata) do not show consolidation of their tissues into organs.
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More complex animals are composed of different organs, which have been evolving over time. For example, the liver evolved in the stem vertebrates more than 500 million years ago, while the gut and brain are even more ancient, arising in the ancestor of vertebrates, insects, and worms more than 600 million years ago.
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Given the ancient origin of most vertebrate organs, researchers have looked for model systems, where organs have evolved more recently, and ideally have evolved multiple times independently. An outstanding model for this kind of research is the placenta, which has evolved more than 100 times independently in vertebrates, has evolved relatively recently in some lineages, and exists in intermediate forms in extant taxa.[22] Studies on the evolution of the placenta have identified a variety of genetic and physiological processes that contribute to the origin and evolution of organs, these include the re-purposing of existing animal tissues, the acquisition of new functional properties by these tissues, and novel interactions of distinct tissue types.[22]
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Sensation is the physical process during which sensory systems respond to stimuli and provide data for perception.[1] A sense is any of the systems involved in sensation. During sensation, sense organs engage in stimulus collection and transduction.[2] Sensation is often differentiated from the related and dependent concept of perception, which processes and integrates sensory information in order to give meaning to and understand detected stimuli, giving rise to subjective perceptual experience, or qualia.[3] Sensation and perception are central to and precede almost all aspects of cognition, behavior and thought.[1]
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In organisms, a sensory organ consists of a group of related sensory cells that respond to a specific type of physical stimulus. Via cranial and spinal nerves, the different types of sensory receptor cells (mechanoreceptors, photoreceptors, chemoreceptors, thermoreceptors) in sensory organs transduct sensory information from sensory organs towards the central nervous system, to the sensory cortices in the brain, where sensory signals are further processed and interpreted (perceived).[1][4][5] Sensory systems, or senses, are often divided into external (exteroception) and internal (interoception) sensory systems.[6][7] Sensory modalities or submodalities refer to the way sensory information is encoded or transduced.[4] Multimodality integrates different senses into one unified perceptual experience. For example, information from one sense has the potential to influence how information from another is perceived.[2] Sensation and perception are studied by a variety of related fields, most notably psychophysics, neurobiology, cognitive psychology, and cognitive science.[1]
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Humans have a multitude of sensory systems. Human external sensation is based on the sensory organs of the eyes, ears, skin, inner ear, nose, and mouth. The corresponding sensory systems of the visual system (sense of vision), auditory system (sense of hearing), somatosensory system (sense of touch), vestibular system (sense of balance), olfactory system (sense of smell), and gustatory system (sense of taste) contribute, respectively, to the perceptions of vision, hearing, touch, spatial orientation, smell, and taste (flavor).[2][1] Internal sensation, or interoception, detects stimuli from internal organs and tissues. Many internal sensory and perceptual systems exist in humans, including proprioception (body position) and nociception (pain). Further internal chemoreception and osmoreception based sensory systems lead to various perceptions, such as hunger, thirst, suffocation, and nausea, or different involuntary behaviors, such as vomiting.[6][7][8]
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Nonhuman animals experience sensation and perception, with varying levels of similarity to and difference from humans and other animal species. For example, mammals, in general, have a stronger sense of smell than humans. Some animal species lack one or more human sensory system analogues, some have sensory systems that are not found in humans, while others process and interpret the same sensory information in very different ways. For example, some animals are able to detect electrical[9] and magnetic fields,[10] air moisture,[11] or polarized light,[12] while others sense and perceive through alternative systems, such as echolocation.[13][14] Recently, it has been suggested that plants and artificial agents may be able to detect and interpret environmental information in an analogous manner to animals.[15][16][17]
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Sensory modality refers to the way that information is encoded, which is similar to the idea of transduction. The main sensory modalities can be described on the basis of how each is transduced. Listing all the different sensory modalities, which can number as many as 17, involves separating the major senses into more specific categories, or submodalities, of the larger sense. An individual sensory modality represents the sensation of a specific type of stimulus. For example, the general sensation and perception of touch, which is known as somatosensation, can be separated into light pressure, deep pressure, vibration, itch, pain, temperature, or hair movement, while the general sensation and perception of taste can be separated into submodalities of sweet, salty, sour, bitter, spicy. and umami, all of which are based on different chemicals binding to sensory neurons.[4]
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Sensory receptors are the cells or structures that detect sensations. Stimuli in the environment activate specialized receptor cells in the peripheral nervous system. During transduction, physical stimulus is converted into action potential by receptors and transmitted towards the central nervous system for processing.[5] Different types of stimuli are sensed by different types of receptor cells. Receptor cells can be classified into types on the basis of three different criteria: cell type, position, and function. Receptors can be classified structurally on the basis of cell type and their position in relation to stimuli they sense. Receptors can further be classified functionally on the basis of the transduction of stimuli, or how the mechanical stimulus, light, or chemical changed the cell membrane potential.[4]
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One way to classify receptors is based on their location relative to the stimuli. An exteroceptor is a receptor that is located near a stimulus of the external environment, such as the somatosensory receptors that are located in the skin. An interoceptor is one that interprets stimuli from internal organs and tissues, such as the receptors that sense the increase in blood pressure in the aorta or carotid sinus.[4]
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The cells that interpret information about the environment can be either (1) a neuron that has a free nerve ending, with dendrites embedded in tissue that would receive a sensation; (2) a neuron that has an encapsulated ending in which the sensory nerve endings are encapsulated in connective tissue that enhances their sensitivity; or (3) a specialized receptor cell, which has distinct structural components that interpret a specific type of stimulus. The pain and temperature receptors in the dermis of the skin are examples of neurons that have free nerve endings (1). Also located in the dermis of the skin are lamellated corpuscles, neurons with encapsulated nerve endings that respond to pressure and touch (2). The cells in the retina that respond to light stimuli are an example of a specialized receptor (3), a photoreceptor.[4]
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A transmembrane protein receptor is a protein in the cell membrane that mediates a physiological change in a neuron, most often through the opening of ion channels or changes in the cell signaling processes. Transmembrane receptors are activated by chemicals called ligands. For example, a molecule in food can serve as a ligand for taste receptors. Other transmembrane proteins, which are not accurately called receptors, are sensitive to mechanical or thermal changes. Physical changes in these proteins increase ion flow across the membrane, and can generate an action potential or a graded potential in the sensory neurons.[4]
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A third classification of receptors is by how the receptor transduces stimuli into membrane potential changes. Stimuli are of three general types. Some stimuli are ions and macromolecules that affect transmembrane receptor proteins when these chemicals diffuse across the cell membrane. Some stimuli are physical variations in the environment that affect receptor cell membrane potentials. Other stimuli include the electromagnetic radiation from visible light. For humans, the only electromagnetic energy that is perceived by our eyes is visible light. Some other organisms have receptors that humans lack, such as the heat sensors of snakes, the ultraviolet light sensors of bees, or magnetic receptors in migratory birds.[4]
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Receptor cells can be further categorized on the basis of the type of stimuli they transduce. The different types of functional receptor cell types are mechanoreceptors, photoreceptors, chemoreceptors (osmoreceptor), thermoreceptors, and nociceptors. Physical stimuli, such as pressure and vibration, as well as the sensation of sound and body position (balance), are interpreted through a mechanoreceptor. Photoreceptors convert light (visible electromagnetic radiation) into signals. Chemical stimuli can be interpreted by a chemoreceptor that interprets chemical stimuli, such as an object's taste or smell, while osmoreceptors respond to a chemical solute concentrations of body fluids. Nociception (pain) interprets the presence of tissue damage, from sensory information from mechano-, chemo-, and thermoreceptors.[18] Another physical stimulus that has its own type of receptor is temperature, which is sensed through a thermoreceptor that is either sensitive to temperatures above (heat) or below (cold) normal body temperature.[4]
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Each sense organ (eyes or nose, for instance) requires a minimal amount of stimulation in order to detect a stimulus. This minimum amount of stimulus is called the absolute threshold.[2] The absolute threshold is defined as the minimum amount of stimulation necessary for the detection of a stimulus 50% of the time.[1] Absolute threshold is measured by using a method called signal detection. This process involves presenting stimuli of varying intensities to a subject in order to determine the level at which the subject can reliably detect stimulation in a given sense.[2]
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Differential threshold or just noticeable difference (JDS) is the smallest detectable difference between two stimuli, or the smallest difference in stimuli that can be judged to be different from each other.[1] Weber's Law is an empirical law that states that the difference threshold is a constant fraction of the comparison stimulus.[1] According to Weber's Law, bigger stimuli require larger differences to be noticed.[2]
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Magnitude estimation is a psychophysical method in which subjects assign perceived values of given stimuli. The relationship between stimulus intensity and perceptive intensity is described by Steven's power law.[1]
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Signal detection theory quantifies the experience of the subject to the presentation of a stimulus in the presence of noise. There is internal noise and there is external noise when it comes to signal detection. The internal noise originates from static in the nervous system. For example, an individual with closed eyes in a dark room still sees something - a blotchy pattern of grey with intermittent brighter flashes -, this is internal noise. External noise is the result of noise in the environment that can interfere with the detection of the stimulus of interest. Noise is only a problem if the magnitude of the noise is large enough to interfere with signal collection. The nervous system calculates a criterion, or an internal threshold, for the detection of a signal in the presence of noise. If a signal is judged to be above the criterion, thus the signal is differentiated from the noise, the signal is sensed and perceived. Errors in signal detection can potentially lead to false positives and false negatives. The sensory criterion might be shifted based on the importance of the detecting the signal. Shifting of the criterion may influence the likelihood of false positives and false negatives.[1]
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Subjective visual and auditory experiences appear to be similar across humans subjects. The same cannot be said about taste. For example, there is a molecule called propylthiouracil (PROP) that some humans experience as bitter, some as almost tasteless, while others experience it as somewhere between tasteless and bitter. There is a genetic basis for this difference between perception given the same sensory stimulus. This subjective difference in taste perception has implications for individuals' food preferences, and consequently, health.[1]
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When a stimulus is constant and unchanging, perceptual sensory adaptation occurs. During this process, the subject becomes less sensitive to the stimulus.[2]
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Biological auditory (hearing), vestibular and spatial, and visual systems (vision) appear to break down real-world complex stimuli into sine wave components, through the mathematical process called Fourier analysis. Many neurons have a strong preference for certain sine frequency components in contrast to others. The way that simpler sounds and images are encoded during sensation can provide insight into how perception of real-world objects happens.[1]
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Perception occurs when nerves that lead from the sensory organs (e.g. eye) to the brain are stimulated, even if that stimulation is unrelated to the target signal of the sensory organ. For example, in the case of the eye, it does not matter whether light or something else stimulates the optic nerve, that stimulation will results in visual perception, even if there was no visual stimulus to begin with. (To prove this point to yourself (and if you are a human), close your eyes (preferably in a dark room) and press gently on the outside corner of one eye through the eyelid. You will see a visual spot toward the inside of your visual field, near your nose.)[1]
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All stimuli received by the receptors are transduced to an action potential, which is carried along one or more afferent neurons towards a specific area (cortex) of the brain. Just as different nerves are dedicated to sensory and motors tasks, different areas of the brain (cortices) are similarly dedicated to different sensory and perceptual tasks. More complex processing is accomplished across primary cortical regions that spread beyond the primary cortices. Every nerve, sensory or motor, has its own signal transmission speed. For example, nerves in the frog's legs have a 90 ft/s (99 km/h) signal transmission speed, while sensory nerves in humans, transmit sensory information at speeds between 165 ft/s (181 km/h) and 330 ft/s (362 km/h).[1]
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Perceptual experience is often multimodal. Multimodality integrates different senses into one unified perceptual experience. Information from one sense has the potential to influence how information from another is perceived.[2] Multimodal perception is qualitatively different from unimodal perception. There has been a growing body of evidence since the mid-1990s on the neural correlates of multimodal perception.[20]
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Historical inquiries into the underlying mechanisms of sensation and perception have lead early researchers to subscribe to various philosophical interpretations of perception and the mind, including panpsychism, dualism, and materialism. The majority of modern scientists who study sensation and perception take on a materialistic view of the mind.[1]
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Some examples of human absolute thresholds for the 9-21 external senses.[21]
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Humans respond more strongly to multimodal stimuli compared to the sum of each single modality together, an effect called the superadditive effect of multisensory integration.[2] Neurons that respond to both visual and auditory stimuli have been identified in the superior temporal sulcus.[20] Additionally, multimodal “what” and “where” pathways have been proposed for auditory and tactile stimuli.[22]
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External receptors that respond to stimuli from outside the body are called extoreceptors.[23] Human external sensation is based on the sensory organs of the eyes, ears, skin, vestibular system, nose, and mouth, which contribute, respectively, to the sensory perceptions of vision, hearing, touch, spatial orientation, smell, and taste. Smell and taste are both responsible for identifying molecules and thus both are types of chemoreceptors. Both olfaction (smell) and gustation (taste) require the transduction of chemical stimuli into electrical potentials.[2][1]
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The visual system, or sense of sight, is based on the transduction of light stimuli received through the eyes and contributes to visual perception. The visual system detects light on photoreceptors in the retina of each eye that generates electrical nerve impulses for the perception of varying colors and brightness. There are two types of photoreceptors: rods and cones. Rods are very sensitive to light but do not distinguish colors. Cones distinguish colors but are less sensitive to dim light.[4]
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At the molecular level, visual stimuli cause changes in the photopigment molecule that lead to changes in membrane potential of the photoreceptor cell. A single unit of light is called a photon, which is described in physics as a packet of energy with properties of both a particle and a wave. The energy of a photon is represented by its wavelength, with each wavelength of visible light corresponding to a particular color. Visible light is electromagnetic radiation with a wavelength between 380 and 720 nm. Wavelengths of electromagnetic radiation longer than 720 nm fall into the infrared range, whereas wavelengths shorter than 380 nm fall into the ultraviolet range. Light with a wavelength of 380 nm is blue whereas light with a wavelength of 720 nm is dark red. All other colors fall between red and blue at various points along the wavelength scale.[4]
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The three types of cone opsins, being sensitive to different wavelengths of light, provide us with color vision. By comparing the activity of the three different cones, the brain can extract color information from visual stimuli. For example, a bright blue light that has a wavelength of approximately 450 nm would activate the “red” cones minimally, the “green” cones marginally, and the “blue” cones predominantly. The relative activation of the three different cones is calculated by the brain, which perceives the color as blue. However, cones cannot react to low-intensity light, and rods do not sense the color of light. Therefore, our low-light vision is—in essence—in grayscale. In other words, in a dark room, everything appears as a shade of gray. If you think that you can see colors in the dark, it is most likely because your brain knows what color something is and is relying on that memory.[4]
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There is some disagreement as to whether the visual system consists of one, two, or three submodalities. Neuroanatomists generally regard it as two submodalities, given that different receptors are responsible for the perception of color and brightness. Some argue[citation needed] that stereopsis, the perception of depth using both eyes, also constitutes a sense, but it is generally regarded as a cognitive (that is, post-sensory) function of the visual cortex of the brain where patterns and objects in images are recognized and interpreted based on previously learned information. This is called visual memory.
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The inability to see is called blindness. Blindness may result from damage to the eyeball, especially to the retina, damage to the optic nerve that connects each eye to the brain, and/or from stroke (infarcts in the brain). Temporary or permanent blindness can be caused by poisons or medications. People who are blind from degradation or damage to the visual cortex, but still have functional eyes, are actually capable of some level of vision and reaction to visual stimuli but not a conscious perception; this is known as blindsight. People with blindsight are usually not aware that they are reacting to visual sources, and instead just unconsciously adapt their behavior to the stimulus.
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On February 14, 2013 researchers developed a neural implant that gives rats the ability to sense infrared light which for the first time provides living creatures with new abilities, instead of simply replacing or augmenting existing abilities.[24]
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Visual Perception in Psychology
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According to Gestalt Psychology, people perceive the whole of something even if it is not there. The Gestalt’s Law of Organization states that people have seven factors that help to group what is seen into patterns or groups: Common Fate, Similarity, Proximity, Closure, Symmetry, Continuity, and Past Experience.[25]
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The Law of Common fate says that objects are led along the smoothest path. People follow the trend of motion as the lines/dots flow. [26]
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The Law of Similarity refers to the grouping of images or objects that are similar to each other in some aspect. This could be due to shade, colour, size, shape, or other qualities you could distinguish.[27]
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The Law of Proximity states that our minds like to group based on how close objects are to each other. We may see 42 objects in a group, but we can also perceive three groups of two lines with seven objects in each line. [26]
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The Law of Closure is the idea that we as humans still see a full picture even if there are gaps within that picture. There could be gaps or parts missing from a section of a shape, but we would still perceive the shape as whole.[27]
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The Law of Symmetry refers to a person's preference to see symmetry around a central point. An example would be when we use parentheses in writing. We tend to perceive all of the words in the parentheses as one section instead of individual words within the parentheses.[27]
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The Law of Continuity tells us that objects are grouped together by their elements and then perceived as a whole. This usually happens when we see overlapping objects.We will see the overlapping objects with no interruptions.[27]
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The Law of Past Experience refers to the tendency humans have to categorize objects according to past experiences under certain circumstances. If two objects are usually perceived together or within close proximity of each other the Law of Past Experience is usually seen.[26]
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Hearing, or audition, is the transduction of sound waves into a neural signal that is made possible by the structures of the ear. The large, fleshy structure on the lateral aspect of the head is known as the auricle. At the end of the auditory canal is the tympanic membrane, or ear drum, which vibrates after it is struck by sound waves. The auricle, ear canal, and tympanic membrane are often referred to as the external ear. The middle ear consists of a space spanned by three small bones called the ossicles. The three ossicles are the malleus, incus, and stapes, which are Latin names that roughly translate to hammer, anvil, and stirrup. The malleus is attached to the tympanic membrane and articulates with the incus. The incus, in turn, articulates with the stapes. The stapes is then attached to the inner ear, where the sound waves will be transduced into a neural signal. The middle ear is connected to the pharynx through the Eustachian tube, which helps equilibrate air pressure across the tympanic membrane. The tube is normally closed but will pop open when the muscles of the pharynx contract during swallowing or yawning.[4]
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Mechanoreceptors turn motion into electrical nerve pulses, which are located in the inner ear. Since sound is vibration, propagating through a medium such as air, the detection of these vibrations, that is the sense of the hearing, is a mechanical sense because these vibrations are mechanically conducted from the eardrum through a series of tiny bones to hair-like fibers in the inner ear, which detect mechanical motion of the fibers within a range of about 20 to 20,000 hertz,[28] with substantial variation between individuals. Hearing at high frequencies declines with an increase in age. Inability to hear is called deafness or hearing impairment. Sound can also be detected as vibrations conducted through the body by tactition. Lower frequencies that can be heard are detected this way. Some deaf people are able to determine the direction and location of vibrations picked up through the feet.[29]
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Studies pertaining to Audition started to increase in number towards the latter end of the nineteenth century. During this time, many laboratories in the United States began to create new models, diagrams, and instruments that all pertained to the ear. [30]
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There is a branch of Cognitive Psychology dedicated strictly to Audition. They call it Auditory Cognitive Psychology. The main point is to understand why humans are able to use sound in thinking outside of actually saying it. [31]
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Relating to Auditory Cognitive Psychology is Psychoacoustics. Psychoacoustics is more pointed to people interested in music.[32] Haptics, a word used to refer to both taction and kinesthesia, has many parallels with psychoacoustics.[32] Most research around these two are focused on the instrument, the listener, and the player of the instrument. [32]
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Somatosensation is considered a general sense, as opposed to the special senses discussed in this section. Somatosensation is the group of sensory modalities that are associated with touch and interoception. The modalities of somatosensation include pressure, vibration, light touch, tickle, itch, temperature, pain, kinesthesia.[4] Somatosensation, also called tactition (adjectival form: tactile) is a perception resulting from activation of neural receptors, generally in the skin including hair follicles, but also in the tongue, throat, and mucosa. A variety of pressure receptors respond to variations in pressure (firm, brushing, sustained, etc.). The touch sense of itching caused by insect bites or allergies involves special itch-specific neurons in the skin and spinal cord.[33] The loss or impairment of the ability to feel anything touched is called tactile anesthesia. Paresthesia is a sensation of tingling, pricking, or numbness of the skin that may result from nerve damage and may be permanent or temporary.
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Two types of somatosensory signals that are transduced by free nerve endings are pain and temperature. These two modalities use thermoreceptors and nociceptors to transduce temperature and pain stimuli, respectively. Temperature receptors are stimulated when local temperatures differ from body temperature. Some thermoreceptors are sensitive to just cold and others to just heat. Nociception is the sensation of potentially damaging stimuli. Mechanical, chemical, or thermal stimuli beyond a set threshold will elicit painful sensations. Stressed or damaged tissues release chemicals that activate receptor proteins in the nociceptors. For example, the sensation of heat associated with spicy foods involves capsaicin, the active molecule in hot peppers.[4]
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Low frequency vibrations are sensed by mechanoreceptors called Merkel cells, also known as type I cutaneous mechanoreceptors. Merkel cells are located in the stratum basale of the epidermis. Deep pressure and vibration is transduced by lamellated (Pacinian) corpuscles, which are receptors with encapsulated endings found deep in the dermis, or subcutaneous tissue. Light touch is transduced by the encapsulated endings known as tactile (Meissner) corpuscles. Follicles are also wrapped in a plexus of nerve endings known as the hair follicle plexus. These nerve endings detect the movement of hair at the surface of the skin, such as when an insect may be walking along the skin. Stretching of the skin is transduced by stretch receptors known as bulbous corpuscles. Bulbous corpuscles are also known as Ruffini corpuscles, or type II cutaneous mechanoreceptors.[4]
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The heat receptors are sensitive to infrared radiation and can occur in specialized organs, for instance in pit vipers. The thermoceptors in the skin are quite different from the homeostatic thermoceptors in the brain (hypothalamus), which provide feedback on internal body temperature.
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The vestibular sense, or sense of balance (equilibrium), is the sense that contributes to the perception of balance (equilibrium), spatial orientation, direction, or acceleration (equilibrioception). Along with audition, the inner ear is responsible for encoding information about equilibrium. A similar mechanoreceptor—a hair cell with stereocilia—senses head position, head movement, and whether our bodies are in motion. These cells are located within the vestibule of the inner ear. Head position is sensed by the utricle and saccule, whereas head movement is sensed by the semicircular canals. The neural signals generated in the vestibular ganglion are transmitted through the vestibulocochlear nerve to the brain stem and cerebellum.[4]
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The semicircular canals are three ring-like extensions of the vestibule. One is oriented in the horizontal plane, whereas the other two are oriented in the vertical plane. The anterior and posterior vertical canals are oriented at approximately 45 degrees relative to the sagittal plane. The base of each semicircular canal, where it meets with the vestibule, connects to an enlarged region known as the ampulla. The ampulla contains the hair cells that respond to rotational movement, such as turning the head while saying “no.” The stereocilia of these hair cells extend into the cupula, a membrane that attaches to the top of the ampulla. As the head rotates in a plane parallel to the semicircular canal, the fluid lags, deflecting the cupula in the direction opposite to the head movement. The semicircular canals contain several ampullae, with some oriented horizontally and others oriented vertically. By comparing the relative movements of both the horizontal and vertical ampullae, the vestibular system can detect the direction of most head movements within three-dimensional (3D) space.[4]
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The vestibular nerve conducts information from sensory receptors in three ampulla that sense motion of fluid in three semicircular canals caused by three-dimensional rotation of the head. The vestibular nerve also conducts information from the utricle and the saccule, which contain hair-like sensory receptors that bend under the weight of otoliths (which are small crystals of calcium carbonate) that provide the inertia needed to detect head rotation, linear acceleration, and the direction of gravitational force.
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The gustatory system or the sense of taste is the sensory system that is partially responsible for the perception of taste (flavor).[34] A few recognized submodalities exist within taste: sweet, salty, sour, bitter, and umami. Very recent research has suggested that there may also be a sixth taste submodality for fats, or lipids.[4] The sense of taste is often confused with the perception of flavor, which is the results of the multimodal integration of gustatory (taste) and olfactory (smell) sensations.[35]
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Within the structure of the lingual papillae are taste buds that contain specialized gustatory receptor cells for the transduction of taste stimuli. These receptor cells are sensitive to the chemicals contained within foods that are ingested, and they release neurotransmitters based on the amount of the chemical in the food. Neurotransmitters from the gustatory cells can activate sensory neurons in the facial, glossopharyngeal, and vagus cranial nerves.[4]
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Salty and sour taste submodalities are triggered by the cations Na+ and H+, respectively. The other taste modalities result from food molecules binding to a G protein–coupled receptor. A G protein signal transduction system ultimately leads to depolarization of the gustatory cell. The sweet taste is the sensitivity of gustatory cells to the presence of glucose (or sugar substitutes) dissolved in the saliva. Bitter taste is similar to sweet in that food molecules bind to G protein–coupled receptors. The taste known as umami is often referred to as the savory taste. Like sweet and bitter, it is based on the activation of G protein–coupled receptors by a specific molecule.[4]
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Once the gustatory cells are activated by the taste molecules, they release neurotransmitters onto the dendrites of sensory neurons. These neurons are part of the facial and glossopharyngeal cranial nerves, as well as a component within the vagus nerve dedicated to the gag reflex. The facial nerve connects to taste buds in the anterior third of the tongue. The glossopharyngeal nerve connects to taste buds in the posterior two thirds of the tongue. The vagus nerve connects to taste buds in the extreme posterior of the tongue, verging on the pharynx, which are more sensitive to noxious stimuli such as bitterness.[4]
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Flavor depends on odor, texture, and temperature as well as on taste. Humans receive tastes through sensory organs called taste buds, or gustatory calyculi, concentrated on the upper surface of the tongue. Other tastes such as calcium[36][37] and free fatty acids[38] may also be basic tastes but have yet to receive widespread acceptance. The inability to taste is called ageusia.
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There is a rare phenomenon when it comes to the Gustatory sense. It is called Lexical-Gustatory Synesthesia. Lexical-Gustatory Synesthesia is when people can “taste” words. [39] They have reported having flavor sensations they aren’t actually eating. When they read words, hear words, or even imagine words. They have reported not only simple flavors, but textures, complex flavors, and temperatures as well. [40]
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Like the sense of taste, the sense of smell, or the olfactiory system, is also responsive to chemical stimuli.[4] Unlike taste, there are hundreds of olfactory receptors (388 according to one source), each binding to a particular molecular feature. Odor molecules possess a variety of features and, thus, excite specific receptors more or less strongly. This combination of excitatory signals from different receptors makes up what humans perceive as the molecule's smell.[41]
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The olfactory receptor neurons are located in a small region within the superior nasal cavity. This region is referred to as the olfactory epithelium and contains bipolar sensory neurons. Each olfactory sensory neuron has dendrites that extend from the apical surface of the epithelium into the mucus lining the cavity. As airborne molecules are inhaled through the nose, they pass over the olfactory epithelial region and dissolve into the mucus. These odorant molecules bind to proteins that keep them dissolved in the mucus and help transport them to the olfactory dendrites. The odorant–protein complex binds to a receptor protein within the cell membrane of an olfactory dendrite. These receptors are G protein–coupled, and will produce a graded membrane potential in the olfactory neurons.[4]
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In the brain, olfaction is processed by the olfactory cortex. Olfactory receptor neurons in the nose differ from most other neurons in that they die and regenerate on a regular basis. The inability to smell is called anosmia. Some neurons in the nose are specialized to detect pheromones.[42] Loss of the sense of smell can result in food tasting bland. A person with an impaired sense of smell may require additional spice and seasoning levels for food to be tasted. Anosmia may also be related to some presentations of mild depression, because the loss of enjoyment of food may lead to a general sense of despair. The ability of olfactory neurons to replace themselves decreases with age, leading to age-related anosmia. This explains why some elderly people salt their food more than younger people do.[4]
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Causes of Olfactory dysfunction can be caused by age, exposure to toxic chemicals, viral infections, epilepsy, some sort of neurodegenerative disease, head trauma, or as a result of another disorder. [5]
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As studies in olfaction have continued, there has been a positive correlation to its dysfunction or degeneration and early signs of Alzheimers and sporadic Parkinson’s disease. Many patients don’t notice the decline in smell before being tested. In Parkinson’s Disease and Alzheimers, an olfactory deficit is present in 85 to 90% of the early onset cases. [5]There is evidence that the decline of this sense can precede the Alzheimers or Parkinson’s Disease by a couple years. Although the deficit is present in these two diseases, as well as others, it is important to make note that the severity or magnitude vary with every disease. This has brought to light some suggestions that olfactory testing could be used in some cases to aid in differentiating many of the neurodegenerative diseases. [5]
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Those who were born without a sense of smell or have a damaged sense of smell usually complain about 1, or more, of 3 things. Our olfactory sense is also used as a warning against bad food. If the sense of smell is damaged or not there, it can lead to a person contracting food poisoning more often. Not having a sense of smell can also lead to damaged relationships or insecurities within the relationships because of the inability for the person to not smell body odor. Lastly, smell influences how food and drink taste. When the olfactory sense is damaged, the satisfaction from eating and drinking is not as prominent.
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Proprioception, the kinesthetic sense, provides the parietal cortex of the brain with information on the movement and relative positions of the parts of the body. Neurologists test this sense by telling patients to close their eyes and touch their own nose with the tip of a finger. Assuming proper proprioceptive function, at no time will the person lose awareness of where the hand actually is, even though it is not being detected by any of the other senses. Proprioception and touch are related in subtle ways, and their impairment results in surprising and deep deficits in perception and action.[43]
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Nociception (physiological pain) signals nerve-damage or damage to tissue. The three types of pain receptors are cutaneous (skin), somatic (joints and bones), and visceral (body organs). It was previously believed that pain was simply the overloading of pressure receptors, but research in the first half of the 20th century indicated that pain is a distinct phenomenon that intertwines with all of the other senses, including touch. Pain was once considered an entirely subjective experience, but recent studies show that pain is registered in the anterior cingulate gyrus of the brain.[44] The main function of pain is to attract our attention to dangers and motivate us to avoid them. For example, humans avoid touching a sharp needle, or hot object, or extending an arm beyond a safe limit because it is dangerous, and thus hurts. Without pain, people could do many dangerous things without being aware of the dangers.
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An internal sensation and perception also known as interoception[45] is "any sense that is normally stimulated from within the body".[46] These involve numerous sensory receptors in internal organs. Interoception is thought to be atypical in clinical conditions such as alexithymia.[47]
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Some examples of specific receptors are:
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Other living organisms have receptors to sense the world around them, including many of the senses listed above for humans. However, the mechanisms and capabilities vary widely.
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An example of smell in non-mammals is that of sharks, which combine their keen sense of smell with timing to determine the direction of a smell. They follow the nostril that first detected the smell.[54] Insects have olfactory receptors on their antennae. Although it is unknown to the degree and magnitude which non-human animals can smell better than humans.[55]
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Many animals (salamanders, reptiles, mammals) have a vomeronasal organ[56] that is connected with the mouth cavity. In mammals it is mainly used to detect pheromones of marked territory, trails, and sexual state. Reptiles like snakes and monitor lizards make extensive use of it as a smelling organ by transferring scent molecules to the vomeronasal organ with the tips of the forked tongue. In reptiles the vomeronasal organ is commonly referred to as Jacobsons organ. In mammals, it is often associated with a special behavior called flehmen characterized by uplifting of the lips. The organ is vestigial in humans, because associated neurons have not been found that give any sensory input in humans.[57]
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Flies and butterflies have taste organs on their feet, allowing them to taste anything they land on. Catfish have taste organs across their entire bodies, and can taste anything they touch, including chemicals in the water.[58]
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Cats have the ability to see in low light, which is due to muscles surrounding their irides–which contract and expand their pupils–as well as to the tapetum lucidum, a reflective membrane that optimizes the image.
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Pit vipers, pythons and some boas have organs that allow them to detect infrared light, such that these snakes are able to sense the body heat of their prey. The common vampire bat may also have an infrared sensor on its nose.[59] It has been found that birds and some other animals are tetrachromats and have the ability to see in the ultraviolet down to 300 nanometers. Bees and dragonflies[60] are also able to see in the ultraviolet. Mantis shrimps can perceive both polarized light and multispectral images and have twelve distinct kinds of color receptors, unlike humans which have three kinds and most mammals which have two kinds.[61]
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Cephalopods have the ability to change color using chromatophores in their skin. Researchers believe that opsins in the skin can sense different wavelengths of light and help the creatures choose a coloration that camouflages them, in addition to light input from the eyes.[62] Other researchers hypothesize that cephalopod eyes in species which only have a single photoreceptor protein may use chromatic aberration to turn monochromatic vision into color vision,[63] explaining pupils shaped like the letter U, the letter W, or a dumbbell, as well as explaining the need for colorful mating displays.[64] Some cephalopods can distinguish the polarization of light.
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Many invertebrates have a statocyst, which is a sensor for acceleration and orientation that works very differently from the mammalian's semi-circular canals.
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In addition, some animals have senses that humans do not, including the following:
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Magnetoception (or magnetoreception) is the ability to detect the direction one is facing based on the Earth's magnetic field. Directional awareness is most commonly observed in birds, which rely on their magnetic sense to navigate during migration.[65][65][66][permanent dead link][67][68] It has also been observed in insects such as bees. Cattle make use of magnetoception to align themselves in a north–south direction.[69] Magnetotactic bacteria build miniature magnets inside themselves and use them to determine their orientation relative to the Earth's magnetic field.[70][71] There has been some recent (tentative) research suggesting that the Rhodopsin in the human eye, which responds particularly well to blue light, can facilitate magnetoception in humans.[72]
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Certain animals, including bats and cetaceans, have the ability to determine orientation to other objects through interpretation of reflected sound (like sonar). They most often use this to navigate through poor lighting conditions or to identify and track prey. There is currently an uncertainty whether this is simply an extremely developed post-sensory interpretation of auditory perceptions or it actually constitutes a separate sense. Resolution of the issue will require brain scans of animals while they actually perform echolocation, a task that has proven difficult in practice.
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Blind people report they are able to navigate and in some cases identify an object by interpreting reflected sounds (especially their own footsteps), a phenomenon known as human echolocation.
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Electroreception (or electroception) is the ability to detect electric fields. Several species of fish, sharks, and rays have the capacity to sense changes in electric fields in their immediate vicinity. For cartilaginous fish this occurs through a specialized organ called the Ampullae of Lorenzini. Some fish passively sense changing nearby electric fields; some generate their own weak electric fields, and sense the pattern of field potentials over their body surface; and some use these electric field generating and sensing capacities for social communication. The mechanisms by which electroceptive fish construct a spatial representation from very small differences in field potentials involve comparisons of spike latencies from different parts of the fish's body.
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The only orders of mammals that are known to demonstrate electroception are the dolphin and monotreme orders. Among these mammals, the platypus[73] has the most acute sense of electroception.
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A dolphin can detect electric fields in water using electroreceptors in vibrissal crypts arrayed in pairs on its snout and which evolved from whisker motion sensors.[74] These electroreceptors can detect electric fields as weak as 4.6 microvolts per centimeter, such as those generated by contracting muscles and pumping gills of potential prey. This permits the dolphin to locate prey from the seafloor where sediment limits visibility and echolocation.
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Spiders have been shown to detect electric fields to determine a suitable time to extend web for 'ballooning'.[75]
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Body modification enthusiasts have experimented with magnetic implants to attempt to replicate this sense.[76] However, in general humans (and it is presumed other mammals) can detect electric fields only indirectly by detecting the effect they have on hairs. An electrically charged balloon, for instance, will exert a force on human arm hairs, which can be felt through tactition and identified as coming from a static charge (and not from wind or the like). This is not electroreception, as it is a post-sensory cognitive action.
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Hygroreception is the ability to detect changes in the moisture content of the environment.[11][77]
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The ability to sense infrared thermal radiation evolved independently in various families of snakes. Essentially, it allows these reptiles to "see" radiant heat at wavelengths between 5 and 30 μm to a degree of accuracy such that a blind rattlesnake can target vulnerable body parts of the prey at which it strikes.[78] It was previously thought that the organs evolved primarily as prey detectors, but it is now believed that it may also be used in thermoregulatory decision making.[79] The facial pit underwent parallel evolution in pitvipers and some boas and pythons, having evolved once in pitvipers and multiple times in boas and pythons.[80] The electrophysiology of the structure is similar between the two lineages, but they differ in gross structural anatomy. Most superficially, pitvipers possess one large pit organ on either side of the head, between the eye and the nostril (Loreal pit), while boas and pythons have three or more comparatively smaller pits lining the upper and sometimes the lower lip, in or between the scales. Those of the pitvipers are the more advanced, having a suspended sensory membrane as opposed to a simple pit structure. Within the family Viperidae, the pit organ is seen only in the subfamily Crotalinae: the pitvipers. The organ is used extensively to detect and target endothermic prey such as rodents and birds, and it was previously assumed that the organ evolved specifically for that purpose. However, recent evidence shows that the pit organ may also be used for thermoregulation. According to Krochmal et al., pitvipers can use their pits for thermoregulatory decision-making while true vipers (vipers who do not contain heat-sensing pits) cannot.
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In spite of its detection of IR light, the pits' IR detection mechanism is not similar to photoreceptors – while photoreceptors detect light via photochemical reactions, the protein in the pits of snakes is in fact a temperature-sensitive ion channel. It senses infrared signals through a mechanism involving warming of the pit organ, rather than a chemical reaction to light.[81] This is consistent with the thin pit membrane, which allows incoming IR radiation to quickly and precisely warm a given ion channel and trigger a nerve impulse, as well as vascularize the pit membrane in order to rapidly cool the ion channel back to its original "resting" or "inactive" temperature.[81]
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Pressure detection uses the organ of Weber, a system consisting of three appendages of vertebrae transferring changes in shape of the gas bladder to the middle ear. It can be used to regulate the buoyancy of the fish. Fish like the weather fish and other loaches are also known to respond to low pressure areas but they lack a swim bladder.
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Current detection is a detection system of water currents, consisting mostly of vortices, found in the lateral line of fish and aquatic forms of amphibians. The lateral line is also sensitive to low-frequency vibrations. The mechanoreceptors are hair cells, the same mechanoreceptors for vestibular sense and hearing. It is used primarily for navigation, hunting, and schooling. The receptors of the electrical sense are modified hair cells of the lateral line system.
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Polarized light direction/detection is used by bees to orient themselves, especially on cloudy days. Cuttlefish, some beetles, and mantis shrimp can also perceive the polarization of light. Most sighted humans can in fact learn to roughly detect large areas of polarization by an effect called Haidinger's brush, however this is considered an entoptic phenomenon rather than a separate sense.
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Slit sensillae of spiders detect mechanical strain in the exoskeleton, providing information on force and vibrations.
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By using a variety of sense receptors, plants sense light, temperature, humidity, chemical substances, chemical gradients, reorientation, magnetic fields, infections, tissue damage and mechanical pressure. The absence of a nervous system notwithstanding, plants interpret and respond to these stimuli by a variety of hormonal and cell-to-cell communication pathways that result in movement, morphological changes and physiological state alterations at the organism level, that is, result in plant behavior. Such physiological and cognitive functions are generally not believed to give rise to mental phenomena or qualia, however, as these are typically considered the product of nervous system activity. The emergence of mental phenomena from the activity of systems functionally or computationally analogous to that of nervous systems is, however, a hypothetical possibility explored by some schools of thought in the philosophy of mind field, such as functionalism and computationalism.
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However, plants could perceive the world around them,[15] and might be able to emit airborne sounds similar to "screaming" when stressed. Those noises could not be detectable by human ears, but organisms with a hearing range that can hear ultrasonic frequencies—like mice, bats or perhaps other plants—could hear the plants' cries from as far as 15 feet (4.6 m) away.[82]
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Machine perception is the capability of a computer system to interpret data in a manner that is similar to the way humans use their senses to relate to the world around them.[16][17][83] Computers take in and respond to their environment through attached hardware. Until recently, input was limited to a keyboard, joystick or a mouse, but advances in technology, both in hardware and software, have allowed computers to take in sensory input in a way similar to humans.[16][17]
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In the time of William Shakespeare, there were commonly reckoned to be five wits or five senses.[85] At that time, the words "sense" and "wit" were synonyms,[85] so the senses were known as the five outward wits.[86][87] This traditional concept of five senses is common today.
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The traditional five senses are enumerated as the "five material faculties" (pañcannaṃ indriyānaṃ avakanti) in Hindu literature. They appear in allegorical representation as early as in the Katha Upanishad (roughly 6th century BC), as five horses drawing the "chariot" of the body, guided by the mind as "chariot driver".
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Depictions of the five traditional senses as allegory became a popular subject for seventeenth-century artists, especially among Dutch and Flemish Baroque painters. A typical example is Gérard de Lairesse's Allegory of the Five Senses (1668), in which each of the figures in the main group alludes to a sense: Sight is the reclining boy with a convex mirror, hearing is the cupid-like boy with a triangle, smell is represented by the girl with flowers, taste is represented by the woman with the fruit, and touch is represented by the woman holding the bird.
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In Buddhist philosophy, Ayatana or "sense-base" includes the mind as a sense organ, in addition to the traditional five. This addition to the commonly acknowledged senses may arise from the psychological orientation involved in Buddhist thought and practice. The mind considered by itself is seen as the principal gateway to a different spectrum of phenomena that differ from the physical sense data. This way of viewing the human sense system indicates the importance of internal sources of sensation and perception that complements our experience of the external world.[citation needed]
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The reproductive system of an organism, also known as the genital system, is the biological system made up of all the anatomical organs involved in sexual reproduction. Many non-living substances such as fluids, hormones, and pheromones are also important accessories to the reproductive system.[1] Unlike most organ systems, the sexes of differentiated species often have significant differences. These differences allow for a combination of genetic material between two individuals, which allows for the possibility of greater genetic fitness of the offspring.[2]
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In mammals, the major organs of the reproductive system include the external genitalia (penis and vulva) as well as a number of internal organs, including the gamete-producing gonads (testicles and ovaries). Diseases of the human reproductive system are very common and widespread, particularly communicable sexually transmitted diseases.[3]
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Most other vertebrates have generally similar reproductive systems consisting of gonads, ducts, and openings. However, there is a great diversity of physical adaptations as well as reproductive strategies in every group of vertebrates.
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Vertebrates share key elements of their reproductive systems. They all have gamete-producing organs known as gonads. In females, these gonads are then connected by oviducts to an opening to the outside of the body, typically the cloaca, but sometimes to a unique pore such as a vagina or intromittent organ.
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The human reproductive system usually involves internal fertilization by sexual intercourse. During this process, the male inserts his erect penis into the female's vagina and ejaculates semen, which contains sperm. The sperm then travels through the vagina and cervix into the uterus or fallopian tubes for fertilization of the ovum. Upon successful fertilization and implantation, gestation of the fetus then occurs within the female's uterus for approximately nine months, this process is known as pregnancy in humans. Gestation ends with childbirth, delivery following labor. Labor consists of the muscles of the uterus contracting, the cervix dilating, and the baby passing out the vagina (the female genital organ). Human's babies and children are nearly helpless and require high levels of parental care for many years. One important type of parental care is the use of the mammary glands in the female breasts to nurse the baby.[4]
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The female reproductive system has two functions: The first is to produce egg cells, and the second is to protect and nourish the offspring until birth. The male reproductive system has one function, and it is to produce and deposit sperm. Humans have a high level of sexual differentiation. In addition to differences in nearly every reproductive organ, numerous differences typically occur in secondary sexual characteristics.
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The male reproductive system is a series of organs located outside of the body and around the pelvic region of a male that contribute towards the reproduction process. The primary direct function of the male reproductive system is to provide the male sperm for fertilization of the ovum.
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The major reproductive organs of the male can be grouped into three categories. The first category is sperm production and storage. Production takes place in the testes which are housed in the temperature regulating scrotum, immature sperm then travel to the epididymis for development and storage. The second category are the ejaculatory fluid producing glands which include the seminal vesicles, prostate, and the vas deferens. The final category are those used for copulation, and deposition of the spermatozoa (sperm) within the male, these include the penis, urethra, vas deferens, and Cowper's gland.
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Major secondary sexual characteristics includes: larger, more muscular stature, deepened voice, facial and body hair, broad shoulders, and development of an Adam's apple. An important sexual hormone of males is androgen, and particularly testosterone.
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The testes release a hormone that controls the development of sperm. This hormone is also responsible for the development of physical characteristics in men such as facial hair and a deep voice.
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The human female reproductive system is a series of organs primarily located inside of the body and around the pelvic region of a female that contribute towards the reproductive process. The human female reproductive system contains three main parts: the vulva, which leads to the vagina, the vaginal opening, to the uterus; the uterus, which holds the developing fetus; and the ovaries, which produce the female's ova. The breasts are involved during the parenting stage of reproduction, but in most classifications they are not considered to be part of the female reproductive system.
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The vagina meets the outside at the vulva, which also includes the labia, clitoris and urethra; during intercourse this area is lubricated by mucus secreted by the Bartholin's glands. The vagina is attached to the uterus through the cervix, while the uterus is attached to the ovaries via the fallopian tubes. Each ovary contains hundreds of ova (singular ovum).
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Approximately every 28 days, the pituitary gland releases a hormone that stimulates some of the ova to develop and grow. One ovum is released and it passes through the fallopian tube into the uterus. Hormones produced by the ovaries prepare the uterus to receive the ovum. It sita her and awaits the sperm for fertilization to occur. When this does not occur i.e. no sperm for fertilization, the lining of the uterus, called the endometrium, and unfertilized ova are shed each cycle through the process of menstruation. If the ovum is fertilized by sperm, it attaches to the endometrium and the fetus develops.
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Most mammal reproductive systems are similar, however, there are some notable differences between the non-human mammals and humans. For instance, most male mammals have a penis which is stored internally until erect, and most have a penis bone or baculum.[5] Additionally, males of most species do not remain continually sexually fertile as humans do. Like humans, most groups of mammals have descended testicles found within a scrotum, however, others have descended testicles that rest on the ventral body wall, and a few groups of mammals, such as elephants, have undescended testicles found deep within their body cavities near their kidneys.[6]
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The reproductive system of marsupials is unique in that the female has two vaginae, both of which open externally through one orifice but lead to different compartments within the uterus; males usually have a two-pronged penis, which corresponds to the females' two vaginae.[7][8] Marsupials typically develop their offspring in an external pouch containing teats to which their newborn young (joeys) attach themselves for post uterine development. Also, marsupials have a unique prepenial scrotum.[9] The 15mm (5/8 in) long newborn joey instinctively crawls and wriggles the several inches (15 cm), while clinging to fur, on the way to its mother's pouch.
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The uterus and vagina are unique to mammals with no homologue in birds, reptiles, amphibians, or fish.[citation needed] In place of the uterus the other vertebrate groups have an unmodified oviduct leading directly to a cloaca, which is a shared exit-hole for gametes, urine, and feces. Monotremes (i.e. platypus and echidnas), a group of egg-laying mammals, also lack a uterus and vagina, and in that respect have a reproductive system resembling that of a reptile.
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In domestic canines, sexual maturity (puberty) occurs between the ages of 6 to 12 months for both males and females, although this can be delayed until up to two years of age for some large breeds.
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The mare's reproductive system is responsible for controlling gestation, birth, and lactation, as well as her estrous cycle and mating behavior. The stallion's reproductive system is responsible for his sexual behavior and secondary sex characteristics (such as a large crest).
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Male and female birds have a cloaca, an opening through which eggs, sperm, and wastes pass. Intercourse is performed by pressing the lips of the cloacae together, which is sometimes known as intromittent organ which is known as a phallus that is analogous to the mammals' penis. The female lays amniotic eggs in which the young fetus continues to develop after it leaves the female's body. Unlike most vertebrates female birds typically have only one functional ovary and oviduct.[10] As a group, birds, like mammals, are noted for their high level of parental care.
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Reptiles are almost all sexually dimorphic, and exhibit internal fertilization through the cloaca. Some reptiles lay eggs while others are ovoviviparous (animals that deliver live young). Reproductive organs are found within the cloaca of reptiles. Most male reptiles have copulatory organs, which are usually retracted or inverted and stored inside the body. In turtles and crocodilians, the male has a single median penis-like organ, while male snakes and lizards each possess a pair of penis-like organs.
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Most amphibians exhibit external fertilization of eggs, typically within the water, though some amphibians such as caecilians have internal fertilization.[11] All have paired, internal gonads, connected by ducts to the cloaca.
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Fish exhibit a wide range of different reproductive strategies. Most fish, however, are oviparous and exhibit external fertilization. In this process, females use their cloaca to release large quantities of their gametes, called spawn into the water and one or more males release "milt", a white fluid containing many sperm over the unfertilized eggs. Other species of fish are oviparous and have internal fertilization aided by pelvic or anal fins that are modified into an intromittent organ analogous to the human penis.[12] A small portion of fish species are either viviparous or ovoviviparous, and are collectively known as livebearers.[13]
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Fish gonads are typically pairs of either ovaries or testes. Most fish are sexually dimorphic but some species are hermaphroditic or unisexual.[14]
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Invertebrates have an extremely diverse array of reproductive systems, the only commonality may be that they all lay eggs. Also, aside from cephalopods and arthropods, nearly all other invertebrates are hermaphroditic and exhibit external fertilization.
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All cephalopods are sexually dimorphic and reproduce by laying eggs. Most cephalopods have semi-internal fertilization, in which the male places his gametes inside the female's mantle cavity or pallial cavity to fertilize the ova found in the female's single ovary.[15] Likewise, male cephalopods have only a single testicle. In the female of most cephalopods the nidamental glands aid in development of the egg.
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The "penis" in most unshelled male cephalopods (Coleoidea) is a long and muscular end of the gonoduct used to transfer spermatophores to a modified arm called a hectocotylus. That in turn is used to transfer the spermatophores to the female. In species where the hectocotylus is missing, the "penis" is long and able to extend beyond the mantle cavity and transfer the spermatophores directly to the female.
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Most insects reproduce oviparously, i.e. by laying eggs. The eggs are produced by the female in a pair of ovaries. Sperm, produced by the male in one testis or more commonly two, is transmitted to the female during mating by means of external genitalia. The sperm is stored within the female in one or more spermathecae. At the time of fertilization, the eggs travel along oviducts to be fertilized by the sperm and are then expelled from the body ("laid"), in most cases via an ovipositor.
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Arachnids may have one or two gonads, which are located in the abdomen. The genital opening is usually located on the underside of the second abdominal segment. In most species, the male transfers sperm to the female in a package, or spermatophore. Complex courtship rituals have evolved in many arachnids to ensure the safe delivery of the sperm to the female.[16]
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Arachnids usually lay yolky eggs, which hatch into immatures that resemble adults. Scorpions, however, are either ovoviviparous or viviparous, depending on species, and bear live young.
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Among all living organisms, flowers, which are the reproductive structures of angiosperms, are the most varied physically and show a correspondingly great diversity in methods of reproduction.[17] Plants that are not flowering plants (green algae, mosses, liverworts, hornworts, ferns and gymnosperms such as conifers) also have complex interplays between morphological adaptation and environmental factors in their sexual reproduction. The breeding system, or how the sperm from one plant fertilizes the ovum of another, depends on the reproductive morphology, and is the single most important determinant of the genetic structure of nonclonal plant populations. Christian Konrad Sprengel (1793) studied the reproduction of flowering plants and for the first time it was understood that the pollination process involved both biotic and abiotic interactions.
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Fungal reproduction is complex, reflecting the differences in lifestyles and genetic makeup within this diverse kingdom of organisms.[18] It is estimated that a third of all fungi reproduce using more than one method of propagation; for example, reproduction may occur in two well-differentiated stages within the life cycle of a species, the teleomorph and the anamorph.[19] Environmental conditions trigger genetically determined developmental states that lead to the creation of specialized structures for sexual or asexual reproduction. These structures aid reproduction by efficiently dispersing spores or spore-containing propagules.
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The United Nations (UN) is an intergovernmental organization that aims to maintain international peace and security, develop friendly relations among nations, achieve international cooperation, and be a centre for harmonizing the actions of nations.[2] It is the largest, most familiar, most internationally represented and most powerful intergovernmental organization in the world. The UN is headquartered on international territory in New York City, with its other main offices in Geneva, Nairobi, Vienna and The Hague.
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The UN was established after World War II with the aim of preventing future wars, succeeding the ineffective League of Nations.[3] On 25 April 1945, 50 governments met in San Francisco for a conference and started drafting the UN Charter, which was adopted on 25 June 1945 and took effect on 24 October 1945, when the UN began operations. Pursuant to the Charter, the organization's objectives include maintaining international peace and security, protecting human rights, delivering humanitarian aid, promoting sustainable development, and upholding international law.[4] At its founding, the UN had 51 member states; this number grew to 193 in 2011,[5] representing almost all of the world's sovereign states.
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The organization's mission to preserve world peace was complicated in its early decades by the Cold War between the United States and Soviet Union and their respective allies. Its missions have consisted primarily of unarmed military observers and lightly armed troops with primarily monitoring, reporting and confidence-building roles.[6] UN membership grew significantly following widespread decolonization beginning in the 1960s. Since then, 80 former colonies have gained independence, including 11 trust territories that had been monitored by the Trusteeship Council.[7] By the 1970s, the UN's budget for economic and social development programmes far outstripped its spending on peacekeeping. After the end of the Cold War, the UN shifted and expanded its field operations, undertaking a wide variety of complex tasks.[8]
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The UN has six principal organs: the General Assembly; the Security Council; the Economic and Social Council (ECOSOC); the Trusteeship Council; the International Court of Justice; and the UN Secretariat. The UN System includes a multitude of specialized agencies, such as the World Bank Group, the World Health Organization, the World Food Programme, UNESCO, and UNICEF. Additionally, non-governmental organizations may be granted consultative status with ECOSOC and other agencies to participate in the UN's work. The UN's chief administrative officer is the Secretary-General, currently Portuguese politician and diplomat António Guterres, who began his five year-term on 1 January 2017. The organization is financed by assessed and voluntary contributions from its member states.
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The UN, its officers, and its agencies have won many Nobel Peace Prizes, though other evaluations of its effectiveness have been mixed. Some commentators believe the organization to be an important force for peace and human development, while others have called it ineffective, biased, or corrupt.
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In the century prior to the UN's creation, several international treaty organizations such as the International Committee of the Red Cross were formed to ensure protection and assistance for victims of armed conflict and strife.[9]
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In 1914, a political assassination in Sarajevo set off a chain of events that led to the outbreak of World War I. As more and more young men were sent down into the trenches, influential voices in the United States and Britain began calling for the establishment of a permanent international body to maintain peace in the postwar world. President Woodrow Wilson became a vocal advocate of this concept, and in 1918 he included a sketch of the international body in his 14-point proposal to end the war. In November 1918, the Central Powers agreed to an armistice to halt the killing in World War I. Two months later, the Allies met with Germany and Austria-Hungary at Versailles to hammer out formal peace terms. President Wilson wanted peace, but the United Kingdom and France disagreed, forcing harsh war reparations on their former enemies. The League of Nations was approved, and in the summer of 1919 Wilson presented the Treaty of Versailles and the Covenant of the League of Nations to the US Senate for ratification. On 10 January 1920, the League of Nations formally came into being when the Covenant of the League of Nations, ratified by 42 nations in 1919, took effect.[10]
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However, at some point the League became ineffective when it failed to act against the Japanese invasion of Manchuria as in February 1933, 40 nations voted for Japan to withdraw from Manchuria but Japan voted against it and walked out of the League instead of withdrawing from Manchuria.[11] It also failed against the Second Italo-Ethiopian War despite trying to talk to Benito Mussolini as he used the time to send an army to Africa, so the League had a plan for Mussolini to just take a part of Ethiopia, but he ignored the League and invaded Ethiopia, the League tried putting sanctions on Italy, but Italy had already conquered Ethiopia and the League had failed.[12] After Italy conquered Ethiopia, Italy and other nations left the league. But all of them realized that it had failed and they began to re-arm as fast as possible.
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During 1938, Britain and France tried negotiating directly with Hitler but this failed in 1939 when Hitler invaded Czechoslovakia.
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When war broke out in 1939, the League closed down and its headquarters in Geneva remained empty throughout the war.[13] Although the United States never joined the League, the country did support its economic and social missions through the work of private philanthropies and by sending representatives to committees.
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The earliest concrete plan for a new world organization began under the aegis of the U.S. State Department in 1939.[14] The text of the "Declaration by United Nations" was drafted at the White House on 29 December 1941, by President Franklin D. Roosevelt, Prime Minister Winston Churchill, and Roosevelt aide Harry Hopkins. It incorporated Soviet suggestions but left no role for France. "Four Policemen" was coined to refer to four major Allied countries, United States, United Kingdom, Soviet Union, and Republic of China, which emerged in the Declaration by United Nations.[15] Roosevelt first coined the term United Nations to describe the Allied countries.[a] "On New Year's Day 1942, President Roosevelt, Prime Minister Churchill, Maxim Litvinov, of the USSR, and T. V. Soong, of China, signed a short document which later came to be known as the United Nations Declaration, and the next day the representatives of twenty-two other nations added their signatures."[16] The term United Nations was first officially used when 26 governments signed this Declaration. One major change from the Atlantic Charter was the addition of a provision for religious freedom, which Stalin approved after Roosevelt insisted.[17][18] By 1 March 1945, 21 additional states had signed.[19]
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A JOINT DECLARATION BY THE UNITED STATES OF AMERICA, THE UNITED KINGDOM OF GREAT BRITAIN AND NORTHERN IRELAND, THE UNION OF SOVIET SOCIALIST REPUBLICS, CHINA, AUSTRALIA, BELGIUM, CANADA, COSTA RICA, CUBA, CZECHOSLOVAKIA, DOMINICAN REPUBLIC, EL SALVADOR, GREECE, GUATEMALA, HAITI, HONDURAS, INDIA, LUXEMBOURG, NETHERLANDS, NEW ZEALAND, NICARAGUA, NORWAY, PANAMA, POLAND, SOUTH AFRICA, YUGOSLAVIA
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The Governments signatory hereto,
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Having subscribed to a common program of purposes and principles embodied in the Joint Declaration of the President of the United States of America and the Prime Minister of Great Britain dated August 14, 1941, known as the Atlantic Charter,
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Being convinced that complete victory over their enemies is essential to defend life, liberty, independence and religious freedom, and to preserve human rights and justice in their own lands as well as in other lands and that they are now engaged in a common struggle against savage and brutal forces seeking to subjugate the world,
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DECLARE:
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The foregoing declaration may be adhered to by other nations which are, or which may be, rendering material assistance and contributions in the struggle for victory over Hitlerism.
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During the war, "the United Nations" became the official term for the Allies. To join, countries had to sign the Declaration and declare war on the Axis.[20]
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The UN was formulated and negotiated among the delegations from the Allied Big Four (the United States, the United Kingdom, the Soviet Union and China) at the Dumbarton Oaks Conference from 21 September 1944 to 7 October 1944 and they agreed on the aims, structure and functioning of the UN.[21][22][23] After months of planning, the UN Conference on International Organization opened in San Francisco, 25 April 1945, attended by 50 governments and a number of non-governmental organizations involved in drafting the UN Charter.[24][25][26] "The heads of the delegations of the sponsoring countries took turns as chairman of the plenary meetings: Anthony Eden, of Britain, Edward Stettinius, of the United States, T. V. Soong, of China, and Vyacheslav Molotov, of the Soviet Union. At the later meetings, Lord Halifax deputized for Mister Eden, Wellington Koo for T. V. Soong, and Mister Gromyko for Mister Molotov."[27] The UN officially came into existence 24 October 1945, upon ratification of the Charter by the five permanent members of the Security Council—France, the Republic of China, the Soviet Union, the UK and the US—and by a majority of the other 46 signatories.[28]
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The first meetings of the General Assembly, with 51 nations represented,[b] and the Security Council took place in Methodist Central Hall, Westminster, London beginning on 10 January 1946.[28] Debates began at once covering topical issues including the presence of Russian troops in Iranian Azerbaijan, Great Britain's forces in Greece and within days the first veto was cast.[31]
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The General Assembly selected New York City as the site for the headquarters of the UN, construction began on 14 September 1948 and the facility was completed on 9 October 1952. Its site—like UN headquarters buildings in Geneva, Vienna, and Nairobi—is designated as international territory.[32] The Norwegian Foreign Minister, Trygve Lie, was elected as the first UN Secretary-General.[28]
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Though the UN's primary mandate was peacekeeping, the division between the US and USSR often paralysed the organization, generally allowing it to intervene only in conflicts distant from the Cold War.[33] Two notable exceptions were a Security Council resolution on 7 July 1950 authorizing a US-led coalition to repel the North Korean invasion of South Korea, passed in the absence of the USSR,[28][34] and the signing of the Korean Armistice Agreement on 27 July 1953.[35]
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On 29 November 1947, the General Assembly approved a resolution to partition Palestine, approving the creation of the state of Israel.[36] Two years later, Ralph Bunche, a UN official, negotiated an armistice to the resulting conflict.[37] On 7 November 1956, the first UN peacekeeping force was established to end the Suez Crisis;[38] however, the UN was unable to intervene against the USSR's simultaneous invasion of Hungary following that country's revolution.[39]
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On 14 July 1960, the UN established United Nations Operation in the Congo (UNOC), the largest military force of its early decades, to bring order to the breakaway State of Katanga, restoring it to the control of the Democratic Republic of the Congo by 11 May 1964.[40] While traveling to meet rebel leader Moise Tshombe during the conflict, Dag Hammarskjöld, often named as one of the UN's most effective Secretaries-General,[41] died in a plane crash; months later he was posthumously awarded the Nobel Peace Prize.[42] In 1964, Hammarskjöld's successor, U Thant, deployed the UN Peacekeeping Force in Cyprus, which would become one of the UN's longest-running peacekeeping missions.[43]
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With the spread of decolonization in the 1960s, the organization's membership saw an influx of newly independent nations. In 1960 alone, 17 new states joined the UN, 16 of them from Africa.[38] On 25 October 1971, with opposition from the United States, but with the support of many Third World nations, the mainland, communist People's Republic of China was given the Chinese seat on the Security Council in place of the Republic of China that occupied Taiwan; the vote was widely seen as a sign of waning US influence in the organization.[44] Third World nations organized into the Group of 77 coalition under the leadership of Algeria, which briefly became a dominant power at the UN.[45] On 10 November 1975, a bloc comprising the USSR and Third World nations passed a resolution, over the strenuous US and Israeli opposition, declaring Zionism to be racism; the resolution was repealed on 16 December 1991, shortly after the end of the Cold War.[46][47]
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With an increasing Third World presence and the failure of UN mediation in conflicts in the Middle East, Vietnam, and Kashmir, the UN increasingly shifted its attention to its ostensibly secondary goals of economic development and cultural exchange.[48] By the 1970s, the UN budget for social and economic development was far greater than its peacekeeping budget.
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After the Cold War, the UN saw a radical expansion in its peacekeeping duties, taking on more missions in ten years than it had in the previous four decades.[49] Between 1988 and 2000, the number of adopted Security Council resolutions more than doubled, and the peacekeeping budget increased more than tenfold.[50][51][52] The UN negotiated an end to the Salvadoran Civil War, launched a successful peacekeeping mission in Namibia, and oversaw democratic elections in post-apartheid South Africa and post-Khmer Rouge Cambodia.[53] In 1991, the UN authorized a US-led coalition that repulsed the Iraqi invasion of Kuwait.[54] Brian Urquhart, Under-Secretary-General from 1971 to 1985, later described the hopes raised by these successes as a "false renaissance" for the organization, given the more troubled missions that followed.[55]
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Though the UN Charter had been written primarily to prevent aggression by one nation against another, in the early 1990s the UN faced a number of simultaneous, serious crises within nations such as Somalia, Haiti, Mozambique, and the former Yugoslavia.[56] The UN mission in Somalia was widely viewed as a failure after the US withdrawal following casualties in the Battle of Mogadishu, and the UN mission to Bosnia faced "worldwide ridicule" for its indecisive and confused mission in the face of ethnic cleansing.[57] In 1994, the UN Assistance Mission for Rwanda failed to intervene in the Rwandan genocide amid indecision in the Security Council.[58]
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Beginning in the last decades of the Cold War, American and European critics of the UN condemned the organization for perceived mismanagement and corruption.[59] In 1984, US President Ronald Reagan, withdrew his nation's funding from United Nations Educational, Scientific and Cultural Organization (UNESCO) over allegations of mismanagement, followed by the UK and Singapore.[60][61] Boutros Boutros-Ghali, Secretary-General from 1992 to 1996, initiated a reform of the Secretariat, reducing the size of the organization somewhat.[62][63] His successor, Kofi Annan (1997–2006), initiated further management reforms in the face of threats from the US to withhold its UN dues.[63]
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From the late 1990s to the early 2000s, international interventions authorized by the UN took a wider variety of forms. The UN mission in the Sierra Leone Civil War of 1991–2002 was supplemented by British Royal Marines, and the invasion of Afghanistan in 2001 was overseen by NATO.[64] In 2003, the United States invaded Iraq despite failing to pass a UN Security Council resolution for authorization, prompting a new round of questioning of the organization's effectiveness.[65] Under the eighth Secretary-General, Ban Ki-moon, the UN intervened with peacekeepers in crises such as the War in Darfur in Sudan and the Kivu conflict in the Democratic Republic of Congo and sent observers and chemical weapons inspectors to the Syrian Civil War.[66] In 2013, an internal review of UN actions in the final battles of the Sri Lankan Civil War in 2009 concluded that the organization had suffered "systemic failure".[67] In 2010, the organization suffered the worst loss of life in its history, when 101 personnel died in the Haiti earthquake[68]
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The Millennium Summit was held in 2000 to discuss the UN's role in the 21st century.[69] The three day meeting was the largest gathering of world leaders in history, and culminated in the adoption by all member states of the Millennium Development Goals (MDGs), a commitment to achieve international development in areas such as poverty reduction, gender equality, and public health. Progress towards these goals, which were to be met by 2015, was ultimately uneven. The 2005 World Summit reaffirmed the UN's focus on promoting development, peacekeeping, human rights, and global security.[70] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71]
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In addition to addressing global challenges, the UN has sought to improve its accountability and democratic legitimacy by engaging more with civil society and fostering a global constituency.[72] In an effort to enhance transparency, in 2016 the organization held its first public debate between candidates for Secretary-General.[73] On 1 January 2017, Portuguese diplomat António Guterres, who previously served as UN High Commissioner for Refugees, became the ninth Secretary-General. Guterres has highlighted several key goals for his administration, including an emphasis on diplomacy for preventing conflicts, more effective peacekeeping efforts, and streamlining the organization to be more responsive and versatile to global needs.[74]
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The UN system is based on five principal organs: the General Assembly, the Security Council, the Economic and Social Council (ECOSOC), the International Court of Justice and the UN Secretariat.[75] A sixth principal organ, the Trusteeship Council, suspended operations on 1 November 1994, upon the independence of Palau, the last remaining UN trustee territory.[76]
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Four of the five principal organs are located at the main UN Headquarters in New York City.[77] The International Court of Justice is located in The Hague, while other major agencies are based in the UN offices at Geneva,[78] Vienna,[79] and Nairobi.[80] Other UN institutions are located throughout the world. The six official languages of the UN, used in intergovernmental meetings and documents, are Arabic, Chinese, English, French, Russian, and Spanish.[81] On the basis of the Convention on the Privileges and Immunities of the United Nations, the UN and its agencies are immune from the laws of the countries where they operate, safeguarding the UN's impartiality with regard to the host and member countries.[82]
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Below the six organs sit, in the words of the author Linda Fasulo, "an amazing collection of entities and organizations, some of which are actually older than the UN itself and operate with almost complete independence from it".[83] These include specialized agencies, research, and training institutions, programs and funds, and other UN entities.[84]
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The UN obeys the Noblemaire principle, which is binding on any organization that belongs to the UN system. This principle calls for salaries that will draw and keep citizens of countries where salaries are highest, and also calls for equal pay for work of equal value independent of the employee's nationality.[85][86] In practice, the ICSC takes reference to the highest-paying national civil service.[87] Staff salaries are subject to an internal tax that is administered by the UN organizations.[85][88]
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The General Assembly is the main deliberative assembly of the UN. Composed of all UN member states, the assembly meets in regular yearly sessions, but emergency sessions can also be called.[90] The assembly is led by a president, elected from among the member states on a rotating regional basis, and 21 vice-presidents.[91] The first session convened 10 January 1946 in the Methodist Central Hall in London and included representatives of 51 nations.[28]
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When the General Assembly decides on important questions such as those on peace and security, admission of new members and budgetary matters, a two-thirds majority of those present and voting is required.[92][93] All other questions are decided by a majority vote. Each member country has one vote. Apart from the approval of budgetary matters, resolutions are not binding on the members. The Assembly may make recommendations on any matters within the scope of the UN, except matters of peace and security that are under consideration by the Security Council.[90]
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Draft resolutions can be forwarded to the General Assembly by its six main committees:[94]
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As well as by the following two committees:
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The Security Council is charged with maintaining peace and security among countries. While other organs of the UN can only make "recommendations" to member states, the Security Council has the power to make binding decisions that member states have agreed to carry out, under the terms of Charter Article 25.[95] The decisions of the council are known as United Nations Security Council resolutions.[96]
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The Security Council is made up of fifteen member states, consisting of five permanent members—China, France, Russia, the United Kingdom, and the United States—and ten non-permanent members elected for two-year terms by the General Assembly (with end of term date)—Belgium (term ends 2020), Dominican Republic (2020), Estonia (2021), Germany (2020), Indonesia (2020), Niger (2021), Saint Vincent and the Grenadines (2021), South Africa (2020), Tunisia (2021), and Vietnam (2021).[97] The five permanent members hold veto power over UN resolutions, allowing a permanent member to block adoption of a resolution, though not debate. The ten temporary seats are held for two-year terms, with five member states per year voted in by the General Assembly on a regional basis.[98] The presidency of the Security Council rotates alphabetically each month.[99]
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The UN Secretariat is headed by the secretary-general, assisted by the deputy secretary-general and a staff of international civil servants worldwide.[100] It provides studies, information, and facilities needed by UN bodies for their meetings. It also carries out tasks as directed by the Security Council, the General Assembly, the Economic and Social Council, and other UN bodies.[101]
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The secretary-general acts as the de facto spokesperson and leader of the UN. The position is defined in the UN Charter as the organization's "chief administrative officer".[102] Article 99 of the charter states that the secretary-general can bring to the Security Council's attention "any matter which in his opinion may threaten the maintenance of international peace and security", a phrase that Secretaries-General since Trygve Lie have interpreted as giving the position broad scope for action on the world stage.[103] The office has evolved into a dual role of an administrator of the UN organization and a diplomat and mediator addressing disputes between member states and finding consensus to global issues.[104]
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The secretary-general is appointed by the General Assembly, after being recommended by the Security Council, where the permanent members have veto power. There are no specific criteria for the post, but over the years it has become accepted that the post shall be held for one or two terms of five years.[105] The current Secretary-General is António Guterres, who replaced Ban Ki-moon in 2017.
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The International Court of Justice (ICJ), located in The Hague, in the Netherlands, is the primary judicial organ of the UN. Established in 1945 by the UN Charter, the Court began work in 1946 as the successor to the Permanent Court of International Justice. The ICJ is composed of 15 judges who serve 9-year terms and are appointed by the General Assembly; every sitting judge must be from a different nation.[107][108]
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It is based in the Peace Palace in The Hague, sharing the building with the Hague Academy of International Law, a private centre for the study of international law. The ICJ's primary purpose is to adjudicate disputes among states. The court has heard cases related to war crimes, illegal state interference, ethnic cleansing, and other issues.[109] The ICJ can also be called upon by other UN organs to provide advisory opinions.[107]
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It is the only organ that is not located in New York.
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The Economic and Social Council (ECOSOC) assists the General Assembly in promoting international economic and social co-operation and development. ECOSOC has 54 members, who are elected by the General Assembly for a three-year term. The president is elected for a one-year term and chosen amongst the small or middle powers represented on ECOSOC. The council has one annual meeting in July, held in either New York or Geneva. Viewed as separate from the specialized bodies it co-ordinates, ECOSOC's functions include information gathering, advising member nations, and making recommendations.[110][111] Owing to its broad mandate of co-ordinating many agencies, ECOSOC has at times been criticized as unfocused or irrelevant.[110][112]
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ECOSOC's subsidiary bodies include the United Nations Permanent Forum on Indigenous Issues, which advises UN agencies on issues relating to indigenous peoples; the United Nations Forum on Forests, which coordinates and promotes sustainable forest management; the United Nations Statistical Commission, which co-ordinates information-gathering efforts between agencies; and the Commission on Sustainable Development, which co-ordinates efforts between UN agencies and NGOs working towards sustainable development. ECOSOC may also grant consultative status to non-governmental organizations;[110] by 2004, more than 2,200 organizations had received this status.[113]
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The UN Charter stipulates that each primary organ of the United Nations can establish various specialized agencies to fulfil its duties.[114] Some best-known agencies are the International Atomic Energy Agency, the Food and Agriculture Organization, UNESCO (United Nations Educational, Scientific and Cultural Organization), the World Bank, and the World Health Organization (WHO). The UN performs most of its humanitarian work through these agencies. Examples include mass vaccination programmes (through WHO), the avoidance of famine and malnutrition (through the work of the WFP), and the protection of vulnerable and displaced people (for example, by UNHCR).[115]
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With the addition of South Sudan 14 July 2011,[5] there are 193 UN member states, including all undisputed independent states apart from Vatican City.[116][c]
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The UN Charter outlines the rules for membership:
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In addition, there are two non-member observer states of the United Nations General Assembly: the Holy See (which holds sovereignty over Vatican City) and the State of Palestine.[118] The Cook Islands and Niue, both states in free association with New Zealand, are full members of several UN specialized agencies and have had their "full treaty-making capacity" recognized by the Secretariat.[119]
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The Group of 77 (G77) at the UN is a loose coalition of developing nations, designed to promote its members' collective economic interests and create an enhanced joint negotiating capacity in the UN. Seventy-seven nations founded the organization, but by November 2013 the organization had since expanded to 133 member countries.[120] The group was founded 15 June 1964 by the "Joint Declaration of the Seventy-Seven Countries" issued at the United Nations Conference on Trade and Development (UNCTAD). The group held its first major meeting in Algiers in 1967, where it adopted the Charter of Algiers and established the basis for permanent institutional structures.[121] With the adoption of the New International Economic Order by developing countries in the 1970s, the work of the G77 spread throughout the UN system.
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The UN, after approval by the Security Council, sends peacekeepers to regions where armed conflict has recently ceased or paused to enforce the terms of peace agreements and to discourage combatants from resuming hostilities. Since the UN does not maintain its own military, peacekeeping forces are voluntarily provided by member states. These soldiers are sometimes nicknamed "Blue Helmets" for their distinctive gear.[122][123] The peacekeeping force as a whole received the Nobel Peace Prize in 1988.[124]
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In September 2013, the UN had peacekeeping soldiers deployed on 15 missions. The largest was the United Nations Organization Stabilization Mission in the Democratic Republic of the Congo (MONUSCO), which included 20,688 uniformed personnel. The smallest, United Nations Military Observer Group in India and Pakistan (UNMOGIP), included 42 uniformed personnel responsible for monitoring the ceasefire in Jammu and Kashmir. UN peacekeepers with the United Nations Truce Supervision Organization (UNTSO) have been stationed in the Middle East since 1948, the longest-running active peacekeeping mission.[125]
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A study by the RAND Corporation in 2005 found the UN to be successful in two out of three peacekeeping efforts. It compared efforts at nation-building by the UN to those of the United States, and found that seven out of eight UN cases are at peace, as compared with four out of eight US cases at peace.[126] Also in 2005, the Human Security Report documented a decline in the number of wars, genocides, and human rights abuses since the end of the Cold War, and presented evidence, albeit circumstantial, that international activism—mostly spearheaded by the UN—has been the main cause of the decline in armed conflict in that period.[127] Situations in which the UN has not only acted to keep the peace but also intervened include the Korean War (1950–53) and the authorization of intervention in Iraq after the Gulf War (1990–91).[128]
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The UN has also drawn criticism for perceived failures. In many cases, member states have shown reluctance to achieve or enforce Security Council resolutions. Disagreements in the Security Council about military action and intervention are seen as having failed to prevent the Bangladesh genocide in 1971,[129] the Cambodian genocide in the 1970s,[130] and the Rwandan genocide in 1994.[131] Similarly, UN inaction is blamed for failing to either prevent the Srebrenica massacre in 1995 or complete the peacekeeping operations in 1992–93 during the Somali Civil War.[132] UN peacekeepers have also been accused of child rape, soliciting prostitutes, and sexual abuse during various peacekeeping missions in the Democratic Republic of the Congo,[133] Haiti,[134] Liberia,[135] Sudan and what is now South Sudan,[136] Burundi, and Ivory Coast.[137] Scientists cited UN peacekeepers from Nepal as the likely source of the 2010–13 Haiti cholera outbreak, which killed more than 8,000 Haitians following the 2010 Haiti earthquake.[138]
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In addition to peacekeeping, the UN is also active in encouraging disarmament. Regulation of armaments was included in the writing of the UN Charter in 1945 and was envisioned as a way of limiting the use of human and economic resources for their creation.[95] The advent of nuclear weapons came only weeks after the signing of the charter, resulting in the first resolution of the first General Assembly meeting calling for specific proposals for "the elimination from national armaments of atomic weapons and of all other major weapons adaptable to mass destruction".[139] The UN has been involved with arms-limitation treaties, such as the Outer Space Treaty (1967), the Treaty on the Non-Proliferation of Nuclear Weapons (1968), the Seabed Arms Control Treaty (1971), the Biological Weapons Convention (1972), the Chemical Weapons Convention (1992), and the Ottawa Treaty (1997), which prohibits landmines.[140] Three UN bodies oversee arms proliferation issues: the International Atomic Energy Agency, the Organisation for the Prohibition of Chemical Weapons, and the Comprehensive Nuclear-Test-Ban Treaty Organization Preparatory Commission.[141]
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One of the UN's primary purposes is "promoting and encouraging respect for human rights and for fundamental freedoms for all without distinction as to race, sex, language, or religion", and member states pledge to undertake "joint and separate action" to protect these rights.[114][142]
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In 1948, the General Assembly adopted a Universal Declaration of Human Rights, drafted by a committee headed by American diplomat and activist Eleanor Roosevelt, and including the French lawyer René Cassin. The document proclaims basic civil, political, and economic rights common to all human beings, though its effectiveness towards achieving these ends has been disputed since its drafting.[143] The Declaration serves as a "common standard of achievement for all peoples and all nations" rather than a legally binding document, but it has become the basis of two binding treaties, the 1966 International Covenant on Civil and Political Rights and International Covenant on Economic, Social and Cultural Rights.[144] In practice, the UN is unable to take significant action against human rights abuses without a Security Council resolution, though it does substantial work in investigating and reporting abuses.[145]
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In 1979, the General Assembly adopted the Convention on the Elimination of All Forms of Discrimination against Women, followed by the Convention on the Rights of the Child in 1989.[146] With the end of the Cold War, the push for human rights action took on new impetus.[147] The United Nations Commission on Human Rights was formed in 1993 to oversee human rights issues for the UN, following the recommendation of that year's World Conference on Human Rights. Jacques Fomerand, a scholar of the UN, describes this organization's mandate as "broad and vague", with only "meagre" resources to carry it out.[148] In 2006, it was replaced by a Human Rights Council consisting of 47 nations.[149] Also in 2006, the General Assembly passed a Declaration on the Rights of Indigenous Peoples,[150] and in 2011 it passed its first resolution recognizing the rights of LGBT people.[151]
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Other UN bodies responsible for women's rights issues include United Nations Commission on the Status of Women, a commission of ECOSOC founded in 1946; the United Nations Development Fund for Women, created in 1976; and the United Nations International Research and Training Institute for the Advancement of Women, founded in 1979.[152] The UN Permanent Forum on Indigenous Issues, one of three bodies with a mandate to oversee issues related to indigenous peoples, held its first session in 2002.[153]
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Millennium Development Goals[154]
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Another primary purpose of the UN is "to achieve international cooperation in solving international problems of an economic, social, cultural, or humanitarian character".[142] Numerous bodies have been created to work towards this goal, primarily under the authority of the General Assembly and ECOSOC.[155] In 2000, the 192 UN member states agreed to achieve eight Millennium Development Goals by 2015.[156] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71] The SDGs have an associated financing framework called the Addis Ababa Action Agenda.
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The UN Development Programme (UNDP), an organization for grant-based technical assistance founded in 1945, is one of the leading bodies in the field of international development. The organization also publishes the UN Human Development Index, a comparative measure ranking countries by poverty, literacy, education, life expectancy, and other factors.[157][158] The Food and Agriculture Organization (FAO), also founded in 1945, promotes agricultural development and food security.[159] UNICEF (the United Nations Children's Fund) was created in 1946 to aid European children after the Second World War and expanded its mission to provide aid around the world and to uphold the convention on the Rights of the Child.[160][161]
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The World Bank Group and International Monetary Fund (IMF) are independent, specialized agencies and observers within the UN framework, according to a 1947 agreement. They were initially formed separately from the UN through the Bretton Woods Agreement in 1944.[162] The World Bank provides loans for international development, while the IMF promotes international economic co-operation and gives emergency loans to indebted countries.[163]
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The World Health Organization (WHO), which focuses on international health issues and disease eradication, is another of the UN's largest agencies. In 1980, the agency announced that the eradication of smallpox had been completed. In subsequent decades, WHO largely eradicated polio, river blindness, and leprosy.[164] The Joint United Nations Programme on HIV/AIDS (UNAIDS), begun in 1996, co-ordinates the organization's response to the AIDS epidemic.[165] The UN Population Fund, which also dedicates part of its resources to combating HIV, is the world's largest source of funding for reproductive health and family planning services.[166]
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Along with the International Red Cross and Red Crescent Movement, the UN often takes a leading role in co-ordinating emergency relief.[167] The World Food Programme (WFP), created in 1961, provides food aid in response to famine, natural disasters, and armed conflict. The organization reports that it feeds an average of 90 million people in 80 nations each year.[167][168] The Office of the United Nations High Commissioner for Refugees (UNHCR), established in 1950, works to protect the rights of refugees, asylum seekers, and stateless people.[169] UNHCR and WFP programmes are funded by voluntary contributions from governments, corporations, and individuals, though the UNHCR's administrative costs are paid for by the UN's primary budget.[170]
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Since the UN's creation, over 80 colonies have attained independence. The General Assembly adopted the Declaration on the Granting of Independence to Colonial Countries and Peoples in 1960 with no votes against but abstentions from all major colonial powers. The UN works towards decolonization through groups including the UN Committee on Decolonization, created in 1962.[171] The committee lists seventeen remaining "Non-Self-Governing Territories", the largest and most populous of which is Western Sahara.[172]
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Beginning with the formation of the UN Environmental Programme (UNEP) in 1972, the UN has made environmental issues a prominent part of its agenda. A lack of success in the first two decades of UN work in this area led to the 1992 Earth Summit in Rio de Janeiro, Brazil, which sought to give new impetus to these efforts.[173] In 1988, the UNEP and the World Meteorological Organization (WMO), another UN organization, established the Intergovernmental Panel on Climate Change, which assesses and reports on research on global warming.[174] The UN-sponsored Kyoto Protocol, signed in 1997, set legally binding emissions reduction targets for ratifying states.[175]
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The UN also declares and co-ordinates international observances, periods of time to observe issues of international interest or concern. Examples include World Tuberculosis Day, Earth Day, and the International Year of Deserts and Desertification.[176]
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The UN is financed from assessed and voluntary contributions from member states. The General Assembly approves the regular budget and determines the assessment for each member. This is broadly based on the relative capacity of each country to pay, as measured by its gross national income (GNI), with adjustments for external debt and low per capita income.[178] The two-year budget for 2012–13 was $5.512 billion in total.[179]
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The Assembly has established the principle that the UN should not be unduly dependent on any one member to finance its operations. Thus, there is a "ceiling" rate, setting the maximum amount that any member can be assessed for the regular budget. In December 2000, the Assembly revised the scale of assessments in response to pressure from the United States. As part of that revision, the regular budget ceiling was reduced from 25% to 22%.[180] For the least developed countries (LDCs), a ceiling rate of 0.01% is applied.[178] In addition to the ceiling rates, the minimum amount assessed to any member nation (or "floor" rate) is set at 0.001% of the UN budget ($55,120 for the two year budget 2013–2014).[181]
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A large share of the UN's expenditure addresses its core mission of peace and security, and this budget is assessed separately from the main organizational budget.[182] The peacekeeping budget for the 2015–16 fiscal year was $8.27 billion, supporting 82,318 troops deployed in 15 missions around the world.[125] UN peace operations are funded by assessments, using a formula derived from the regular funding scale that includes a weighted surcharge for the five permanent Security Council members, who must approve all peacekeeping operations. This surcharge serves to offset discounted peacekeeping assessment rates for less developed countries.
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the largest contributors for the UN peacekeeping financial operations for the period 2019–2021 are : the United States 27.89% China 15.21%, Japan 8.56%, Germany 6.09% , the United Kingdom 5.78%, France 5.61%, Italy 3.30% and the Russian Federation 3.04%.[183]
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Special UN programmes not included in the regular budget, such as UNICEF and the World Food Programme, are financed by voluntary contributions from member governments, corporations, and private individuals.[184][185]
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A number of agencies and individuals associated with the UN have won the Nobel Peace Prize in recognition of their work. Two Secretaries-General, Dag Hammarskjöld, and Kofi Annan, were each awarded the prize (in 1961 and 2001, respectively), as were Ralph Bunche (1950), a UN negotiator, René Cassin (1968), a contributor to the Universal Declaration of Human Rights, and the US Secretary of State Cordell Hull (1945), the latter for his role in the organization's founding. Lester B. Pearson, the Canadian Secretary of State for External Affairs, was awarded the prize in 1957 for his role in organizing the UN's first peacekeeping force to resolve the Suez Crisis. UNICEF won the prize in 1965, the International Labour Organization in 1969, the UN Peace-Keeping Forces in 1988, the International Atomic Energy Agency (which reports to the UN) in 2005, and the UN-supported Organisation for the Prohibition of Chemical Weapons in 2013. The UN High Commissioner for Refugees was awarded in 1954 and 1981, becoming one of only two recipients to win the prize twice. The UN as a whole was awarded the prize in 2001, sharing it with Annan.[186] In 2007, IPCC received the prize "for their efforts to build up and disseminate greater knowledge about man-made climate change, and to lay the foundations for the measures that are needed to counteract such change."[187]
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Since its founding, there have been many calls for reform of the UN but little consensus on how to do so. Some want the UN to play a greater or more effective role in world affairs, while others want its role reduced to humanitarian work. There have also been numerous calls for the UN Security Council's membership to be increased, for different ways of electing the UN's Secretary-General, and for a UN Parliamentary Assembly. Jacques Fomerand states the most enduring divide in views of the UN is "the North–South split" between richer Northern nations and developing Southern nations. Southern nations tend to favour a more empowered UN with a stronger General Assembly, allowing them a greater voice in world affairs, while Northern nations prefer an economically laissez-faire UN that focuses on transnational threats such as terrorism.[188]
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After World War II, the French Committee of National Liberation was late to be recognized by the US as the government of France, and so the country was initially excluded from the conferences that created the new organization. The future French president Charles de Gaulle criticized the UN, famously calling it a machin ("contraption"), and was not convinced that a global security alliance would help maintain world peace, preferring direct defence treaties between countries.[189] Throughout the Cold War, both the US and USSR repeatedly accused the UN of favouring the other. In 1953, the USSR effectively forced the resignation of Trygve Lie, the Secretary-General, through its refusal to deal with him, while in the 1950s and 1960s, a popular US bumper sticker read, "You can't spell communism without U.N."[190] In a sometimes-misquoted statement, President George W. Bush stated in February 2003 (referring to UN uncertainty towards Iraqi provocations under the Saddam Hussein regime) that "free nations will not allow the UN to fade into history as an ineffective, irrelevant debating society."[191][192][193] In contrast, the French President, François Hollande, stated in 2012 that "France trusts the United Nations. She knows that no state, no matter how powerful, can solve urgent problems, fight for development and bring an end to all crises ... France wants the UN to be the centre of global governance."[194] Critics such as Dore Gold, an Israeli diplomat, Robert S. Wistrich, a British scholar, Alan Dershowitz, an American legal scholar, Mark Dreyfus, an Australian politician, and the Anti-Defamation League consider UN attention to Israel's treatment of Palestinians to be excessive.[195] In September 2015, Saudi Arabia's Faisal bin Hassan Trad has been elected Chair of the UN Human Rights Council panel that appoints independent experts,[196] a move criticized by human rights groups.[197][198]
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Since 1971, the Republic of China on Taiwan has been excluded from the UN and since then has always been rejected in new applications. Taiwanese citizens are also not allowed to enter the buildings of the United Nations with ROC passports. In this way, critics agree that the UN is failing its own development goals and guidelines. This criticism also brought pressure from the People's Republic of China, which regards the territories administered by the ROC as their own territory.[199][200]
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Critics have also accused the UN of bureaucratic inefficiency, waste, and corruption. In 1976, the General Assembly established the Joint Inspection Unit to seek out inefficiencies within the UN system. During the 1990s, the US withheld dues citing inefficiency and only started repayment on the condition that a major reforms initiative be introduced. In 1994, the Office of Internal Oversight Services (OIOS) was established by the General Assembly to serve as an efficiency watchdog.[201] In 1994, former Special Representative of the Secretary-General of the UN to Somalia Mohamed Sahnoun published "Somalia: The Missed Opportunities",[202] a book in which he analyses the reasons for the failure of the 1992 UN intervention in Somalia, showing that, between the start of the Somali civil war in 1988 and the fall of the Siad Barre regime in January 1991, the UN missed at least three opportunities to prevent major human tragedies; when the UN tried to provide humanitarian assistance, they were totally outperformed by NGOs, whose competence and dedication sharply contrasted with the UN's excessive caution and bureaucratic inefficiencies. If radical reform were not undertaken, warned Mohamed Sahnoun, then the UN would continue to respond to such crises with inept improvisation.[203] In 2004, the UN faced accusations that its recently ended Oil-for-Food Programme — in which Iraq had been allowed to trade oil for basic needs to relieve the pressure of sanctions — had suffered from widespread corruption, including billions of dollars of kickbacks. An independent inquiry created by the UN found that many of its officials had been involved, as well as raising "significant" questions about the role of Kojo Annan, the son of Kofi Annan.[204]
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In evaluating the UN as a whole, Jacques Fomerand writes that the "accomplishments of the United Nations in the last 60 years are impressive in their own terms. Progress in human development during the 20th century has been dramatic, and the UN and its agencies have certainly helped the world become a more hospitable and livable place for millions."[205] Evaluating the first 50 years of the UN's history, the author Stanley Meisler writes that "the United Nations never fulfilled the hopes of its founders, but it accomplished a great deal nevertheless", citing its role in decolonization and its many successful peacekeeping efforts.[206] The British historian Paul Kennedy states that while the organization has suffered some major setbacks, "when all its aspects are considered, the UN has brought great benefits to our generation and ... will bring benefits to our children's and grandchildren's generations as well."[207]
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Core features of the UN apparatus, such as the veto privileges of some nations in the Security Council, are often described as fundamentally undemocratic, contrary to the UN mission, and as a main cause of inaction on genocides and crimes against humanity.[208][209]
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The United Nations has inspired the extracurricular activity Model United Nations (MUN). MUN is a simulation of United Nations activity based on the UN agenda and following UN procedure. MUN is usually attended by high school and university students who organize conferences to simulate the various UN committees to discuss important issues of the day.[210] Today Model United Nations educates tens of thousands on United Nations activity around the world. Model United Nations has many famous and notable alumni, such as former Secretary-General of the United Nations Ban Ki-moon.[211]
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The United Nations (UN) is an intergovernmental organization that aims to maintain international peace and security, develop friendly relations among nations, achieve international cooperation, and be a centre for harmonizing the actions of nations.[2] It is the largest, most familiar, most internationally represented and most powerful intergovernmental organization in the world. The UN is headquartered on international territory in New York City, with its other main offices in Geneva, Nairobi, Vienna and The Hague.
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The UN was established after World War II with the aim of preventing future wars, succeeding the ineffective League of Nations.[3] On 25 April 1945, 50 governments met in San Francisco for a conference and started drafting the UN Charter, which was adopted on 25 June 1945 and took effect on 24 October 1945, when the UN began operations. Pursuant to the Charter, the organization's objectives include maintaining international peace and security, protecting human rights, delivering humanitarian aid, promoting sustainable development, and upholding international law.[4] At its founding, the UN had 51 member states; this number grew to 193 in 2011,[5] representing almost all of the world's sovereign states.
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The organization's mission to preserve world peace was complicated in its early decades by the Cold War between the United States and Soviet Union and their respective allies. Its missions have consisted primarily of unarmed military observers and lightly armed troops with primarily monitoring, reporting and confidence-building roles.[6] UN membership grew significantly following widespread decolonization beginning in the 1960s. Since then, 80 former colonies have gained independence, including 11 trust territories that had been monitored by the Trusteeship Council.[7] By the 1970s, the UN's budget for economic and social development programmes far outstripped its spending on peacekeeping. After the end of the Cold War, the UN shifted and expanded its field operations, undertaking a wide variety of complex tasks.[8]
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The UN has six principal organs: the General Assembly; the Security Council; the Economic and Social Council (ECOSOC); the Trusteeship Council; the International Court of Justice; and the UN Secretariat. The UN System includes a multitude of specialized agencies, such as the World Bank Group, the World Health Organization, the World Food Programme, UNESCO, and UNICEF. Additionally, non-governmental organizations may be granted consultative status with ECOSOC and other agencies to participate in the UN's work. The UN's chief administrative officer is the Secretary-General, currently Portuguese politician and diplomat António Guterres, who began his five year-term on 1 January 2017. The organization is financed by assessed and voluntary contributions from its member states.
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The UN, its officers, and its agencies have won many Nobel Peace Prizes, though other evaluations of its effectiveness have been mixed. Some commentators believe the organization to be an important force for peace and human development, while others have called it ineffective, biased, or corrupt.
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In the century prior to the UN's creation, several international treaty organizations such as the International Committee of the Red Cross were formed to ensure protection and assistance for victims of armed conflict and strife.[9]
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In 1914, a political assassination in Sarajevo set off a chain of events that led to the outbreak of World War I. As more and more young men were sent down into the trenches, influential voices in the United States and Britain began calling for the establishment of a permanent international body to maintain peace in the postwar world. President Woodrow Wilson became a vocal advocate of this concept, and in 1918 he included a sketch of the international body in his 14-point proposal to end the war. In November 1918, the Central Powers agreed to an armistice to halt the killing in World War I. Two months later, the Allies met with Germany and Austria-Hungary at Versailles to hammer out formal peace terms. President Wilson wanted peace, but the United Kingdom and France disagreed, forcing harsh war reparations on their former enemies. The League of Nations was approved, and in the summer of 1919 Wilson presented the Treaty of Versailles and the Covenant of the League of Nations to the US Senate for ratification. On 10 January 1920, the League of Nations formally came into being when the Covenant of the League of Nations, ratified by 42 nations in 1919, took effect.[10]
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However, at some point the League became ineffective when it failed to act against the Japanese invasion of Manchuria as in February 1933, 40 nations voted for Japan to withdraw from Manchuria but Japan voted against it and walked out of the League instead of withdrawing from Manchuria.[11] It also failed against the Second Italo-Ethiopian War despite trying to talk to Benito Mussolini as he used the time to send an army to Africa, so the League had a plan for Mussolini to just take a part of Ethiopia, but he ignored the League and invaded Ethiopia, the League tried putting sanctions on Italy, but Italy had already conquered Ethiopia and the League had failed.[12] After Italy conquered Ethiopia, Italy and other nations left the league. But all of them realized that it had failed and they began to re-arm as fast as possible.
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During 1938, Britain and France tried negotiating directly with Hitler but this failed in 1939 when Hitler invaded Czechoslovakia.
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When war broke out in 1939, the League closed down and its headquarters in Geneva remained empty throughout the war.[13] Although the United States never joined the League, the country did support its economic and social missions through the work of private philanthropies and by sending representatives to committees.
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The earliest concrete plan for a new world organization began under the aegis of the U.S. State Department in 1939.[14] The text of the "Declaration by United Nations" was drafted at the White House on 29 December 1941, by President Franklin D. Roosevelt, Prime Minister Winston Churchill, and Roosevelt aide Harry Hopkins. It incorporated Soviet suggestions but left no role for France. "Four Policemen" was coined to refer to four major Allied countries, United States, United Kingdom, Soviet Union, and Republic of China, which emerged in the Declaration by United Nations.[15] Roosevelt first coined the term United Nations to describe the Allied countries.[a] "On New Year's Day 1942, President Roosevelt, Prime Minister Churchill, Maxim Litvinov, of the USSR, and T. V. Soong, of China, signed a short document which later came to be known as the United Nations Declaration, and the next day the representatives of twenty-two other nations added their signatures."[16] The term United Nations was first officially used when 26 governments signed this Declaration. One major change from the Atlantic Charter was the addition of a provision for religious freedom, which Stalin approved after Roosevelt insisted.[17][18] By 1 March 1945, 21 additional states had signed.[19]
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A JOINT DECLARATION BY THE UNITED STATES OF AMERICA, THE UNITED KINGDOM OF GREAT BRITAIN AND NORTHERN IRELAND, THE UNION OF SOVIET SOCIALIST REPUBLICS, CHINA, AUSTRALIA, BELGIUM, CANADA, COSTA RICA, CUBA, CZECHOSLOVAKIA, DOMINICAN REPUBLIC, EL SALVADOR, GREECE, GUATEMALA, HAITI, HONDURAS, INDIA, LUXEMBOURG, NETHERLANDS, NEW ZEALAND, NICARAGUA, NORWAY, PANAMA, POLAND, SOUTH AFRICA, YUGOSLAVIA
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The Governments signatory hereto,
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Having subscribed to a common program of purposes and principles embodied in the Joint Declaration of the President of the United States of America and the Prime Minister of Great Britain dated August 14, 1941, known as the Atlantic Charter,
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Being convinced that complete victory over their enemies is essential to defend life, liberty, independence and religious freedom, and to preserve human rights and justice in their own lands as well as in other lands and that they are now engaged in a common struggle against savage and brutal forces seeking to subjugate the world,
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DECLARE:
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The foregoing declaration may be adhered to by other nations which are, or which may be, rendering material assistance and contributions in the struggle for victory over Hitlerism.
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During the war, "the United Nations" became the official term for the Allies. To join, countries had to sign the Declaration and declare war on the Axis.[20]
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The UN was formulated and negotiated among the delegations from the Allied Big Four (the United States, the United Kingdom, the Soviet Union and China) at the Dumbarton Oaks Conference from 21 September 1944 to 7 October 1944 and they agreed on the aims, structure and functioning of the UN.[21][22][23] After months of planning, the UN Conference on International Organization opened in San Francisco, 25 April 1945, attended by 50 governments and a number of non-governmental organizations involved in drafting the UN Charter.[24][25][26] "The heads of the delegations of the sponsoring countries took turns as chairman of the plenary meetings: Anthony Eden, of Britain, Edward Stettinius, of the United States, T. V. Soong, of China, and Vyacheslav Molotov, of the Soviet Union. At the later meetings, Lord Halifax deputized for Mister Eden, Wellington Koo for T. V. Soong, and Mister Gromyko for Mister Molotov."[27] The UN officially came into existence 24 October 1945, upon ratification of the Charter by the five permanent members of the Security Council—France, the Republic of China, the Soviet Union, the UK and the US—and by a majority of the other 46 signatories.[28]
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The first meetings of the General Assembly, with 51 nations represented,[b] and the Security Council took place in Methodist Central Hall, Westminster, London beginning on 10 January 1946.[28] Debates began at once covering topical issues including the presence of Russian troops in Iranian Azerbaijan, Great Britain's forces in Greece and within days the first veto was cast.[31]
|
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The General Assembly selected New York City as the site for the headquarters of the UN, construction began on 14 September 1948 and the facility was completed on 9 October 1952. Its site—like UN headquarters buildings in Geneva, Vienna, and Nairobi—is designated as international territory.[32] The Norwegian Foreign Minister, Trygve Lie, was elected as the first UN Secretary-General.[28]
|
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+
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Though the UN's primary mandate was peacekeeping, the division between the US and USSR often paralysed the organization, generally allowing it to intervene only in conflicts distant from the Cold War.[33] Two notable exceptions were a Security Council resolution on 7 July 1950 authorizing a US-led coalition to repel the North Korean invasion of South Korea, passed in the absence of the USSR,[28][34] and the signing of the Korean Armistice Agreement on 27 July 1953.[35]
|
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On 29 November 1947, the General Assembly approved a resolution to partition Palestine, approving the creation of the state of Israel.[36] Two years later, Ralph Bunche, a UN official, negotiated an armistice to the resulting conflict.[37] On 7 November 1956, the first UN peacekeeping force was established to end the Suez Crisis;[38] however, the UN was unable to intervene against the USSR's simultaneous invasion of Hungary following that country's revolution.[39]
|
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On 14 July 1960, the UN established United Nations Operation in the Congo (UNOC), the largest military force of its early decades, to bring order to the breakaway State of Katanga, restoring it to the control of the Democratic Republic of the Congo by 11 May 1964.[40] While traveling to meet rebel leader Moise Tshombe during the conflict, Dag Hammarskjöld, often named as one of the UN's most effective Secretaries-General,[41] died in a plane crash; months later he was posthumously awarded the Nobel Peace Prize.[42] In 1964, Hammarskjöld's successor, U Thant, deployed the UN Peacekeeping Force in Cyprus, which would become one of the UN's longest-running peacekeeping missions.[43]
|
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+
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+
With the spread of decolonization in the 1960s, the organization's membership saw an influx of newly independent nations. In 1960 alone, 17 new states joined the UN, 16 of them from Africa.[38] On 25 October 1971, with opposition from the United States, but with the support of many Third World nations, the mainland, communist People's Republic of China was given the Chinese seat on the Security Council in place of the Republic of China that occupied Taiwan; the vote was widely seen as a sign of waning US influence in the organization.[44] Third World nations organized into the Group of 77 coalition under the leadership of Algeria, which briefly became a dominant power at the UN.[45] On 10 November 1975, a bloc comprising the USSR and Third World nations passed a resolution, over the strenuous US and Israeli opposition, declaring Zionism to be racism; the resolution was repealed on 16 December 1991, shortly after the end of the Cold War.[46][47]
|
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With an increasing Third World presence and the failure of UN mediation in conflicts in the Middle East, Vietnam, and Kashmir, the UN increasingly shifted its attention to its ostensibly secondary goals of economic development and cultural exchange.[48] By the 1970s, the UN budget for social and economic development was far greater than its peacekeeping budget.
|
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+
|
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+
After the Cold War, the UN saw a radical expansion in its peacekeeping duties, taking on more missions in ten years than it had in the previous four decades.[49] Between 1988 and 2000, the number of adopted Security Council resolutions more than doubled, and the peacekeeping budget increased more than tenfold.[50][51][52] The UN negotiated an end to the Salvadoran Civil War, launched a successful peacekeeping mission in Namibia, and oversaw democratic elections in post-apartheid South Africa and post-Khmer Rouge Cambodia.[53] In 1991, the UN authorized a US-led coalition that repulsed the Iraqi invasion of Kuwait.[54] Brian Urquhart, Under-Secretary-General from 1971 to 1985, later described the hopes raised by these successes as a "false renaissance" for the organization, given the more troubled missions that followed.[55]
|
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+
|
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Though the UN Charter had been written primarily to prevent aggression by one nation against another, in the early 1990s the UN faced a number of simultaneous, serious crises within nations such as Somalia, Haiti, Mozambique, and the former Yugoslavia.[56] The UN mission in Somalia was widely viewed as a failure after the US withdrawal following casualties in the Battle of Mogadishu, and the UN mission to Bosnia faced "worldwide ridicule" for its indecisive and confused mission in the face of ethnic cleansing.[57] In 1994, the UN Assistance Mission for Rwanda failed to intervene in the Rwandan genocide amid indecision in the Security Council.[58]
|
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+
|
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+
Beginning in the last decades of the Cold War, American and European critics of the UN condemned the organization for perceived mismanagement and corruption.[59] In 1984, US President Ronald Reagan, withdrew his nation's funding from United Nations Educational, Scientific and Cultural Organization (UNESCO) over allegations of mismanagement, followed by the UK and Singapore.[60][61] Boutros Boutros-Ghali, Secretary-General from 1992 to 1996, initiated a reform of the Secretariat, reducing the size of the organization somewhat.[62][63] His successor, Kofi Annan (1997–2006), initiated further management reforms in the face of threats from the US to withhold its UN dues.[63]
|
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+
|
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From the late 1990s to the early 2000s, international interventions authorized by the UN took a wider variety of forms. The UN mission in the Sierra Leone Civil War of 1991–2002 was supplemented by British Royal Marines, and the invasion of Afghanistan in 2001 was overseen by NATO.[64] In 2003, the United States invaded Iraq despite failing to pass a UN Security Council resolution for authorization, prompting a new round of questioning of the organization's effectiveness.[65] Under the eighth Secretary-General, Ban Ki-moon, the UN intervened with peacekeepers in crises such as the War in Darfur in Sudan and the Kivu conflict in the Democratic Republic of Congo and sent observers and chemical weapons inspectors to the Syrian Civil War.[66] In 2013, an internal review of UN actions in the final battles of the Sri Lankan Civil War in 2009 concluded that the organization had suffered "systemic failure".[67] In 2010, the organization suffered the worst loss of life in its history, when 101 personnel died in the Haiti earthquake[68]
|
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+
|
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+
The Millennium Summit was held in 2000 to discuss the UN's role in the 21st century.[69] The three day meeting was the largest gathering of world leaders in history, and culminated in the adoption by all member states of the Millennium Development Goals (MDGs), a commitment to achieve international development in areas such as poverty reduction, gender equality, and public health. Progress towards these goals, which were to be met by 2015, was ultimately uneven. The 2005 World Summit reaffirmed the UN's focus on promoting development, peacekeeping, human rights, and global security.[70] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71]
|
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+
|
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+
In addition to addressing global challenges, the UN has sought to improve its accountability and democratic legitimacy by engaging more with civil society and fostering a global constituency.[72] In an effort to enhance transparency, in 2016 the organization held its first public debate between candidates for Secretary-General.[73] On 1 January 2017, Portuguese diplomat António Guterres, who previously served as UN High Commissioner for Refugees, became the ninth Secretary-General. Guterres has highlighted several key goals for his administration, including an emphasis on diplomacy for preventing conflicts, more effective peacekeeping efforts, and streamlining the organization to be more responsive and versatile to global needs.[74]
|
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+
|
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+
The UN system is based on five principal organs: the General Assembly, the Security Council, the Economic and Social Council (ECOSOC), the International Court of Justice and the UN Secretariat.[75] A sixth principal organ, the Trusteeship Council, suspended operations on 1 November 1994, upon the independence of Palau, the last remaining UN trustee territory.[76]
|
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+
|
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+
Four of the five principal organs are located at the main UN Headquarters in New York City.[77] The International Court of Justice is located in The Hague, while other major agencies are based in the UN offices at Geneva,[78] Vienna,[79] and Nairobi.[80] Other UN institutions are located throughout the world. The six official languages of the UN, used in intergovernmental meetings and documents, are Arabic, Chinese, English, French, Russian, and Spanish.[81] On the basis of the Convention on the Privileges and Immunities of the United Nations, the UN and its agencies are immune from the laws of the countries where they operate, safeguarding the UN's impartiality with regard to the host and member countries.[82]
|
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Below the six organs sit, in the words of the author Linda Fasulo, "an amazing collection of entities and organizations, some of which are actually older than the UN itself and operate with almost complete independence from it".[83] These include specialized agencies, research, and training institutions, programs and funds, and other UN entities.[84]
|
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|
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+
The UN obeys the Noblemaire principle, which is binding on any organization that belongs to the UN system. This principle calls for salaries that will draw and keep citizens of countries where salaries are highest, and also calls for equal pay for work of equal value independent of the employee's nationality.[85][86] In practice, the ICSC takes reference to the highest-paying national civil service.[87] Staff salaries are subject to an internal tax that is administered by the UN organizations.[85][88]
|
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|
71 |
+
|
72 |
+
|
73 |
+
The General Assembly is the main deliberative assembly of the UN. Composed of all UN member states, the assembly meets in regular yearly sessions, but emergency sessions can also be called.[90] The assembly is led by a president, elected from among the member states on a rotating regional basis, and 21 vice-presidents.[91] The first session convened 10 January 1946 in the Methodist Central Hall in London and included representatives of 51 nations.[28]
|
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+
|
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When the General Assembly decides on important questions such as those on peace and security, admission of new members and budgetary matters, a two-thirds majority of those present and voting is required.[92][93] All other questions are decided by a majority vote. Each member country has one vote. Apart from the approval of budgetary matters, resolutions are not binding on the members. The Assembly may make recommendations on any matters within the scope of the UN, except matters of peace and security that are under consideration by the Security Council.[90]
|
76 |
+
|
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Draft resolutions can be forwarded to the General Assembly by its six main committees:[94]
|
78 |
+
|
79 |
+
As well as by the following two committees:
|
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+
|
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+
The Security Council is charged with maintaining peace and security among countries. While other organs of the UN can only make "recommendations" to member states, the Security Council has the power to make binding decisions that member states have agreed to carry out, under the terms of Charter Article 25.[95] The decisions of the council are known as United Nations Security Council resolutions.[96]
|
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|
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+
The Security Council is made up of fifteen member states, consisting of five permanent members—China, France, Russia, the United Kingdom, and the United States—and ten non-permanent members elected for two-year terms by the General Assembly (with end of term date)—Belgium (term ends 2020), Dominican Republic (2020), Estonia (2021), Germany (2020), Indonesia (2020), Niger (2021), Saint Vincent and the Grenadines (2021), South Africa (2020), Tunisia (2021), and Vietnam (2021).[97] The five permanent members hold veto power over UN resolutions, allowing a permanent member to block adoption of a resolution, though not debate. The ten temporary seats are held for two-year terms, with five member states per year voted in by the General Assembly on a regional basis.[98] The presidency of the Security Council rotates alphabetically each month.[99]
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|
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+
The UN Secretariat is headed by the secretary-general, assisted by the deputy secretary-general and a staff of international civil servants worldwide.[100] It provides studies, information, and facilities needed by UN bodies for their meetings. It also carries out tasks as directed by the Security Council, the General Assembly, the Economic and Social Council, and other UN bodies.[101]
|
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+
|
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+
The secretary-general acts as the de facto spokesperson and leader of the UN. The position is defined in the UN Charter as the organization's "chief administrative officer".[102] Article 99 of the charter states that the secretary-general can bring to the Security Council's attention "any matter which in his opinion may threaten the maintenance of international peace and security", a phrase that Secretaries-General since Trygve Lie have interpreted as giving the position broad scope for action on the world stage.[103] The office has evolved into a dual role of an administrator of the UN organization and a diplomat and mediator addressing disputes between member states and finding consensus to global issues.[104]
|
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+
|
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+
The secretary-general is appointed by the General Assembly, after being recommended by the Security Council, where the permanent members have veto power. There are no specific criteria for the post, but over the years it has become accepted that the post shall be held for one or two terms of five years.[105] The current Secretary-General is António Guterres, who replaced Ban Ki-moon in 2017.
|
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|
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+
The International Court of Justice (ICJ), located in The Hague, in the Netherlands, is the primary judicial organ of the UN. Established in 1945 by the UN Charter, the Court began work in 1946 as the successor to the Permanent Court of International Justice. The ICJ is composed of 15 judges who serve 9-year terms and are appointed by the General Assembly; every sitting judge must be from a different nation.[107][108]
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|
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It is based in the Peace Palace in The Hague, sharing the building with the Hague Academy of International Law, a private centre for the study of international law. The ICJ's primary purpose is to adjudicate disputes among states. The court has heard cases related to war crimes, illegal state interference, ethnic cleansing, and other issues.[109] The ICJ can also be called upon by other UN organs to provide advisory opinions.[107]
|
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It is the only organ that is not located in New York.
|
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+
|
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+
The Economic and Social Council (ECOSOC) assists the General Assembly in promoting international economic and social co-operation and development. ECOSOC has 54 members, who are elected by the General Assembly for a three-year term. The president is elected for a one-year term and chosen amongst the small or middle powers represented on ECOSOC. The council has one annual meeting in July, held in either New York or Geneva. Viewed as separate from the specialized bodies it co-ordinates, ECOSOC's functions include information gathering, advising member nations, and making recommendations.[110][111] Owing to its broad mandate of co-ordinating many agencies, ECOSOC has at times been criticized as unfocused or irrelevant.[110][112]
|
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|
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+
ECOSOC's subsidiary bodies include the United Nations Permanent Forum on Indigenous Issues, which advises UN agencies on issues relating to indigenous peoples; the United Nations Forum on Forests, which coordinates and promotes sustainable forest management; the United Nations Statistical Commission, which co-ordinates information-gathering efforts between agencies; and the Commission on Sustainable Development, which co-ordinates efforts between UN agencies and NGOs working towards sustainable development. ECOSOC may also grant consultative status to non-governmental organizations;[110] by 2004, more than 2,200 organizations had received this status.[113]
|
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|
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+
The UN Charter stipulates that each primary organ of the United Nations can establish various specialized agencies to fulfil its duties.[114] Some best-known agencies are the International Atomic Energy Agency, the Food and Agriculture Organization, UNESCO (United Nations Educational, Scientific and Cultural Organization), the World Bank, and the World Health Organization (WHO). The UN performs most of its humanitarian work through these agencies. Examples include mass vaccination programmes (through WHO), the avoidance of famine and malnutrition (through the work of the WFP), and the protection of vulnerable and displaced people (for example, by UNHCR).[115]
|
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+
|
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+
With the addition of South Sudan 14 July 2011,[5] there are 193 UN member states, including all undisputed independent states apart from Vatican City.[116][c]
|
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+
The UN Charter outlines the rules for membership:
|
104 |
+
|
105 |
+
In addition, there are two non-member observer states of the United Nations General Assembly: the Holy See (which holds sovereignty over Vatican City) and the State of Palestine.[118] The Cook Islands and Niue, both states in free association with New Zealand, are full members of several UN specialized agencies and have had their "full treaty-making capacity" recognized by the Secretariat.[119]
|
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|
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+
The Group of 77 (G77) at the UN is a loose coalition of developing nations, designed to promote its members' collective economic interests and create an enhanced joint negotiating capacity in the UN. Seventy-seven nations founded the organization, but by November 2013 the organization had since expanded to 133 member countries.[120] The group was founded 15 June 1964 by the "Joint Declaration of the Seventy-Seven Countries" issued at the United Nations Conference on Trade and Development (UNCTAD). The group held its first major meeting in Algiers in 1967, where it adopted the Charter of Algiers and established the basis for permanent institutional structures.[121] With the adoption of the New International Economic Order by developing countries in the 1970s, the work of the G77 spread throughout the UN system.
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+
The UN, after approval by the Security Council, sends peacekeepers to regions where armed conflict has recently ceased or paused to enforce the terms of peace agreements and to discourage combatants from resuming hostilities. Since the UN does not maintain its own military, peacekeeping forces are voluntarily provided by member states. These soldiers are sometimes nicknamed "Blue Helmets" for their distinctive gear.[122][123] The peacekeeping force as a whole received the Nobel Peace Prize in 1988.[124]
|
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|
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+
In September 2013, the UN had peacekeeping soldiers deployed on 15 missions. The largest was the United Nations Organization Stabilization Mission in the Democratic Republic of the Congo (MONUSCO), which included 20,688 uniformed personnel. The smallest, United Nations Military Observer Group in India and Pakistan (UNMOGIP), included 42 uniformed personnel responsible for monitoring the ceasefire in Jammu and Kashmir. UN peacekeepers with the United Nations Truce Supervision Organization (UNTSO) have been stationed in the Middle East since 1948, the longest-running active peacekeeping mission.[125]
|
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+
|
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+
A study by the RAND Corporation in 2005 found the UN to be successful in two out of three peacekeeping efforts. It compared efforts at nation-building by the UN to those of the United States, and found that seven out of eight UN cases are at peace, as compared with four out of eight US cases at peace.[126] Also in 2005, the Human Security Report documented a decline in the number of wars, genocides, and human rights abuses since the end of the Cold War, and presented evidence, albeit circumstantial, that international activism—mostly spearheaded by the UN—has been the main cause of the decline in armed conflict in that period.[127] Situations in which the UN has not only acted to keep the peace but also intervened include the Korean War (1950–53) and the authorization of intervention in Iraq after the Gulf War (1990–91).[128]
|
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+
|
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+
|
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+
|
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+
The UN has also drawn criticism for perceived failures. In many cases, member states have shown reluctance to achieve or enforce Security Council resolutions. Disagreements in the Security Council about military action and intervention are seen as having failed to prevent the Bangladesh genocide in 1971,[129] the Cambodian genocide in the 1970s,[130] and the Rwandan genocide in 1994.[131] Similarly, UN inaction is blamed for failing to either prevent the Srebrenica massacre in 1995 or complete the peacekeeping operations in 1992–93 during the Somali Civil War.[132] UN peacekeepers have also been accused of child rape, soliciting prostitutes, and sexual abuse during various peacekeeping missions in the Democratic Republic of the Congo,[133] Haiti,[134] Liberia,[135] Sudan and what is now South Sudan,[136] Burundi, and Ivory Coast.[137] Scientists cited UN peacekeepers from Nepal as the likely source of the 2010–13 Haiti cholera outbreak, which killed more than 8,000 Haitians following the 2010 Haiti earthquake.[138]
|
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+
In addition to peacekeeping, the UN is also active in encouraging disarmament. Regulation of armaments was included in the writing of the UN Charter in 1945 and was envisioned as a way of limiting the use of human and economic resources for their creation.[95] The advent of nuclear weapons came only weeks after the signing of the charter, resulting in the first resolution of the first General Assembly meeting calling for specific proposals for "the elimination from national armaments of atomic weapons and of all other major weapons adaptable to mass destruction".[139] The UN has been involved with arms-limitation treaties, such as the Outer Space Treaty (1967), the Treaty on the Non-Proliferation of Nuclear Weapons (1968), the Seabed Arms Control Treaty (1971), the Biological Weapons Convention (1972), the Chemical Weapons Convention (1992), and the Ottawa Treaty (1997), which prohibits landmines.[140] Three UN bodies oversee arms proliferation issues: the International Atomic Energy Agency, the Organisation for the Prohibition of Chemical Weapons, and the Comprehensive Nuclear-Test-Ban Treaty Organization Preparatory Commission.[141]
|
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+
One of the UN's primary purposes is "promoting and encouraging respect for human rights and for fundamental freedoms for all without distinction as to race, sex, language, or religion", and member states pledge to undertake "joint and separate action" to protect these rights.[114][142]
|
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In 1948, the General Assembly adopted a Universal Declaration of Human Rights, drafted by a committee headed by American diplomat and activist Eleanor Roosevelt, and including the French lawyer René Cassin. The document proclaims basic civil, political, and economic rights common to all human beings, though its effectiveness towards achieving these ends has been disputed since its drafting.[143] The Declaration serves as a "common standard of achievement for all peoples and all nations" rather than a legally binding document, but it has become the basis of two binding treaties, the 1966 International Covenant on Civil and Political Rights and International Covenant on Economic, Social and Cultural Rights.[144] In practice, the UN is unable to take significant action against human rights abuses without a Security Council resolution, though it does substantial work in investigating and reporting abuses.[145]
|
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|
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+
In 1979, the General Assembly adopted the Convention on the Elimination of All Forms of Discrimination against Women, followed by the Convention on the Rights of the Child in 1989.[146] With the end of the Cold War, the push for human rights action took on new impetus.[147] The United Nations Commission on Human Rights was formed in 1993 to oversee human rights issues for the UN, following the recommendation of that year's World Conference on Human Rights. Jacques Fomerand, a scholar of the UN, describes this organization's mandate as "broad and vague", with only "meagre" resources to carry it out.[148] In 2006, it was replaced by a Human Rights Council consisting of 47 nations.[149] Also in 2006, the General Assembly passed a Declaration on the Rights of Indigenous Peoples,[150] and in 2011 it passed its first resolution recognizing the rights of LGBT people.[151]
|
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+
Other UN bodies responsible for women's rights issues include United Nations Commission on the Status of Women, a commission of ECOSOC founded in 1946; the United Nations Development Fund for Women, created in 1976; and the United Nations International Research and Training Institute for the Advancement of Women, founded in 1979.[152] The UN Permanent Forum on Indigenous Issues, one of three bodies with a mandate to oversee issues related to indigenous peoples, held its first session in 2002.[153]
|
128 |
+
|
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+
Millennium Development Goals[154]
|
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Another primary purpose of the UN is "to achieve international cooperation in solving international problems of an economic, social, cultural, or humanitarian character".[142] Numerous bodies have been created to work towards this goal, primarily under the authority of the General Assembly and ECOSOC.[155] In 2000, the 192 UN member states agreed to achieve eight Millennium Development Goals by 2015.[156] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71] The SDGs have an associated financing framework called the Addis Ababa Action Agenda.
|
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+
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The UN Development Programme (UNDP), an organization for grant-based technical assistance founded in 1945, is one of the leading bodies in the field of international development. The organization also publishes the UN Human Development Index, a comparative measure ranking countries by poverty, literacy, education, life expectancy, and other factors.[157][158] The Food and Agriculture Organization (FAO), also founded in 1945, promotes agricultural development and food security.[159] UNICEF (the United Nations Children's Fund) was created in 1946 to aid European children after the Second World War and expanded its mission to provide aid around the world and to uphold the convention on the Rights of the Child.[160][161]
|
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The World Bank Group and International Monetary Fund (IMF) are independent, specialized agencies and observers within the UN framework, according to a 1947 agreement. They were initially formed separately from the UN through the Bretton Woods Agreement in 1944.[162] The World Bank provides loans for international development, while the IMF promotes international economic co-operation and gives emergency loans to indebted countries.[163]
|
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The World Health Organization (WHO), which focuses on international health issues and disease eradication, is another of the UN's largest agencies. In 1980, the agency announced that the eradication of smallpox had been completed. In subsequent decades, WHO largely eradicated polio, river blindness, and leprosy.[164] The Joint United Nations Programme on HIV/AIDS (UNAIDS), begun in 1996, co-ordinates the organization's response to the AIDS epidemic.[165] The UN Population Fund, which also dedicates part of its resources to combating HIV, is the world's largest source of funding for reproductive health and family planning services.[166]
|
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+
Along with the International Red Cross and Red Crescent Movement, the UN often takes a leading role in co-ordinating emergency relief.[167] The World Food Programme (WFP), created in 1961, provides food aid in response to famine, natural disasters, and armed conflict. The organization reports that it feeds an average of 90 million people in 80 nations each year.[167][168] The Office of the United Nations High Commissioner for Refugees (UNHCR), established in 1950, works to protect the rights of refugees, asylum seekers, and stateless people.[169] UNHCR and WFP programmes are funded by voluntary contributions from governments, corporations, and individuals, though the UNHCR's administrative costs are paid for by the UN's primary budget.[170]
|
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+
Since the UN's creation, over 80 colonies have attained independence. The General Assembly adopted the Declaration on the Granting of Independence to Colonial Countries and Peoples in 1960 with no votes against but abstentions from all major colonial powers. The UN works towards decolonization through groups including the UN Committee on Decolonization, created in 1962.[171] The committee lists seventeen remaining "Non-Self-Governing Territories", the largest and most populous of which is Western Sahara.[172]
|
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+
Beginning with the formation of the UN Environmental Programme (UNEP) in 1972, the UN has made environmental issues a prominent part of its agenda. A lack of success in the first two decades of UN work in this area led to the 1992 Earth Summit in Rio de Janeiro, Brazil, which sought to give new impetus to these efforts.[173] In 1988, the UNEP and the World Meteorological Organization (WMO), another UN organization, established the Intergovernmental Panel on Climate Change, which assesses and reports on research on global warming.[174] The UN-sponsored Kyoto Protocol, signed in 1997, set legally binding emissions reduction targets for ratifying states.[175]
|
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+
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The UN also declares and co-ordinates international observances, periods of time to observe issues of international interest or concern. Examples include World Tuberculosis Day, Earth Day, and the International Year of Deserts and Desertification.[176]
|
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+
|
147 |
+
The UN is financed from assessed and voluntary contributions from member states. The General Assembly approves the regular budget and determines the assessment for each member. This is broadly based on the relative capacity of each country to pay, as measured by its gross national income (GNI), with adjustments for external debt and low per capita income.[178] The two-year budget for 2012–13 was $5.512 billion in total.[179]
|
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|
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+
The Assembly has established the principle that the UN should not be unduly dependent on any one member to finance its operations. Thus, there is a "ceiling" rate, setting the maximum amount that any member can be assessed for the regular budget. In December 2000, the Assembly revised the scale of assessments in response to pressure from the United States. As part of that revision, the regular budget ceiling was reduced from 25% to 22%.[180] For the least developed countries (LDCs), a ceiling rate of 0.01% is applied.[178] In addition to the ceiling rates, the minimum amount assessed to any member nation (or "floor" rate) is set at 0.001% of the UN budget ($55,120 for the two year budget 2013–2014).[181]
|
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+
A large share of the UN's expenditure addresses its core mission of peace and security, and this budget is assessed separately from the main organizational budget.[182] The peacekeeping budget for the 2015–16 fiscal year was $8.27 billion, supporting 82,318 troops deployed in 15 missions around the world.[125] UN peace operations are funded by assessments, using a formula derived from the regular funding scale that includes a weighted surcharge for the five permanent Security Council members, who must approve all peacekeeping operations. This surcharge serves to offset discounted peacekeeping assessment rates for less developed countries.
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the largest contributors for the UN peacekeeping financial operations for the period 2019–2021 are : the United States 27.89% China 15.21%, Japan 8.56%, Germany 6.09% , the United Kingdom 5.78%, France 5.61%, Italy 3.30% and the Russian Federation 3.04%.[183]
|
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+
|
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+
Special UN programmes not included in the regular budget, such as UNICEF and the World Food Programme, are financed by voluntary contributions from member governments, corporations, and private individuals.[184][185]
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+
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A number of agencies and individuals associated with the UN have won the Nobel Peace Prize in recognition of their work. Two Secretaries-General, Dag Hammarskjöld, and Kofi Annan, were each awarded the prize (in 1961 and 2001, respectively), as were Ralph Bunche (1950), a UN negotiator, René Cassin (1968), a contributor to the Universal Declaration of Human Rights, and the US Secretary of State Cordell Hull (1945), the latter for his role in the organization's founding. Lester B. Pearson, the Canadian Secretary of State for External Affairs, was awarded the prize in 1957 for his role in organizing the UN's first peacekeeping force to resolve the Suez Crisis. UNICEF won the prize in 1965, the International Labour Organization in 1969, the UN Peace-Keeping Forces in 1988, the International Atomic Energy Agency (which reports to the UN) in 2005, and the UN-supported Organisation for the Prohibition of Chemical Weapons in 2013. The UN High Commissioner for Refugees was awarded in 1954 and 1981, becoming one of only two recipients to win the prize twice. The UN as a whole was awarded the prize in 2001, sharing it with Annan.[186] In 2007, IPCC received the prize "for their efforts to build up and disseminate greater knowledge about man-made climate change, and to lay the foundations for the measures that are needed to counteract such change."[187]
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Since its founding, there have been many calls for reform of the UN but little consensus on how to do so. Some want the UN to play a greater or more effective role in world affairs, while others want its role reduced to humanitarian work. There have also been numerous calls for the UN Security Council's membership to be increased, for different ways of electing the UN's Secretary-General, and for a UN Parliamentary Assembly. Jacques Fomerand states the most enduring divide in views of the UN is "the North–South split" between richer Northern nations and developing Southern nations. Southern nations tend to favour a more empowered UN with a stronger General Assembly, allowing them a greater voice in world affairs, while Northern nations prefer an economically laissez-faire UN that focuses on transnational threats such as terrorism.[188]
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After World War II, the French Committee of National Liberation was late to be recognized by the US as the government of France, and so the country was initially excluded from the conferences that created the new organization. The future French president Charles de Gaulle criticized the UN, famously calling it a machin ("contraption"), and was not convinced that a global security alliance would help maintain world peace, preferring direct defence treaties between countries.[189] Throughout the Cold War, both the US and USSR repeatedly accused the UN of favouring the other. In 1953, the USSR effectively forced the resignation of Trygve Lie, the Secretary-General, through its refusal to deal with him, while in the 1950s and 1960s, a popular US bumper sticker read, "You can't spell communism without U.N."[190] In a sometimes-misquoted statement, President George W. Bush stated in February 2003 (referring to UN uncertainty towards Iraqi provocations under the Saddam Hussein regime) that "free nations will not allow the UN to fade into history as an ineffective, irrelevant debating society."[191][192][193] In contrast, the French President, François Hollande, stated in 2012 that "France trusts the United Nations. She knows that no state, no matter how powerful, can solve urgent problems, fight for development and bring an end to all crises ... France wants the UN to be the centre of global governance."[194] Critics such as Dore Gold, an Israeli diplomat, Robert S. Wistrich, a British scholar, Alan Dershowitz, an American legal scholar, Mark Dreyfus, an Australian politician, and the Anti-Defamation League consider UN attention to Israel's treatment of Palestinians to be excessive.[195] In September 2015, Saudi Arabia's Faisal bin Hassan Trad has been elected Chair of the UN Human Rights Council panel that appoints independent experts,[196] a move criticized by human rights groups.[197][198]
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Since 1971, the Republic of China on Taiwan has been excluded from the UN and since then has always been rejected in new applications. Taiwanese citizens are also not allowed to enter the buildings of the United Nations with ROC passports. In this way, critics agree that the UN is failing its own development goals and guidelines. This criticism also brought pressure from the People's Republic of China, which regards the territories administered by the ROC as their own territory.[199][200]
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Critics have also accused the UN of bureaucratic inefficiency, waste, and corruption. In 1976, the General Assembly established the Joint Inspection Unit to seek out inefficiencies within the UN system. During the 1990s, the US withheld dues citing inefficiency and only started repayment on the condition that a major reforms initiative be introduced. In 1994, the Office of Internal Oversight Services (OIOS) was established by the General Assembly to serve as an efficiency watchdog.[201] In 1994, former Special Representative of the Secretary-General of the UN to Somalia Mohamed Sahnoun published "Somalia: The Missed Opportunities",[202] a book in which he analyses the reasons for the failure of the 1992 UN intervention in Somalia, showing that, between the start of the Somali civil war in 1988 and the fall of the Siad Barre regime in January 1991, the UN missed at least three opportunities to prevent major human tragedies; when the UN tried to provide humanitarian assistance, they were totally outperformed by NGOs, whose competence and dedication sharply contrasted with the UN's excessive caution and bureaucratic inefficiencies. If radical reform were not undertaken, warned Mohamed Sahnoun, then the UN would continue to respond to such crises with inept improvisation.[203] In 2004, the UN faced accusations that its recently ended Oil-for-Food Programme — in which Iraq had been allowed to trade oil for basic needs to relieve the pressure of sanctions — had suffered from widespread corruption, including billions of dollars of kickbacks. An independent inquiry created by the UN found that many of its officials had been involved, as well as raising "significant" questions about the role of Kojo Annan, the son of Kofi Annan.[204]
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In evaluating the UN as a whole, Jacques Fomerand writes that the "accomplishments of the United Nations in the last 60 years are impressive in their own terms. Progress in human development during the 20th century has been dramatic, and the UN and its agencies have certainly helped the world become a more hospitable and livable place for millions."[205] Evaluating the first 50 years of the UN's history, the author Stanley Meisler writes that "the United Nations never fulfilled the hopes of its founders, but it accomplished a great deal nevertheless", citing its role in decolonization and its many successful peacekeeping efforts.[206] The British historian Paul Kennedy states that while the organization has suffered some major setbacks, "when all its aspects are considered, the UN has brought great benefits to our generation and ... will bring benefits to our children's and grandchildren's generations as well."[207]
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Core features of the UN apparatus, such as the veto privileges of some nations in the Security Council, are often described as fundamentally undemocratic, contrary to the UN mission, and as a main cause of inaction on genocides and crimes against humanity.[208][209]
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The United Nations has inspired the extracurricular activity Model United Nations (MUN). MUN is a simulation of United Nations activity based on the UN agenda and following UN procedure. MUN is usually attended by high school and university students who organize conferences to simulate the various UN committees to discuss important issues of the day.[210] Today Model United Nations educates tens of thousands on United Nations activity around the world. Model United Nations has many famous and notable alumni, such as former Secretary-General of the United Nations Ban Ki-moon.[211]
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The United Nations (UN) is an intergovernmental organization that aims to maintain international peace and security, develop friendly relations among nations, achieve international cooperation, and be a centre for harmonizing the actions of nations.[2] It is the largest, most familiar, most internationally represented and most powerful intergovernmental organization in the world. The UN is headquartered on international territory in New York City, with its other main offices in Geneva, Nairobi, Vienna and The Hague.
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The UN was established after World War II with the aim of preventing future wars, succeeding the ineffective League of Nations.[3] On 25 April 1945, 50 governments met in San Francisco for a conference and started drafting the UN Charter, which was adopted on 25 June 1945 and took effect on 24 October 1945, when the UN began operations. Pursuant to the Charter, the organization's objectives include maintaining international peace and security, protecting human rights, delivering humanitarian aid, promoting sustainable development, and upholding international law.[4] At its founding, the UN had 51 member states; this number grew to 193 in 2011,[5] representing almost all of the world's sovereign states.
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The organization's mission to preserve world peace was complicated in its early decades by the Cold War between the United States and Soviet Union and their respective allies. Its missions have consisted primarily of unarmed military observers and lightly armed troops with primarily monitoring, reporting and confidence-building roles.[6] UN membership grew significantly following widespread decolonization beginning in the 1960s. Since then, 80 former colonies have gained independence, including 11 trust territories that had been monitored by the Trusteeship Council.[7] By the 1970s, the UN's budget for economic and social development programmes far outstripped its spending on peacekeeping. After the end of the Cold War, the UN shifted and expanded its field operations, undertaking a wide variety of complex tasks.[8]
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The UN has six principal organs: the General Assembly; the Security Council; the Economic and Social Council (ECOSOC); the Trusteeship Council; the International Court of Justice; and the UN Secretariat. The UN System includes a multitude of specialized agencies, such as the World Bank Group, the World Health Organization, the World Food Programme, UNESCO, and UNICEF. Additionally, non-governmental organizations may be granted consultative status with ECOSOC and other agencies to participate in the UN's work. The UN's chief administrative officer is the Secretary-General, currently Portuguese politician and diplomat António Guterres, who began his five year-term on 1 January 2017. The organization is financed by assessed and voluntary contributions from its member states.
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The UN, its officers, and its agencies have won many Nobel Peace Prizes, though other evaluations of its effectiveness have been mixed. Some commentators believe the organization to be an important force for peace and human development, while others have called it ineffective, biased, or corrupt.
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In the century prior to the UN's creation, several international treaty organizations such as the International Committee of the Red Cross were formed to ensure protection and assistance for victims of armed conflict and strife.[9]
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In 1914, a political assassination in Sarajevo set off a chain of events that led to the outbreak of World War I. As more and more young men were sent down into the trenches, influential voices in the United States and Britain began calling for the establishment of a permanent international body to maintain peace in the postwar world. President Woodrow Wilson became a vocal advocate of this concept, and in 1918 he included a sketch of the international body in his 14-point proposal to end the war. In November 1918, the Central Powers agreed to an armistice to halt the killing in World War I. Two months later, the Allies met with Germany and Austria-Hungary at Versailles to hammer out formal peace terms. President Wilson wanted peace, but the United Kingdom and France disagreed, forcing harsh war reparations on their former enemies. The League of Nations was approved, and in the summer of 1919 Wilson presented the Treaty of Versailles and the Covenant of the League of Nations to the US Senate for ratification. On 10 January 1920, the League of Nations formally came into being when the Covenant of the League of Nations, ratified by 42 nations in 1919, took effect.[10]
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However, at some point the League became ineffective when it failed to act against the Japanese invasion of Manchuria as in February 1933, 40 nations voted for Japan to withdraw from Manchuria but Japan voted against it and walked out of the League instead of withdrawing from Manchuria.[11] It also failed against the Second Italo-Ethiopian War despite trying to talk to Benito Mussolini as he used the time to send an army to Africa, so the League had a plan for Mussolini to just take a part of Ethiopia, but he ignored the League and invaded Ethiopia, the League tried putting sanctions on Italy, but Italy had already conquered Ethiopia and the League had failed.[12] After Italy conquered Ethiopia, Italy and other nations left the league. But all of them realized that it had failed and they began to re-arm as fast as possible.
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During 1938, Britain and France tried negotiating directly with Hitler but this failed in 1939 when Hitler invaded Czechoslovakia.
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When war broke out in 1939, the League closed down and its headquarters in Geneva remained empty throughout the war.[13] Although the United States never joined the League, the country did support its economic and social missions through the work of private philanthropies and by sending representatives to committees.
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The earliest concrete plan for a new world organization began under the aegis of the U.S. State Department in 1939.[14] The text of the "Declaration by United Nations" was drafted at the White House on 29 December 1941, by President Franklin D. Roosevelt, Prime Minister Winston Churchill, and Roosevelt aide Harry Hopkins. It incorporated Soviet suggestions but left no role for France. "Four Policemen" was coined to refer to four major Allied countries, United States, United Kingdom, Soviet Union, and Republic of China, which emerged in the Declaration by United Nations.[15] Roosevelt first coined the term United Nations to describe the Allied countries.[a] "On New Year's Day 1942, President Roosevelt, Prime Minister Churchill, Maxim Litvinov, of the USSR, and T. V. Soong, of China, signed a short document which later came to be known as the United Nations Declaration, and the next day the representatives of twenty-two other nations added their signatures."[16] The term United Nations was first officially used when 26 governments signed this Declaration. One major change from the Atlantic Charter was the addition of a provision for religious freedom, which Stalin approved after Roosevelt insisted.[17][18] By 1 March 1945, 21 additional states had signed.[19]
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A JOINT DECLARATION BY THE UNITED STATES OF AMERICA, THE UNITED KINGDOM OF GREAT BRITAIN AND NORTHERN IRELAND, THE UNION OF SOVIET SOCIALIST REPUBLICS, CHINA, AUSTRALIA, BELGIUM, CANADA, COSTA RICA, CUBA, CZECHOSLOVAKIA, DOMINICAN REPUBLIC, EL SALVADOR, GREECE, GUATEMALA, HAITI, HONDURAS, INDIA, LUXEMBOURG, NETHERLANDS, NEW ZEALAND, NICARAGUA, NORWAY, PANAMA, POLAND, SOUTH AFRICA, YUGOSLAVIA
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The Governments signatory hereto,
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Having subscribed to a common program of purposes and principles embodied in the Joint Declaration of the President of the United States of America and the Prime Minister of Great Britain dated August 14, 1941, known as the Atlantic Charter,
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Being convinced that complete victory over their enemies is essential to defend life, liberty, independence and religious freedom, and to preserve human rights and justice in their own lands as well as in other lands and that they are now engaged in a common struggle against savage and brutal forces seeking to subjugate the world,
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DECLARE:
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The foregoing declaration may be adhered to by other nations which are, or which may be, rendering material assistance and contributions in the struggle for victory over Hitlerism.
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During the war, "the United Nations" became the official term for the Allies. To join, countries had to sign the Declaration and declare war on the Axis.[20]
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The UN was formulated and negotiated among the delegations from the Allied Big Four (the United States, the United Kingdom, the Soviet Union and China) at the Dumbarton Oaks Conference from 21 September 1944 to 7 October 1944 and they agreed on the aims, structure and functioning of the UN.[21][22][23] After months of planning, the UN Conference on International Organization opened in San Francisco, 25 April 1945, attended by 50 governments and a number of non-governmental organizations involved in drafting the UN Charter.[24][25][26] "The heads of the delegations of the sponsoring countries took turns as chairman of the plenary meetings: Anthony Eden, of Britain, Edward Stettinius, of the United States, T. V. Soong, of China, and Vyacheslav Molotov, of the Soviet Union. At the later meetings, Lord Halifax deputized for Mister Eden, Wellington Koo for T. V. Soong, and Mister Gromyko for Mister Molotov."[27] The UN officially came into existence 24 October 1945, upon ratification of the Charter by the five permanent members of the Security Council—France, the Republic of China, the Soviet Union, the UK and the US—and by a majority of the other 46 signatories.[28]
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The first meetings of the General Assembly, with 51 nations represented,[b] and the Security Council took place in Methodist Central Hall, Westminster, London beginning on 10 January 1946.[28] Debates began at once covering topical issues including the presence of Russian troops in Iranian Azerbaijan, Great Britain's forces in Greece and within days the first veto was cast.[31]
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The General Assembly selected New York City as the site for the headquarters of the UN, construction began on 14 September 1948 and the facility was completed on 9 October 1952. Its site—like UN headquarters buildings in Geneva, Vienna, and Nairobi—is designated as international territory.[32] The Norwegian Foreign Minister, Trygve Lie, was elected as the first UN Secretary-General.[28]
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Though the UN's primary mandate was peacekeeping, the division between the US and USSR often paralysed the organization, generally allowing it to intervene only in conflicts distant from the Cold War.[33] Two notable exceptions were a Security Council resolution on 7 July 1950 authorizing a US-led coalition to repel the North Korean invasion of South Korea, passed in the absence of the USSR,[28][34] and the signing of the Korean Armistice Agreement on 27 July 1953.[35]
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On 29 November 1947, the General Assembly approved a resolution to partition Palestine, approving the creation of the state of Israel.[36] Two years later, Ralph Bunche, a UN official, negotiated an armistice to the resulting conflict.[37] On 7 November 1956, the first UN peacekeeping force was established to end the Suez Crisis;[38] however, the UN was unable to intervene against the USSR's simultaneous invasion of Hungary following that country's revolution.[39]
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On 14 July 1960, the UN established United Nations Operation in the Congo (UNOC), the largest military force of its early decades, to bring order to the breakaway State of Katanga, restoring it to the control of the Democratic Republic of the Congo by 11 May 1964.[40] While traveling to meet rebel leader Moise Tshombe during the conflict, Dag Hammarskjöld, often named as one of the UN's most effective Secretaries-General,[41] died in a plane crash; months later he was posthumously awarded the Nobel Peace Prize.[42] In 1964, Hammarskjöld's successor, U Thant, deployed the UN Peacekeeping Force in Cyprus, which would become one of the UN's longest-running peacekeeping missions.[43]
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With the spread of decolonization in the 1960s, the organization's membership saw an influx of newly independent nations. In 1960 alone, 17 new states joined the UN, 16 of them from Africa.[38] On 25 October 1971, with opposition from the United States, but with the support of many Third World nations, the mainland, communist People's Republic of China was given the Chinese seat on the Security Council in place of the Republic of China that occupied Taiwan; the vote was widely seen as a sign of waning US influence in the organization.[44] Third World nations organized into the Group of 77 coalition under the leadership of Algeria, which briefly became a dominant power at the UN.[45] On 10 November 1975, a bloc comprising the USSR and Third World nations passed a resolution, over the strenuous US and Israeli opposition, declaring Zionism to be racism; the resolution was repealed on 16 December 1991, shortly after the end of the Cold War.[46][47]
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With an increasing Third World presence and the failure of UN mediation in conflicts in the Middle East, Vietnam, and Kashmir, the UN increasingly shifted its attention to its ostensibly secondary goals of economic development and cultural exchange.[48] By the 1970s, the UN budget for social and economic development was far greater than its peacekeeping budget.
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After the Cold War, the UN saw a radical expansion in its peacekeeping duties, taking on more missions in ten years than it had in the previous four decades.[49] Between 1988 and 2000, the number of adopted Security Council resolutions more than doubled, and the peacekeeping budget increased more than tenfold.[50][51][52] The UN negotiated an end to the Salvadoran Civil War, launched a successful peacekeeping mission in Namibia, and oversaw democratic elections in post-apartheid South Africa and post-Khmer Rouge Cambodia.[53] In 1991, the UN authorized a US-led coalition that repulsed the Iraqi invasion of Kuwait.[54] Brian Urquhart, Under-Secretary-General from 1971 to 1985, later described the hopes raised by these successes as a "false renaissance" for the organization, given the more troubled missions that followed.[55]
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Though the UN Charter had been written primarily to prevent aggression by one nation against another, in the early 1990s the UN faced a number of simultaneous, serious crises within nations such as Somalia, Haiti, Mozambique, and the former Yugoslavia.[56] The UN mission in Somalia was widely viewed as a failure after the US withdrawal following casualties in the Battle of Mogadishu, and the UN mission to Bosnia faced "worldwide ridicule" for its indecisive and confused mission in the face of ethnic cleansing.[57] In 1994, the UN Assistance Mission for Rwanda failed to intervene in the Rwandan genocide amid indecision in the Security Council.[58]
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Beginning in the last decades of the Cold War, American and European critics of the UN condemned the organization for perceived mismanagement and corruption.[59] In 1984, US President Ronald Reagan, withdrew his nation's funding from United Nations Educational, Scientific and Cultural Organization (UNESCO) over allegations of mismanagement, followed by the UK and Singapore.[60][61] Boutros Boutros-Ghali, Secretary-General from 1992 to 1996, initiated a reform of the Secretariat, reducing the size of the organization somewhat.[62][63] His successor, Kofi Annan (1997–2006), initiated further management reforms in the face of threats from the US to withhold its UN dues.[63]
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From the late 1990s to the early 2000s, international interventions authorized by the UN took a wider variety of forms. The UN mission in the Sierra Leone Civil War of 1991–2002 was supplemented by British Royal Marines, and the invasion of Afghanistan in 2001 was overseen by NATO.[64] In 2003, the United States invaded Iraq despite failing to pass a UN Security Council resolution for authorization, prompting a new round of questioning of the organization's effectiveness.[65] Under the eighth Secretary-General, Ban Ki-moon, the UN intervened with peacekeepers in crises such as the War in Darfur in Sudan and the Kivu conflict in the Democratic Republic of Congo and sent observers and chemical weapons inspectors to the Syrian Civil War.[66] In 2013, an internal review of UN actions in the final battles of the Sri Lankan Civil War in 2009 concluded that the organization had suffered "systemic failure".[67] In 2010, the organization suffered the worst loss of life in its history, when 101 personnel died in the Haiti earthquake[68]
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The Millennium Summit was held in 2000 to discuss the UN's role in the 21st century.[69] The three day meeting was the largest gathering of world leaders in history, and culminated in the adoption by all member states of the Millennium Development Goals (MDGs), a commitment to achieve international development in areas such as poverty reduction, gender equality, and public health. Progress towards these goals, which were to be met by 2015, was ultimately uneven. The 2005 World Summit reaffirmed the UN's focus on promoting development, peacekeeping, human rights, and global security.[70] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71]
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In addition to addressing global challenges, the UN has sought to improve its accountability and democratic legitimacy by engaging more with civil society and fostering a global constituency.[72] In an effort to enhance transparency, in 2016 the organization held its first public debate between candidates for Secretary-General.[73] On 1 January 2017, Portuguese diplomat António Guterres, who previously served as UN High Commissioner for Refugees, became the ninth Secretary-General. Guterres has highlighted several key goals for his administration, including an emphasis on diplomacy for preventing conflicts, more effective peacekeeping efforts, and streamlining the organization to be more responsive and versatile to global needs.[74]
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The UN system is based on five principal organs: the General Assembly, the Security Council, the Economic and Social Council (ECOSOC), the International Court of Justice and the UN Secretariat.[75] A sixth principal organ, the Trusteeship Council, suspended operations on 1 November 1994, upon the independence of Palau, the last remaining UN trustee territory.[76]
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Four of the five principal organs are located at the main UN Headquarters in New York City.[77] The International Court of Justice is located in The Hague, while other major agencies are based in the UN offices at Geneva,[78] Vienna,[79] and Nairobi.[80] Other UN institutions are located throughout the world. The six official languages of the UN, used in intergovernmental meetings and documents, are Arabic, Chinese, English, French, Russian, and Spanish.[81] On the basis of the Convention on the Privileges and Immunities of the United Nations, the UN and its agencies are immune from the laws of the countries where they operate, safeguarding the UN's impartiality with regard to the host and member countries.[82]
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Below the six organs sit, in the words of the author Linda Fasulo, "an amazing collection of entities and organizations, some of which are actually older than the UN itself and operate with almost complete independence from it".[83] These include specialized agencies, research, and training institutions, programs and funds, and other UN entities.[84]
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The UN obeys the Noblemaire principle, which is binding on any organization that belongs to the UN system. This principle calls for salaries that will draw and keep citizens of countries where salaries are highest, and also calls for equal pay for work of equal value independent of the employee's nationality.[85][86] In practice, the ICSC takes reference to the highest-paying national civil service.[87] Staff salaries are subject to an internal tax that is administered by the UN organizations.[85][88]
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The General Assembly is the main deliberative assembly of the UN. Composed of all UN member states, the assembly meets in regular yearly sessions, but emergency sessions can also be called.[90] The assembly is led by a president, elected from among the member states on a rotating regional basis, and 21 vice-presidents.[91] The first session convened 10 January 1946 in the Methodist Central Hall in London and included representatives of 51 nations.[28]
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When the General Assembly decides on important questions such as those on peace and security, admission of new members and budgetary matters, a two-thirds majority of those present and voting is required.[92][93] All other questions are decided by a majority vote. Each member country has one vote. Apart from the approval of budgetary matters, resolutions are not binding on the members. The Assembly may make recommendations on any matters within the scope of the UN, except matters of peace and security that are under consideration by the Security Council.[90]
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Draft resolutions can be forwarded to the General Assembly by its six main committees:[94]
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As well as by the following two committees:
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The Security Council is charged with maintaining peace and security among countries. While other organs of the UN can only make "recommendations" to member states, the Security Council has the power to make binding decisions that member states have agreed to carry out, under the terms of Charter Article 25.[95] The decisions of the council are known as United Nations Security Council resolutions.[96]
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The Security Council is made up of fifteen member states, consisting of five permanent members—China, France, Russia, the United Kingdom, and the United States—and ten non-permanent members elected for two-year terms by the General Assembly (with end of term date)—Belgium (term ends 2020), Dominican Republic (2020), Estonia (2021), Germany (2020), Indonesia (2020), Niger (2021), Saint Vincent and the Grenadines (2021), South Africa (2020), Tunisia (2021), and Vietnam (2021).[97] The five permanent members hold veto power over UN resolutions, allowing a permanent member to block adoption of a resolution, though not debate. The ten temporary seats are held for two-year terms, with five member states per year voted in by the General Assembly on a regional basis.[98] The presidency of the Security Council rotates alphabetically each month.[99]
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The UN Secretariat is headed by the secretary-general, assisted by the deputy secretary-general and a staff of international civil servants worldwide.[100] It provides studies, information, and facilities needed by UN bodies for their meetings. It also carries out tasks as directed by the Security Council, the General Assembly, the Economic and Social Council, and other UN bodies.[101]
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The secretary-general acts as the de facto spokesperson and leader of the UN. The position is defined in the UN Charter as the organization's "chief administrative officer".[102] Article 99 of the charter states that the secretary-general can bring to the Security Council's attention "any matter which in his opinion may threaten the maintenance of international peace and security", a phrase that Secretaries-General since Trygve Lie have interpreted as giving the position broad scope for action on the world stage.[103] The office has evolved into a dual role of an administrator of the UN organization and a diplomat and mediator addressing disputes between member states and finding consensus to global issues.[104]
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The secretary-general is appointed by the General Assembly, after being recommended by the Security Council, where the permanent members have veto power. There are no specific criteria for the post, but over the years it has become accepted that the post shall be held for one or two terms of five years.[105] The current Secretary-General is António Guterres, who replaced Ban Ki-moon in 2017.
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The International Court of Justice (ICJ), located in The Hague, in the Netherlands, is the primary judicial organ of the UN. Established in 1945 by the UN Charter, the Court began work in 1946 as the successor to the Permanent Court of International Justice. The ICJ is composed of 15 judges who serve 9-year terms and are appointed by the General Assembly; every sitting judge must be from a different nation.[107][108]
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It is based in the Peace Palace in The Hague, sharing the building with the Hague Academy of International Law, a private centre for the study of international law. The ICJ's primary purpose is to adjudicate disputes among states. The court has heard cases related to war crimes, illegal state interference, ethnic cleansing, and other issues.[109] The ICJ can also be called upon by other UN organs to provide advisory opinions.[107]
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It is the only organ that is not located in New York.
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The Economic and Social Council (ECOSOC) assists the General Assembly in promoting international economic and social co-operation and development. ECOSOC has 54 members, who are elected by the General Assembly for a three-year term. The president is elected for a one-year term and chosen amongst the small or middle powers represented on ECOSOC. The council has one annual meeting in July, held in either New York or Geneva. Viewed as separate from the specialized bodies it co-ordinates, ECOSOC's functions include information gathering, advising member nations, and making recommendations.[110][111] Owing to its broad mandate of co-ordinating many agencies, ECOSOC has at times been criticized as unfocused or irrelevant.[110][112]
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ECOSOC's subsidiary bodies include the United Nations Permanent Forum on Indigenous Issues, which advises UN agencies on issues relating to indigenous peoples; the United Nations Forum on Forests, which coordinates and promotes sustainable forest management; the United Nations Statistical Commission, which co-ordinates information-gathering efforts between agencies; and the Commission on Sustainable Development, which co-ordinates efforts between UN agencies and NGOs working towards sustainable development. ECOSOC may also grant consultative status to non-governmental organizations;[110] by 2004, more than 2,200 organizations had received this status.[113]
|
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The UN Charter stipulates that each primary organ of the United Nations can establish various specialized agencies to fulfil its duties.[114] Some best-known agencies are the International Atomic Energy Agency, the Food and Agriculture Organization, UNESCO (United Nations Educational, Scientific and Cultural Organization), the World Bank, and the World Health Organization (WHO). The UN performs most of its humanitarian work through these agencies. Examples include mass vaccination programmes (through WHO), the avoidance of famine and malnutrition (through the work of the WFP), and the protection of vulnerable and displaced people (for example, by UNHCR).[115]
|
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With the addition of South Sudan 14 July 2011,[5] there are 193 UN member states, including all undisputed independent states apart from Vatican City.[116][c]
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The UN Charter outlines the rules for membership:
|
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In addition, there are two non-member observer states of the United Nations General Assembly: the Holy See (which holds sovereignty over Vatican City) and the State of Palestine.[118] The Cook Islands and Niue, both states in free association with New Zealand, are full members of several UN specialized agencies and have had their "full treaty-making capacity" recognized by the Secretariat.[119]
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The Group of 77 (G77) at the UN is a loose coalition of developing nations, designed to promote its members' collective economic interests and create an enhanced joint negotiating capacity in the UN. Seventy-seven nations founded the organization, but by November 2013 the organization had since expanded to 133 member countries.[120] The group was founded 15 June 1964 by the "Joint Declaration of the Seventy-Seven Countries" issued at the United Nations Conference on Trade and Development (UNCTAD). The group held its first major meeting in Algiers in 1967, where it adopted the Charter of Algiers and established the basis for permanent institutional structures.[121] With the adoption of the New International Economic Order by developing countries in the 1970s, the work of the G77 spread throughout the UN system.
|
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The UN, after approval by the Security Council, sends peacekeepers to regions where armed conflict has recently ceased or paused to enforce the terms of peace agreements and to discourage combatants from resuming hostilities. Since the UN does not maintain its own military, peacekeeping forces are voluntarily provided by member states. These soldiers are sometimes nicknamed "Blue Helmets" for their distinctive gear.[122][123] The peacekeeping force as a whole received the Nobel Peace Prize in 1988.[124]
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In September 2013, the UN had peacekeeping soldiers deployed on 15 missions. The largest was the United Nations Organization Stabilization Mission in the Democratic Republic of the Congo (MONUSCO), which included 20,688 uniformed personnel. The smallest, United Nations Military Observer Group in India and Pakistan (UNMOGIP), included 42 uniformed personnel responsible for monitoring the ceasefire in Jammu and Kashmir. UN peacekeepers with the United Nations Truce Supervision Organization (UNTSO) have been stationed in the Middle East since 1948, the longest-running active peacekeeping mission.[125]
|
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A study by the RAND Corporation in 2005 found the UN to be successful in two out of three peacekeeping efforts. It compared efforts at nation-building by the UN to those of the United States, and found that seven out of eight UN cases are at peace, as compared with four out of eight US cases at peace.[126] Also in 2005, the Human Security Report documented a decline in the number of wars, genocides, and human rights abuses since the end of the Cold War, and presented evidence, albeit circumstantial, that international activism—mostly spearheaded by the UN—has been the main cause of the decline in armed conflict in that period.[127] Situations in which the UN has not only acted to keep the peace but also intervened include the Korean War (1950–53) and the authorization of intervention in Iraq after the Gulf War (1990–91).[128]
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The UN has also drawn criticism for perceived failures. In many cases, member states have shown reluctance to achieve or enforce Security Council resolutions. Disagreements in the Security Council about military action and intervention are seen as having failed to prevent the Bangladesh genocide in 1971,[129] the Cambodian genocide in the 1970s,[130] and the Rwandan genocide in 1994.[131] Similarly, UN inaction is blamed for failing to either prevent the Srebrenica massacre in 1995 or complete the peacekeeping operations in 1992–93 during the Somali Civil War.[132] UN peacekeepers have also been accused of child rape, soliciting prostitutes, and sexual abuse during various peacekeeping missions in the Democratic Republic of the Congo,[133] Haiti,[134] Liberia,[135] Sudan and what is now South Sudan,[136] Burundi, and Ivory Coast.[137] Scientists cited UN peacekeepers from Nepal as the likely source of the 2010–13 Haiti cholera outbreak, which killed more than 8,000 Haitians following the 2010 Haiti earthquake.[138]
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In addition to peacekeeping, the UN is also active in encouraging disarmament. Regulation of armaments was included in the writing of the UN Charter in 1945 and was envisioned as a way of limiting the use of human and economic resources for their creation.[95] The advent of nuclear weapons came only weeks after the signing of the charter, resulting in the first resolution of the first General Assembly meeting calling for specific proposals for "the elimination from national armaments of atomic weapons and of all other major weapons adaptable to mass destruction".[139] The UN has been involved with arms-limitation treaties, such as the Outer Space Treaty (1967), the Treaty on the Non-Proliferation of Nuclear Weapons (1968), the Seabed Arms Control Treaty (1971), the Biological Weapons Convention (1972), the Chemical Weapons Convention (1992), and the Ottawa Treaty (1997), which prohibits landmines.[140] Three UN bodies oversee arms proliferation issues: the International Atomic Energy Agency, the Organisation for the Prohibition of Chemical Weapons, and the Comprehensive Nuclear-Test-Ban Treaty Organization Preparatory Commission.[141]
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One of the UN's primary purposes is "promoting and encouraging respect for human rights and for fundamental freedoms for all without distinction as to race, sex, language, or religion", and member states pledge to undertake "joint and separate action" to protect these rights.[114][142]
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In 1948, the General Assembly adopted a Universal Declaration of Human Rights, drafted by a committee headed by American diplomat and activist Eleanor Roosevelt, and including the French lawyer René Cassin. The document proclaims basic civil, political, and economic rights common to all human beings, though its effectiveness towards achieving these ends has been disputed since its drafting.[143] The Declaration serves as a "common standard of achievement for all peoples and all nations" rather than a legally binding document, but it has become the basis of two binding treaties, the 1966 International Covenant on Civil and Political Rights and International Covenant on Economic, Social and Cultural Rights.[144] In practice, the UN is unable to take significant action against human rights abuses without a Security Council resolution, though it does substantial work in investigating and reporting abuses.[145]
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In 1979, the General Assembly adopted the Convention on the Elimination of All Forms of Discrimination against Women, followed by the Convention on the Rights of the Child in 1989.[146] With the end of the Cold War, the push for human rights action took on new impetus.[147] The United Nations Commission on Human Rights was formed in 1993 to oversee human rights issues for the UN, following the recommendation of that year's World Conference on Human Rights. Jacques Fomerand, a scholar of the UN, describes this organization's mandate as "broad and vague", with only "meagre" resources to carry it out.[148] In 2006, it was replaced by a Human Rights Council consisting of 47 nations.[149] Also in 2006, the General Assembly passed a Declaration on the Rights of Indigenous Peoples,[150] and in 2011 it passed its first resolution recognizing the rights of LGBT people.[151]
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Other UN bodies responsible for women's rights issues include United Nations Commission on the Status of Women, a commission of ECOSOC founded in 1946; the United Nations Development Fund for Women, created in 1976; and the United Nations International Research and Training Institute for the Advancement of Women, founded in 1979.[152] The UN Permanent Forum on Indigenous Issues, one of three bodies with a mandate to oversee issues related to indigenous peoples, held its first session in 2002.[153]
|
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Millennium Development Goals[154]
|
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Another primary purpose of the UN is "to achieve international cooperation in solving international problems of an economic, social, cultural, or humanitarian character".[142] Numerous bodies have been created to work towards this goal, primarily under the authority of the General Assembly and ECOSOC.[155] In 2000, the 192 UN member states agreed to achieve eight Millennium Development Goals by 2015.[156] The Sustainable Development Goals were launched in 2015 to succeed the Millennium Development Goals.[71] The SDGs have an associated financing framework called the Addis Ababa Action Agenda.
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The UN Development Programme (UNDP), an organization for grant-based technical assistance founded in 1945, is one of the leading bodies in the field of international development. The organization also publishes the UN Human Development Index, a comparative measure ranking countries by poverty, literacy, education, life expectancy, and other factors.[157][158] The Food and Agriculture Organization (FAO), also founded in 1945, promotes agricultural development and food security.[159] UNICEF (the United Nations Children's Fund) was created in 1946 to aid European children after the Second World War and expanded its mission to provide aid around the world and to uphold the convention on the Rights of the Child.[160][161]
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The World Bank Group and International Monetary Fund (IMF) are independent, specialized agencies and observers within the UN framework, according to a 1947 agreement. They were initially formed separately from the UN through the Bretton Woods Agreement in 1944.[162] The World Bank provides loans for international development, while the IMF promotes international economic co-operation and gives emergency loans to indebted countries.[163]
|
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The World Health Organization (WHO), which focuses on international health issues and disease eradication, is another of the UN's largest agencies. In 1980, the agency announced that the eradication of smallpox had been completed. In subsequent decades, WHO largely eradicated polio, river blindness, and leprosy.[164] The Joint United Nations Programme on HIV/AIDS (UNAIDS), begun in 1996, co-ordinates the organization's response to the AIDS epidemic.[165] The UN Population Fund, which also dedicates part of its resources to combating HIV, is the world's largest source of funding for reproductive health and family planning services.[166]
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Along with the International Red Cross and Red Crescent Movement, the UN often takes a leading role in co-ordinating emergency relief.[167] The World Food Programme (WFP), created in 1961, provides food aid in response to famine, natural disasters, and armed conflict. The organization reports that it feeds an average of 90 million people in 80 nations each year.[167][168] The Office of the United Nations High Commissioner for Refugees (UNHCR), established in 1950, works to protect the rights of refugees, asylum seekers, and stateless people.[169] UNHCR and WFP programmes are funded by voluntary contributions from governments, corporations, and individuals, though the UNHCR's administrative costs are paid for by the UN's primary budget.[170]
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Since the UN's creation, over 80 colonies have attained independence. The General Assembly adopted the Declaration on the Granting of Independence to Colonial Countries and Peoples in 1960 with no votes against but abstentions from all major colonial powers. The UN works towards decolonization through groups including the UN Committee on Decolonization, created in 1962.[171] The committee lists seventeen remaining "Non-Self-Governing Territories", the largest and most populous of which is Western Sahara.[172]
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Beginning with the formation of the UN Environmental Programme (UNEP) in 1972, the UN has made environmental issues a prominent part of its agenda. A lack of success in the first two decades of UN work in this area led to the 1992 Earth Summit in Rio de Janeiro, Brazil, which sought to give new impetus to these efforts.[173] In 1988, the UNEP and the World Meteorological Organization (WMO), another UN organization, established the Intergovernmental Panel on Climate Change, which assesses and reports on research on global warming.[174] The UN-sponsored Kyoto Protocol, signed in 1997, set legally binding emissions reduction targets for ratifying states.[175]
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The UN also declares and co-ordinates international observances, periods of time to observe issues of international interest or concern. Examples include World Tuberculosis Day, Earth Day, and the International Year of Deserts and Desertification.[176]
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The UN is financed from assessed and voluntary contributions from member states. The General Assembly approves the regular budget and determines the assessment for each member. This is broadly based on the relative capacity of each country to pay, as measured by its gross national income (GNI), with adjustments for external debt and low per capita income.[178] The two-year budget for 2012–13 was $5.512 billion in total.[179]
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The Assembly has established the principle that the UN should not be unduly dependent on any one member to finance its operations. Thus, there is a "ceiling" rate, setting the maximum amount that any member can be assessed for the regular budget. In December 2000, the Assembly revised the scale of assessments in response to pressure from the United States. As part of that revision, the regular budget ceiling was reduced from 25% to 22%.[180] For the least developed countries (LDCs), a ceiling rate of 0.01% is applied.[178] In addition to the ceiling rates, the minimum amount assessed to any member nation (or "floor" rate) is set at 0.001% of the UN budget ($55,120 for the two year budget 2013–2014).[181]
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A large share of the UN's expenditure addresses its core mission of peace and security, and this budget is assessed separately from the main organizational budget.[182] The peacekeeping budget for the 2015–16 fiscal year was $8.27 billion, supporting 82,318 troops deployed in 15 missions around the world.[125] UN peace operations are funded by assessments, using a formula derived from the regular funding scale that includes a weighted surcharge for the five permanent Security Council members, who must approve all peacekeeping operations. This surcharge serves to offset discounted peacekeeping assessment rates for less developed countries.
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the largest contributors for the UN peacekeeping financial operations for the period 2019–2021 are : the United States 27.89% China 15.21%, Japan 8.56%, Germany 6.09% , the United Kingdom 5.78%, France 5.61%, Italy 3.30% and the Russian Federation 3.04%.[183]
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Special UN programmes not included in the regular budget, such as UNICEF and the World Food Programme, are financed by voluntary contributions from member governments, corporations, and private individuals.[184][185]
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A number of agencies and individuals associated with the UN have won the Nobel Peace Prize in recognition of their work. Two Secretaries-General, Dag Hammarskjöld, and Kofi Annan, were each awarded the prize (in 1961 and 2001, respectively), as were Ralph Bunche (1950), a UN negotiator, René Cassin (1968), a contributor to the Universal Declaration of Human Rights, and the US Secretary of State Cordell Hull (1945), the latter for his role in the organization's founding. Lester B. Pearson, the Canadian Secretary of State for External Affairs, was awarded the prize in 1957 for his role in organizing the UN's first peacekeeping force to resolve the Suez Crisis. UNICEF won the prize in 1965, the International Labour Organization in 1969, the UN Peace-Keeping Forces in 1988, the International Atomic Energy Agency (which reports to the UN) in 2005, and the UN-supported Organisation for the Prohibition of Chemical Weapons in 2013. The UN High Commissioner for Refugees was awarded in 1954 and 1981, becoming one of only two recipients to win the prize twice. The UN as a whole was awarded the prize in 2001, sharing it with Annan.[186] In 2007, IPCC received the prize "for their efforts to build up and disseminate greater knowledge about man-made climate change, and to lay the foundations for the measures that are needed to counteract such change."[187]
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Since its founding, there have been many calls for reform of the UN but little consensus on how to do so. Some want the UN to play a greater or more effective role in world affairs, while others want its role reduced to humanitarian work. There have also been numerous calls for the UN Security Council's membership to be increased, for different ways of electing the UN's Secretary-General, and for a UN Parliamentary Assembly. Jacques Fomerand states the most enduring divide in views of the UN is "the North–South split" between richer Northern nations and developing Southern nations. Southern nations tend to favour a more empowered UN with a stronger General Assembly, allowing them a greater voice in world affairs, while Northern nations prefer an economically laissez-faire UN that focuses on transnational threats such as terrorism.[188]
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After World War II, the French Committee of National Liberation was late to be recognized by the US as the government of France, and so the country was initially excluded from the conferences that created the new organization. The future French president Charles de Gaulle criticized the UN, famously calling it a machin ("contraption"), and was not convinced that a global security alliance would help maintain world peace, preferring direct defence treaties between countries.[189] Throughout the Cold War, both the US and USSR repeatedly accused the UN of favouring the other. In 1953, the USSR effectively forced the resignation of Trygve Lie, the Secretary-General, through its refusal to deal with him, while in the 1950s and 1960s, a popular US bumper sticker read, "You can't spell communism without U.N."[190] In a sometimes-misquoted statement, President George W. Bush stated in February 2003 (referring to UN uncertainty towards Iraqi provocations under the Saddam Hussein regime) that "free nations will not allow the UN to fade into history as an ineffective, irrelevant debating society."[191][192][193] In contrast, the French President, François Hollande, stated in 2012 that "France trusts the United Nations. She knows that no state, no matter how powerful, can solve urgent problems, fight for development and bring an end to all crises ... France wants the UN to be the centre of global governance."[194] Critics such as Dore Gold, an Israeli diplomat, Robert S. Wistrich, a British scholar, Alan Dershowitz, an American legal scholar, Mark Dreyfus, an Australian politician, and the Anti-Defamation League consider UN attention to Israel's treatment of Palestinians to be excessive.[195] In September 2015, Saudi Arabia's Faisal bin Hassan Trad has been elected Chair of the UN Human Rights Council panel that appoints independent experts,[196] a move criticized by human rights groups.[197][198]
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Since 1971, the Republic of China on Taiwan has been excluded from the UN and since then has always been rejected in new applications. Taiwanese citizens are also not allowed to enter the buildings of the United Nations with ROC passports. In this way, critics agree that the UN is failing its own development goals and guidelines. This criticism also brought pressure from the People's Republic of China, which regards the territories administered by the ROC as their own territory.[199][200]
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Critics have also accused the UN of bureaucratic inefficiency, waste, and corruption. In 1976, the General Assembly established the Joint Inspection Unit to seek out inefficiencies within the UN system. During the 1990s, the US withheld dues citing inefficiency and only started repayment on the condition that a major reforms initiative be introduced. In 1994, the Office of Internal Oversight Services (OIOS) was established by the General Assembly to serve as an efficiency watchdog.[201] In 1994, former Special Representative of the Secretary-General of the UN to Somalia Mohamed Sahnoun published "Somalia: The Missed Opportunities",[202] a book in which he analyses the reasons for the failure of the 1992 UN intervention in Somalia, showing that, between the start of the Somali civil war in 1988 and the fall of the Siad Barre regime in January 1991, the UN missed at least three opportunities to prevent major human tragedies; when the UN tried to provide humanitarian assistance, they were totally outperformed by NGOs, whose competence and dedication sharply contrasted with the UN's excessive caution and bureaucratic inefficiencies. If radical reform were not undertaken, warned Mohamed Sahnoun, then the UN would continue to respond to such crises with inept improvisation.[203] In 2004, the UN faced accusations that its recently ended Oil-for-Food Programme — in which Iraq had been allowed to trade oil for basic needs to relieve the pressure of sanctions — had suffered from widespread corruption, including billions of dollars of kickbacks. An independent inquiry created by the UN found that many of its officials had been involved, as well as raising "significant" questions about the role of Kojo Annan, the son of Kofi Annan.[204]
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In evaluating the UN as a whole, Jacques Fomerand writes that the "accomplishments of the United Nations in the last 60 years are impressive in their own terms. Progress in human development during the 20th century has been dramatic, and the UN and its agencies have certainly helped the world become a more hospitable and livable place for millions."[205] Evaluating the first 50 years of the UN's history, the author Stanley Meisler writes that "the United Nations never fulfilled the hopes of its founders, but it accomplished a great deal nevertheless", citing its role in decolonization and its many successful peacekeeping efforts.[206] The British historian Paul Kennedy states that while the organization has suffered some major setbacks, "when all its aspects are considered, the UN has brought great benefits to our generation and ... will bring benefits to our children's and grandchildren's generations as well."[207]
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Core features of the UN apparatus, such as the veto privileges of some nations in the Security Council, are often described as fundamentally undemocratic, contrary to the UN mission, and as a main cause of inaction on genocides and crimes against humanity.[208][209]
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The United Nations has inspired the extracurricular activity Model United Nations (MUN). MUN is a simulation of United Nations activity based on the UN agenda and following UN procedure. MUN is usually attended by high school and university students who organize conferences to simulate the various UN committees to discuss important issues of the day.[210] Today Model United Nations educates tens of thousands on United Nations activity around the world. Model United Nations has many famous and notable alumni, such as former Secretary-General of the United Nations Ban Ki-moon.[211]
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1 |
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The United Nations Educational, Scientific and Cultural Organisation (UNESCO;[1] French: Organisation des Nations unies pour l'éducation, la science et la culture) is a specialised agency of the United Nations (UN) aimed at contributing "to the building of peace, the eradication of poverty, sustainable development and intercultural dialogue through education, the sciences, culture, communication and information.[2][3] It is the successor of the League of Nations' International Committee on Intellectual Cooperation.[4]
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5 |
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UNESCO has 193 member states and 11 associate members.[5] Based in Paris, France, most of its field offices are "cluster" offices that cover three or more countries; national and regional offices also exist.
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6 |
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7 |
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UNESCO seeks to build a culture of peace and inclusive knowledge societies through information and communication.[6] To that end, it pursues its objectives through five major program areas: education, natural sciences, social/human sciences, culture and communication/information. It sponsors projects related to literacy, technical training, education, the advancement of science, promoting independent media and freedom of the press, preserving regional and cultural history, and promoting cultural diversity. UNESCO assists in translating and disseminating world literature, establishing international cooperation agreements to secure "World Heritage Sites" of cultural and natural importance, preserving human rights, and bridging the worldwide digital divide. It also launched and leads the Education For All movement and lifelong learning.
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9 |
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UNESCO is a member of the United Nations Development Group,[7] a coalition of UN agencies and organisations aimed at fulfilling the Sustainable Development Goals.
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11 |
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UNESCO and its mandate for international cooperation can be traced back to a League of Nations resolution on 21 September 1921, to elect a Commission to study feasibility.[8][9] This new body, the International Committee on Intellectual Cooperation (ICIC) was created in 1922. On 18 December 1925, the International Bureau of Education (IBE) began work as a non-governmental organization in the service of international educational development.[10] However, the onset of World War II largely interrupted the work of these predecessor organizations.[citation needed]
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After the signing of the Atlantic Charter and the Declaration of the United Nations, the Conference of Allied Ministers of Education (CAME) began meetings in London which continued from 16 November 1942 to 5 December 1945. On 30 October 1943, the necessity for an international organization was expressed in the Moscow Declaration, agreed upon by China, the United Kingdom, the United States and the USSR. This was followed by the Dumbarton Oaks Conference proposals of 9 October 1944. Upon the proposal of CAME and in accordance with the recommendations of the United Nations Conference on International Organization (UNCIO), held in San Francisco in April–June 1945, a United Nations Conference for the establishment of an educational and cultural organization (ECO/CONF) was convened in London 1–16 November 1945 with 44 governments represented. The idea of UNESCO was largely developed by Rab Butler, the Minister of Education for the United Kingdom, who had a great deal of influence in its development.[11] At the ECO/CONF, the Constitution of UNESCO was introduced and signed by 37 countries, and a Preparatory Commission was established.[12] The Preparatory Commission operated between 16 November 1945, and 4 November 1946—the date when UNESCO's Constitution came into force with the deposit of the twentieth ratification by a member state.[13]
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15 |
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The first General Conference took place from 19 November to 10 December 1946, and elected Dr. Julian Huxley to Director-General.[14] The Constitution was amended in November 1954 when the General Conference resolved that members of the Executive Board would be representatives of the governments of the States of which they are nationals and would not, as before, act in their personal capacity.[15] This change in governance distinguished UNESCO from its predecessor, the ICIC, in how member states would work together in the organization's fields of competence. As member states worked together over time to realize UNESCO's mandate, political and historical factors have shaped the organization's operations in particular during the Cold War, the decolonization process, and the dissolution of the USSR.[citation needed]
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17 |
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Among the major achievements of the organization is its work against racism, for example through influential statements on race starting with a declaration of anthropologists (among them was Claude Lévi-Strauss) and other scientists in 1950[16] and concluding with the 1978 Declaration on Race and Racial Prejudice.[17]
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In 1956, the Republic of South Africa withdrew from UNESCO saying that some of the organization's publications amounted to "interference" in the country's "racial problems".[18] South Africa rejoined the organization in 1994 under the leadership of Nelson Mandela.[citation needed]
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19 |
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20 |
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UNESCO's early work in the field of education included the pilot project on fundamental education in the Marbial Valley, Haiti, started in 1947.[19]
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21 |
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This project was followed by expert missions to other countries, including, for example, a mission to Afghanistan in 1949.[20]
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22 |
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In 1948, UNESCO recommended that Member States should make free primary education compulsory and universal.[21]
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23 |
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In 1990, the World Conference on Education for All, in Jomtien, Thailand, launched a global movement to provide basic education for all children, youths and adults.[22]
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24 |
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Ten years later, the 2000 World Education Forum held in Dakar, Senegal, led member governments to commit to achieving basic education for all by 2015.[23]
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25 |
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26 |
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UNESCO's early activities in culture included the Nubia Campaign, launched in 1960.[24]
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27 |
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The purpose of the campaign was to move the Great Temple of Abu Simbel to keep it from being swamped by the Nile after the construction of the Aswan Dam. During the 20-year campaign, 22 monuments and architectural complexes were relocated. This was the first and largest in a series of campaigns including Mohenjo-daro (Pakistan), Fes (Morocco), Kathmandu (Nepal), Borobudur (Indonesia) and the Acropolis (Greece).[citation needed]
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28 |
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The organization's work on heritage led to the adoption, in 1972, of the Convention concerning the Protection of the World Cultural and Natural Heritage.[25]
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29 |
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The World Heritage Committee was established in 1976 and the first sites inscribed on the World Heritage List in 1978.[26]
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30 |
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Since then important legal instruments on cultural heritage and diversity have been adopted by UNESCO member states in 2003 (Convention for the Safeguarding of the Intangible Cultural Heritage[27]) and 2005 (Convention on the Protection and Promotion of the Diversity of Cultural Expressions[28]).
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31 |
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|
32 |
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An intergovernmental meeting of UNESCO in Paris in December 1951 led to the creation of the European Council for Nuclear Research, which was responsible for establishing the European Organization for Nuclear Research (CERN)[29] later on, in 1954.
|
33 |
+
|
34 |
+
Arid Zone programming, 1948–1966, is another example of an early major UNESCO project in the field of natural sciences.[30]
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35 |
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In 1968, UNESCO organized the first intergovernmental conference aimed at reconciling the environment and development, a problem that continues to be addressed in the field of sustainable development. The main outcome of the 1968 conference was the creation of UNESCO's Man and the Biosphere Programme.[31]
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36 |
+
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37 |
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In the field of communication, the "free flow of ideas by word and image" has been in UNESCO's constitution from its beginnings, following the experience of the Second World War when control of information was a factor in indoctrinating populations for aggression.[32]
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38 |
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In the years immediately following World War II, efforts were concentrated on reconstruction and on the identification of needs for means of mass communication around the world. UNESCO started organizing training and education for journalists in the 1950s.[33]
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39 |
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In response to calls for a "New World Information and Communication Order" in the late 1970s, UNESCO established the International Commission for the Study of Communication Problems,[34] which produced the 1980 MacBride report (named after the Chair of the Commission, the Nobel Peace Prize laureate Seán MacBride).[35]
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40 |
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The same year, UNESCO created the International Programme for the Development of Communication (IPDC), a multilateral forum designed to promote media development in developing countries.[36][37] In 1991, UNESCO's General Conference endorsed the Windhoek Declaration on media independence and pluralism, which led the UN General Assembly to declare the date of its adoption, 3 May, as World Press Freedom Day.[38] Since 1997, UNESCO has awarded the UNESCO / Guillermo Cano World Press Freedom Prize every 3 May. In the lead up to the World Summit on the Information Society in 2003 (Geneva) and 2005 (Tunis), UNESCO introduced the Information for All Programme.[citation needed]
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41 |
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|
42 |
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UNESCO admitted Palestine as a member in 2011.[39][40]
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43 |
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Laws passed in the United States after Palestine applied for UNESCO and WHO membership in April 1989[41][42] mean that the US cannot contribute financially to any UN organisation that accepts Palestine as a full member.[43][44] As a result, the US withdrew its funding, which had accounted for about 22% of UNESCO's budget.[45] Israel also reacted to Palestine's admittance to UNESCO by freezing Israeli payments to UNESCO and imposing sanctions on the Palestinian Authority,[46] stating that Palestine's admittance would be detrimental "to potential peace talks".[47] Two years after they stopped paying their dues to UNESCO, the US and Israel lost UNESCO voting rights in 2013 without losing the right to be elected; thus, the US was elected as a member of the Executive Board for the period 2016–19.[48] In 2019, Israel left UNESCO after 69 years of membership, with Israel's ambassador to the UN Danny Danon writing: "UNESCO is the body that continually rewrites history, including by erasing the Jewish connection to Jerusalem... it is corrupted and manipulated by Israel's enemies... we are not going to be a member of an organisation that deliberately acts against us".[49]
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44 |
+
|
45 |
+
UNESCO implements its activities through the five program areas: education, natural sciences, social and human sciences, culture, and communication and information.
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46 |
+
|
47 |
+
UNESCO does not accredit institutions of higher learning.[50]
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48 |
+
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49 |
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The UNESCO transparency portal has been designed to enable public access to information regarding Organization's activities, such as its aggregate budget for a biennium, as well as links to relevant programmatic and financial documents. These two distinct sets of information are published on the IATI registry, respectively based on the IATI Activity Standard and the IATI Organization Standard.
|
50 |
+
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51 |
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There have been proposals to establish two new UNESCO lists. The first proposed list will focus on movable cultural heritage such as artifacts, paintings, and biofacts. The list may include cultural objects, such as the Jōmon Venus of Japan, the Mona Lisa of France, the Gebel el-Arak Knife of Egypt, The Ninth Wave of Russia, the Seated Woman of Çatalhöyük of Turkey, the David (Michelangelo) of Italy, the Mathura Herakles of India, the Manunggul Jar of the Philippines, the Crown of Baekje of South Korea, The Hay Wain of the United Kingdom and the Benin Bronzes of Nigeria. The second proposed list will focus on the world's living species, such as the komodo dragon of Indonesia, the panda of China, the bald eagle of North American countries, the aye-aye of Madagascar, the Asiatic lion of India, the kakapo of New Zealand, and the mountain tapir of Colombia, Ecuador and Peru.[67][68]
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52 |
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53 |
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UNESCO and its specialized institutions issue a number of magazines.
|
54 |
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55 |
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The UNESCO Courier magazine states its mission to "promote UNESCO's ideals, maintain a platform for the dialogue between cultures and provide a forum for international debate". Since March 2006 it is available online, with limited printed issues. Its articles express the opinions of the authors which are not necessarily the opinions of UNESCO. There was a hiatus in publishing between 2012 and 2017.[69]
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56 |
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57 |
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In 1950, UNESCO initiated the quarterly review Impact of Science on Society (also known as Impact) to discuss the influence of science on society. The journal ceased publication in 1992.[70] UNESCO also published Museum International Quarterly from the year 1948.
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58 |
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59 |
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UNESCO has official relations with 322 international non-governmental organizations (NGOs).[71] Most of these are what UNESCO calls "operational"; a select few are "formal".[72]
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60 |
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The highest form of affiliation to UNESCO is "formal associate", and the 22 NGOs[73] with formal associate (ASC) relations occupying offices at UNESCO are:
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61 |
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62 |
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The institutes are specialized departments of the organization that support UNESCO's programme, providing specialized support for cluster and national offices.
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63 |
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64 |
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UNESCO awards 22 prizes[85] in education, science, culture and peace:
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65 |
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66 |
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International Days observed at UNESCO is provided in the table given below[86]
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67 |
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68 |
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As of January 2019, UNESCO has 193 member states and 11 associate members.[88] Some members are not independent states and some members have additional National Organizing Committees from some of their dependent territories.[89] UNESCO state parties are the United Nations member states (except Liechtenstein, United States[90] and Israel[91]), as well as Cook Islands, Niue and Palestine.[92][93] The United States and Israel left UNESCO on 31 December 2018.[94]
|
69 |
+
|
70 |
+
There has been no elected UNESCO Director-General from Southeast Asia, South Asia, Central and North Asia, Middle East, North Africa, East Africa, Central Africa, South Africa, Australia-Oceania, and South America since inception.
|
71 |
+
|
72 |
+
The Directors-General of UNESCO came from West Europe (5), Central America (1), North America (2), West Africa (1), East Asia (1), and East Europe (1). Out of the 11 Directors-General since inception, women have held the position only twice. Qatar, the Philippines, and Iran are proposing for a Director-General bid by 2021 or 2025. There have never been a Middle Eastern or Southeast Asian UNESCO Director-General since inception. The ASEAN bloc and some Pacific and Latin American nations support the possible bid of the Philippines, which is culturally Asian, Oceanic, and Latin. Qatar and Iran, on the other hand, have fragmented support in the Middle East. Egypt, Israel, and Madagascar are also vying for the position but have yet to express a direct or indirect proposal. Both Qatar and Egypt lost in the 2017 bid against France.
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73 |
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The list of the Directors-General of UNESCO since its establishment in 1946 is as follows:[95]
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75 |
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|
76 |
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This is the list of the sessions of the UNESCO General Conference held since 1946:[96]
|
77 |
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|
78 |
+
The 40th General Conference On 12th To 27 November 2019.
|
79 |
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|
80 |
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Italy
|
81 |
+
|
82 |
+
Netherlands
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83 |
+
|
84 |
+
Spain
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85 |
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|
86 |
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Switzerland
|
87 |
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|
88 |
+
Poland
|
89 |
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|
90 |
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Russia
|
91 |
+
|
92 |
+
Serbia
|
93 |
+
|
94 |
+
Brazil
|
95 |
+
|
96 |
+
Dominican Republic
|
97 |
+
|
98 |
+
Uruguay
|
99 |
+
|
100 |
+
Kyrgyzstan
|
101 |
+
|
102 |
+
Philippines
|
103 |
+
|
104 |
+
Pakistan
|
105 |
+
|
106 |
+
South Korea
|
107 |
+
|
108 |
+
Thailand
|
109 |
+
|
110 |
+
Congo
|
111 |
+
|
112 |
+
Guinea
|
113 |
+
|
114 |
+
Ghana
|
115 |
+
|
116 |
+
Kenya
|
117 |
+
|
118 |
+
Namibia
|
119 |
+
|
120 |
+
Senegal
|
121 |
+
|
122 |
+
Togo
|
123 |
+
|
124 |
+
UAE
|
125 |
+
|
126 |
+
Tunisia
|
127 |
+
|
128 |
+
France
|
129 |
+
Greece
|
130 |
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Italy
|
131 |
+
Spain
|
132 |
+
United Kingdom
|
133 |
+
|
134 |
+
Lithuania
|
135 |
+
Russia
|
136 |
+
Serbia
|
137 |
+
Slovenia
|
138 |
+
|
139 |
+
Brazil
|
140 |
+
Haiti
|
141 |
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Mexico
|
142 |
+
Nicaragua
|
143 |
+
Paraguay
|
144 |
+
|
145 |
+
India
|
146 |
+
Iran
|
147 |
+
Malaysia
|
148 |
+
Pakistan
|
149 |
+
South Korea
|
150 |
+
Sri Lanka
|
151 |
+
Vietnam
|
152 |
+
|
153 |
+
Cameroon
|
154 |
+
Ivory Coast
|
155 |
+
Ghana
|
156 |
+
Kenya
|
157 |
+
Nigeria
|
158 |
+
Senegal
|
159 |
+
South Africa
|
160 |
+
|
161 |
+
Lebanon
|
162 |
+
Oman
|
163 |
+
Qatar
|
164 |
+
Sudan
|
165 |
+
|
166 |
+
Germany
|
167 |
+
Netherlands
|
168 |
+
Sweden
|
169 |
+
|
170 |
+
Albania
|
171 |
+
Estonia
|
172 |
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Ukraine
|
173 |
+
|
174 |
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Argentina
|
175 |
+
Belize
|
176 |
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Dominican Republic
|
177 |
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El Salvador
|
178 |
+
Saint Kitts and Nevis
|
179 |
+
Trinidad and Tobago
|
180 |
+
|
181 |
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Bangladesh
|
182 |
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China
|
183 |
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India
|
184 |
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Japan
|
185 |
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Nepal
|
186 |
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Turkmenistan
|
187 |
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|
188 |
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Chad
|
189 |
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Guinea
|
190 |
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Mauritius
|
191 |
+
Mozambique
|
192 |
+
Togo
|
193 |
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Uganda
|
194 |
+
|
195 |
+
Algeria
|
196 |
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Egypt
|
197 |
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Kuwait
|
198 |
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Morocco
|
199 |
+
|
200 |
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Austria
|
201 |
+
France
|
202 |
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Italy
|
203 |
+
India
|
204 |
+
Spain
|
205 |
+
United Kingdom
|
206 |
+
United States
|
207 |
+
|
208 |
+
Czech Republic
|
209 |
+
Montenegro
|
210 |
+
North Macedonia
|
211 |
+
|
212 |
+
Russia
|
213 |
+
|
214 |
+
Brazil
|
215 |
+
Cuba
|
216 |
+
Ecuador
|
217 |
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Mexico
|
218 |
+
|
219 |
+
Afghanistan
|
220 |
+
Indonesia
|
221 |
+
Pakistan
|
222 |
+
Papua New Guinea
|
223 |
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South Korea
|
224 |
+
Thailand
|
225 |
+
|
226 |
+
Angola
|
227 |
+
Ethiopia
|
228 |
+
Gabon
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229 |
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Gambia
|
230 |
+
Malawi
|
231 |
+
Mali
|
232 |
+
Namibia
|
233 |
+
Nigeria
|
234 |
+
|
235 |
+
Jordan
|
236 |
+
Tunisia
|
237 |
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United Arab Emirates
|
238 |
+
|
239 |
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UNESCO headquarters are located at Place de Fontenoy in Paris, France.
|
240 |
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|
241 |
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UNESCO's field offices across the globe are categorized into four primary office types based upon their function and geographic coverage: cluster offices, national offices, regional bureaus and liaison offices.
|
242 |
+
|
243 |
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The following list of all UNESCO Field Offices is organized geographically by UNESCO Region and identifies the members states and associate members of UNESCO which are served by each office.[103]
|
244 |
+
|
245 |
+
UNESCO has been the centre of controversy in the past, particularly in its relationships with the United States, the United Kingdom, Singapore and the former Soviet Union. During the 1970s and 1980s, UNESCO's support for a "New World Information and Communication Order" and its MacBride report calling for democratization of the media and more egalitarian access to information was condemned in these countries as attempts to curb freedom of the press. UNESCO was perceived as a platform for communists and Third World dictators to attack the West, in contrast to accusations made by the USSR in the late 1940s and early 1950s.[107] In 1984, the United States withheld its contributions and withdrew from the organization in protest, followed by the United Kingdom in 1985.[108] Singapore withdrew also at the end of 1985, citing rising membership fees.[109] Following a change of government in 1997, the UK rejoined. The United States rejoined in 2003, followed by Singapore on 8 October 2007.[110]
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246 |
+
|
247 |
+
Israel was admitted to UNESCO in 1949, one year after its creation. Israel has maintained its membership since 1949.
|
248 |
+
In 2010, Israel designated the Cave of the Patriarchs, Hebron and Rachel's Tomb, Bethlehem as National Heritage Sites and announced restoration work, prompting criticism from the Obama administration and protests from Palestinians.[111] In October 2010, UNESCO's Executive Board voted to declare the sites as "al-Haram al-Ibrahimi/Tomb of the Patriarchs" and "Bilal bin Rabah Mosque/Rachel's Tomb" and stated that they were "an integral part of the occupied Palestinian Territories" and any unilateral Israeli action was a violation of international law.[112]
|
249 |
+
UNESCO described the sites as significant to "people of the Muslim, Christian and Jewish traditions", and accused Israel of highlighting only the Jewish character of the sites.[113]
|
250 |
+
Israel in turn accused UNESCO of "detach[ing] the Nation of Israel from its heritage", and accused it of being politically motivated.[114]
|
251 |
+
The Rabbi of the Western Wall said that Rachel's tomb had not previously been declared a holy Muslim site.[115] Israel partially suspended ties with UNESCO. Israeli Deputy Foreign Minister Danny Ayalon declared that the resolution was a "part of Palestinian escalation".
|
252 |
+
Zevulun Orlev, chairman of the Knesset Education and Culture Committee, referred to the resolutions as an attempt to undermine the mission of UNESCO as a scientific and cultural organization that promotes cooperation throughout the world.[116][117]
|
253 |
+
|
254 |
+
On 28 June 2011, UNESCO's World Heritage Committee, at Jordan's insistence, censured[clarification needed] Israel's decision to demolish and rebuild the Mughrabi Gate Bridge in Jerusalem for safety reasons. Israel stated that Jordan had signed an agreement with Israel stipulating that the existing bridge must be dismantled for safety reasons; Jordan disputed the agreement, saying that it was only signed under U.S. pressure. Israel was also unable to address the UNESCO committee over objections from Egypt.[118]
|
255 |
+
|
256 |
+
In January 2014, days before it was scheduled to open, UNESCO Director-General, Irina Bokova, "indefinitely postponed" and effectively cancelled an exhibit created by the Simon Wiesenthal Center entitled "The People, The Book, The Land: The 3,500-year relationship between the Jewish people and the Land of Israel". The event was scheduled to run from 21 January through 30 January in Paris. Bokova cancelled the event after representatives of Arab states at UNESCO argued that its display would "harm the peace process".[119] The author of the exhibition, Professor Robert Wistrich of the Hebrew University's Vidal Sassoon International Center for the Study of Anti-Semitism, called the cancellation an "appalling act", and characterized Bokova's decision as "an arbitrary act of total cynicism and, really, contempt for the Jewish people and its history". UNESCO amended the decision to cancel the exhibit within the year, and it quickly achieved popularity and was viewed as a great success.[120]
|
257 |
+
|
258 |
+
On January 1, 2019, Israel formally left UNESCO in pursuance of the US withdrawal over the perceived continuous anti-Israel bias.
|
259 |
+
|
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On 13 October 2016, UNESCO passed a resolution on East Jerusalem that condemned Israel for "aggressions" by Israeli police and soldiers and "illegal measures" against the freedom of worship and Muslims' access to their holy sites, while also recognizing Israel as the occupying power. Palestinian leaders welcomed the decision.[121] While the text acknowledged the "importance of the Old City of Jerusalem and its walls for the three monotheistic religions", it referred to the sacred hilltop compound in Jerusalem's Old City only by its Muslim name "Al-Haram al-Sharif", Arabic for Noble Sanctuary. In response, Israel denounced the UNESCO resolution for its omission of the words "Temple Mount" or "Har HaBayit", stating that it denies Jewish ties to the key holy site.[121][122] After receiving criticism from numerous Israeli politicians and diplomats, including Benjamin Netanyahu and Ayelet Shaked, Israel froze all ties with the organization.[123][124] The resolution was condemned by Ban Ki-moon and the Director-General of UNESCO, Irina Bokova, who said that Judaism, Islam and Christianity have clear historical connections to Jerusalem and "to deny, conceal or erase any of the Jewish, Christian or Muslim traditions undermines the integrity of the site.[125][126] "Al-Aqsa Mosque [or] Al-Haram al-Sharif" is also Temple Mount, whose Western Wall is the holiest place in Judaism."[127]. It was also rejected by the Czech Parliament which said the resolution reflects a "hateful anti-Israel sentiment",[128] and hundreds of Italian Jews demonstrated in Rome over Italy's abstention.[128] On 26 October, UNESCO approved a reviewed version of the resolution, which also criticized Israel for its continuous "refusal to let the body's experts access Jerusalem's holy sites to determine their conservation status".[129] Despite containing some softening of language following Israeli protests over a previous version, Israel continued to denounce the text.[130] The resolution refers to the site Jews and Christians refer to as the Temple Mount, or Har HaBayit in Hebrew, only by its Arab name — a significant semantic decision also adopted by UNESCO's executive board, triggering condemnation from Israel and its allies. U.S. Ambassador Crystal Nix Hines stated: "This item should have been defeated. These politicized and one-sided resolutions are damaging the credibility of UNESCO."[131]
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In October 2017, the United States and Israel announced they would withdraw from the organization, citing in-part anti-Israel bias.[132][133]
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In February 2011, an article was published in a Palestinian youth magazine in which a teenage girl described one of her four role-models as Adolf Hitler. In December 2011, UNESCO, which partly funded the magazine, condemned the material and subsequently withdrew support.[134]
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In 2012, UNESCO decided to establish a chair at the Islamic University of Gaza in the field of astronomy, astrophysics, and space sciences,[135] fueling controversy and criticism. Israel bombed the school in 2008 stating that they develop and store weapons there, which Israel restated in criticizing UNESCO's move.[136][137]
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The head, Kamalain Shaath, defended UNESCO, stating that "the Islamic University is a purely academic university that is interested only in education and its development".[138][139][140] Israeli ambassador to UNESCO Nimrod Barkan planned to submit a letter of protest with information about the university's ties to Hamas, especially angry that this was the first Palestinian university that UNESCO chose to cooperate with.[141] The Jewish organization B'nai B'rith criticized the move as well.[142]
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In 2013, UNESCO announced that the collection "The Life and Works of Ernesto Che Guevara" became part of the Memory of the World Register. US Congresswoman Ileana Ros-Lehtinen condemned this decision, saying that the organization acts against its own ideals:[143]
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This decision is more than an insult to the families of those Cubans who were lined up and summarily executed by Che and his merciless cronies but it also serves as a direct contradiction to the UNESCO ideals of encouraging peace and universal respect for human rights.
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UN Watch also condemned this selection by UNESCO.[144]
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In 2015, Japan threatened to halt funding for UNESCO over the organization's decision to include documents relating to the 1937 Nanjing massacre in the latest listing for its "Memory of the World" program.[145]
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In October 2016, Japanese Foreign Minister Fumio Kishida confirmed that Japan's 2016 annual funding of ¥4.4 billion had been suspended although denied any direct link with the Nanjing document controversy.[146]
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The United States withdrew from UNESCO in 1984, citing the "highly politicized" nature of the organisation, its ostensible "hostility toward the basic institutions of a free society, especially a free market and a free press", as well as its "unrestrained budgetary expansion", and poor management under then Director General Amadou-Mahter M'Bow of Senegal.[147]
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On 19 September 1989, former U.S. Congressman Jim Leach stated before a Congressional subcommittee:[148]
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The reasons for the withdrawal of the United States from UNESCO in 1984 are well-known; my view is that we overreacted to the calls of some who wanted to radicalize UNESCO, and the calls of others who wanted the United States to lead in emasculating the UN system. The fact is UNESCO is one of the least dangerous international institutions ever created. While some member countries within UNESCO attempted to push journalistic views antithetical to the values of the west, and engage in Israel bashing, UNESCO itself never adopted such radical postures. The U.S. opted for empty-chair diplomacy, after winning, not losing, the battles we engaged in… It was nuts to get out, and would be nuttier not to rejoin.
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Leach concluded that the record showed Israel bashing, a call for a new world information order, money management, and arms control policy to be the impetus behind the withdrawal; he asserted that before departing from UNESCO, a withdrawal from the IAEA had been pushed on him.[148] On 1 October 2003, the U.S. rejoined UNESCO.[147]
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On 12 October 2017, the United States notified UNESCO that it will again withdraw from the organization on 31 December 2018 and will seek to establish a permanent observer mission beginning in 2019. The Department of State cited "mounting arrears at UNESCO, the need for fundamental reform in the organization, and continuing anti-Israel bias at UNESCO".[132] Israel praised the withdrawal decision as "brave" and "moral".[147]
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The United States has not paid over $600 million in dues[149] since it stopped paying its $80 million annual UNESCO dues when Palestine became a full member in 2011. Israel and the US were among the 14 votes against the membership out of 194 member countries.[150]
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On May 25, 2016, the noted Turkish poet and human rights activist Zülfü Livaneli resigned as Turkey's only UNESCO goodwill ambassador. He highlighted human rights situation in Turkey and destruction of historical Sur district of Diyarbakir, the largest city in Kurdish-majority southeast Turkey, during fighting between the Turkish army and Kurdish militants as the main reasons for his resignation. Livaneli said: "To pontificate on peace while remaining silent against such violations is a contradiction of the fundamental ideals of UNESCO."[151]
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In 1981, UNESCO and the UN celebrated the Atatürk Centennial.
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UNESCO develops, maintains and disseminates, free of charge, two interrelated software packages for database management (CDS/ISIS [not to be confused with UK police software package ISIS]) and data mining/statistical analysis (IDAMS).[153]
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The North Atlantic Treaty Organization (NATO, /ˈneɪtoʊ/; French: Organisation du traité de l'Atlantique nord, OTAN), also called the North Atlantic Alliance, is an intergovernmental military alliance between 30 North American and European countries. The organization implements the North Atlantic Treaty that was signed on 4 April 1949.[3][4] NATO constitutes a system of collective defence whereby its independent member states agree to mutual defence in response to an attack by any external party. NATO's Headquarters are located in Evere, Brussels, Belgium, while the headquarters of Allied Command Operations is near Mons, Belgium.
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Since its founding, the admission of new member states has increased the alliance from the original 12 countries to 30. The most recent member state to be added to NATO was North Macedonia on 27 March 2020. NATO currently recognizes Bosnia and Herzegovina, Georgia, and Ukraine as aspiring members.[5] An additional 20 countries participate in NATO's Partnership for Peace program, with 15 other countries involved in institutionalized dialogue programs. The combined military spending of all NATO members constitutes over 70% of the global total.[6] Members agreed that their aim is to reach or maintain the target defense spending of at least 2% of GDP by 2024.[7][8]
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On 4 March 1947, the Treaty of Dunkirk was signed by France and the United Kingdom as a Treaty of Alliance and Mutual Assistance in the event of a possible attack by Germany or the Soviet Union in the aftermath of World War II. In 1948, this alliance was expanded to include the Benelux countries, in the form of the Western Union, also referred to as the Brussels Treaty Organization (BTO), established by the Treaty of Brussels.[9] Talks for a new military alliance which could also include North America resulted in the signature of the North Atlantic Treaty on 4 April 1949 by the member states of the Western Union plus the United States, Canada, Portugal, Italy, Norway, Denmark and Iceland.[10]
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The North Atlantic Treaty was largely dormant until the Korean War initiated the establishment of NATO to implement it, by means of an integrated military structure: This included the formation of Supreme Headquarters Allied Powers Europe (SHAPE) in 1951, which adopted the Western Union's military structures and plans.[11] In 1952, the post of Secretary General of NATO was established as the organization's chief civilian. That year also saw the first major NATO maritime exercises, Exercise Mainbrace and the accession of Greece and Turkey to the organization.[12][13] Following the London and Paris Conferences, West Germany was permitted to rearm militarily, as they joined NATO in May 1955, which was, in turn, a major factor in the creation of the Soviet-dominated Warsaw Pact, delineating the two opposing sides of the Cold War.
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Doubts over the strength of the relationship between the European states and the United States ebbed and flowed, along with doubts over the credibility of the NATO defense against a prospective Soviet invasion – doubts that led to the development of the independent French nuclear deterrent and the withdrawal of France from NATO's military structure in 1966.[14][15] In 1982, the newly democratic Spain joined the alliance.
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The collapse of the Warsaw Pact in 1989–1991 removed the de facto main adversary of NATO and caused a strategic re-evaluation of NATO's purpose, nature, tasks, and focus on the continent of Europe. This shift started with the 1990 signing in Paris of the Treaty on Conventional Armed Forces in Europe between NATO and the Soviet Union, which mandated specific military reductions across the continent that continued after the dissolution of the Soviet Union in December 1991.[16] At that time, European countries accounted for 34 percent of NATO's military spending; by 2012, this had fallen to 21 percent.[17] NATO also began a gradual expansion to include newly autonomous Central and Eastern European nations, and extended its activities into political and humanitarian situations that had not formerly been NATO concerns.[18]
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After the fall of the Berlin Wall in Germany in 1989, the organization conducted its first military interventions in Bosnia from 1992 to 1995 and later Yugoslavia in 1999 during the break-up of Yugoslavia.[19] Politically, the organization sought better relations with former Warsaw Pact countries, most of which joined the alliance in 1999 and 2004. Article 5 of the North Atlantic treaty, requiring member states to come to the aid of any member state subject to an armed attack, was invoked for the first and only time after the September 11 attacks,[20] after which troops were deployed to Afghanistan under the NATO-led ISAF. The organization has operated a range of additional roles since then, including sending trainers to Iraq, assisting in counter-piracy operations[21] and in 2011 enforcing a no-fly zone over Libya in accordance with UN Security Council Resolution 1973. The less potent Article 4, which merely invokes consultation among NATO members, has been invoked five times following incidents in the Iraq War, Syrian Civil War, and annexation of Crimea.[22]
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The first post-Cold War expansion of NATO came with German reunification on 3 October 1990, when the former East Germany became part of the Federal Republic of Germany and the alliance. As part of post-Cold War restructuring, NATO's military structure was cut back and reorganized, with new forces such as the Headquarters Allied Command Europe Rapid Reaction Corps established. The changes brought about by the collapse of the Soviet Union on the military balance in Europe were recognized in the Adapted Conventional Armed Forces in Europe Treaty, which was signed in 1999. The policies of French President Nicolas Sarkozy resulted in a major reform of France's military position, culminating with the return to full membership on 4 April 2009, which also included France rejoining the NATO Military Command Structure, while maintaining an independent nuclear deterrent.[15][23][24]
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Between 1994 and 1997, wider forums for regional cooperation between NATO and its neighbors were set up, like the Partnership for Peace, the Mediterranean Dialogue initiative and the Euro-Atlantic Partnership Council. In 1998, the NATO–Russia Permanent Joint Council was established. Between 1999 and 2020 NATO incorporated the following Central and Eastern European countries, including several former communist states: the Czech Republic, Hungary, Poland, Bulgaria, Estonia, Latvia, Lithuania, Romania, Slovakia, Slovenia, Albania, Croatia, Montenegro, and North Macedonia.[18]
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The Russian intervention in Crimea in 2014 led to strong condemnation by NATO nations and the creation of a new "spearhead" force of 5,000 troops at bases in Estonia, Lithuania, Latvia, Poland, Romania, and Bulgaria.[25] At the subsequent 2014 Wales summit, the leaders of NATO's member states formally committed for the first time to spend the equivalent of at least 2% of their gross domestic products on defence by 2024, which had previously been only an informal guideline.[26] NATO did not condemn the 2016–present purges in Turkey.[27] NATO members have resisted the UN's Nuclear Weapon Ban Treaty, a binding agreement for negotiations for the total elimination of nuclear weapons, supported by more than 120 nations.[28]
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No military operations were conducted by NATO during the Cold War. Following the end of the Cold War, the first operations, Anchor Guard in 1990 and Ace Guard in 1991, were prompted by the Iraqi invasion of Kuwait. Airborne early warning aircraft were sent to provide coverage of southeastern Turkey, and later a quick-reaction force was deployed to the area.[29]
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The Bosnian War began in 1992, as a result of the break-up of Yugoslavia. The deteriorating situation led to United Nations Security Council Resolution 816 on 9 October 1992, ordering a no-fly zone over central Bosnia and Herzegovina, which NATO began enforcing on 12 April 1993 with Operation Deny Flight. From June 1993 until October 1996, Operation Sharp Guard added maritime enforcement of the arms embargo and economic sanctions against the Federal Republic of Yugoslavia. On 28 February 1994, NATO took its first wartime action by shooting down four Bosnian Serb aircraft violating the no-fly zone.[30]
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On 10 and 11 April 1994, the United Nations Protection Force called in air strikes to protect the Goražde safe area, resulting in the bombing of a Bosnian Serb military command outpost near Goražde by two US F-16 jets acting under NATO direction.[31] In retaliation, Serbs took 150 U.N. personnel hostage on 14 April.[32][33] On 16 April a British Sea Harrier was shot down over Goražde by Serb forces.[34]
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In August 1995, a two-week NATO bombing campaign, Operation Deliberate Force, began against the Army of the Republika Srpska, after the Srebrenica massacre.[35] Further NATO air strikes helped bring the Yugoslav wars to an end, resulting in the Dayton Agreement in November 1995.[35] As part of this agreement, NATO deployed a UN-mandated peacekeeping force, under Operation Joint Endeavor, named IFOR. Almost 60,000 NATO troops were joined by forces from non-NATO nations in this peacekeeping mission. This transitioned into the smaller SFOR, which started with 32,000 troops initially and ran from December 1996 until December 2004, when operations were then passed onto European Union Force Althea.[36] Following the lead of its member nations, NATO began to award a service medal, the NATO Medal, for these operations.[37]
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In an effort to stop Slobodan Milošević's Serbian-led crackdown on KLA separatists and Albanian civilians in Kosovo, the United Nations Security Council passed Resolution 1199 on 23 September 1998 to demand a ceasefire. Negotiations under US Special Envoy Richard Holbrooke broke down on 23 March 1999, and he handed the matter to NATO,[38] which started a 78-day bombing campaign on 24 March 1999.[39] Operation Allied Force targeted the military capabilities of what was then the Federal Republic of Yugoslavia. During the crisis, NATO also deployed one of its international reaction forces, the ACE Mobile Force (Land), to Albania as the Albania Force (AFOR), to deliver humanitarian aid to refugees from Kosovo.[40]
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Though the campaign was criticized for high civilian casualties, including bombing of the Chinese embassy in Belgrade, Milošević finally accepted the terms of an international peace plan on 3 June 1999, ending the Kosovo War. On 11 June, Milošević further accepted UN resolution 1244, under the mandate of which NATO then helped establish the KFOR peacekeeping force. Nearly one million refugees had fled Kosovo, and part of KFOR's mandate was to protect the humanitarian missions, in addition to deterring violence.[40][41] In August–September 2001, the alliance also mounted Operation Essential Harvest, a mission disarming ethnic Albanian militias in the Republic of Macedonia.[42] As of 1 December 2013[update], 4,882 KFOR soldiers, representing 31 countries, continue to operate in the area.[43]
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The US, the UK, and most other NATO countries opposed efforts to require the UN Security Council to approve NATO military strikes, such as the action against Serbia in 1999, while France and some others claimed that the alliance needed UN approval.[44] The US/UK side claimed that this would undermine the authority of the alliance, and they noted that Russia and China would have exercised their Security Council vetoes to block the strike on Yugoslavia, and could do the same in future conflicts where NATO intervention was required, thus nullifying the entire potency and purpose of the organization. Recognizing the post-Cold War military environment, NATO adopted the Alliance Strategic Concept during its Washington summit in April 1999 that emphasized conflict prevention and crisis management.[45]
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The September 11 attacks in the United States caused NATO to invoke Article 5 of the NATO Charter for the first time in the organization's history. The Article states that an attack on any member shall be considered to be an attack on all. The invocation was confirmed on 4 October 2001 when NATO determined that the attacks were indeed eligible under the terms of the North Atlantic Treaty.[46] The eight official actions taken by NATO in response to the attacks included Operation Eagle Assist and Operation Active Endeavour, a naval operation in the Mediterranean Sea designed to prevent the movement of terrorists or weapons of mass destruction, and to enhance the security of shipping in general, which began on 4 October 2001.[47]
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The alliance showed unity: On 16 April 2003, NATO agreed to take command of the International Security Assistance Force (ISAF), which included troops from 42 countries. The decision came at the request of Germany and the Netherlands, the two nations leading ISAF at the time of the agreement, and all nineteen NATO ambassadors approved it unanimously. The handover of control to NATO took place on 11 August, and marked the first time in NATO's history that it took charge of a mission outside the north Atlantic area.[48]
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ISAF was initially charged with securing Kabul and surrounding areas from the Taliban, al Qaeda and factional warlords, so as to allow for the establishment of the Afghan Transitional Administration headed by Hamid Karzai. In October 2003, the UN Security Council authorized the expansion of the ISAF mission throughout Afghanistan,[49] and ISAF subsequently expanded the mission in four main stages over the whole of the country.[50]
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On 31 July 2006, the ISAF additionally took over military operations in the south of Afghanistan from a US-led anti-terrorism coalition.[51] Due to the intensity of the fighting in the south, in 2011 France allowed a squadron of Mirage 2000 fighter/attack aircraft to be moved into the area, to Kandahar, in order to reinforce the alliance's efforts.[52] During its 2012 Chicago Summit, NATO endorsed a plan to end the Afghanistan war and to remove the NATO-led ISAF Forces by the end of December 2014.[53] ISAF was disestablished in December 2014 and replaced by the follow-on training Resolute Support Mission.[54]
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In August 2004, during the Iraq War, NATO formed the NATO Training Mission – Iraq, a training mission to assist the Iraqi security forces in conjunction with the US led MNF-I.[55] The NATO Training Mission-Iraq (NTM-I) was established at the request of the Iraqi Interim Government under the provisions of United Nations Security Council Resolution 1546. The aim of NTM-I was to assist in the development of Iraqi security forces training structures and institutions so that Iraq can build an effective and sustainable capability that addresses the needs of the nation. NTM-I was not a combat mission but is a distinct mission, under the political control of NATO's North Atlantic Council. Its operational emphasis was on training and mentoring. The activities of the mission were coordinated with Iraqi authorities and the US-led Deputy Commanding General Advising and Training, who was also dual-hatted as the Commander of NTM-I. The mission officially concluded on 17 December 2011.[56]
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Turkey invoked the first Article 4 meetings in 2003 at the start of the Iraq War. Turkey also invoked this article twice in 2012 during the Syrian Civil War, after the downing of an unarmed Turkish F-4 reconnaissance jet, and after a mortar was fired at Turkey from Syria,[57] and again in 2015 after threats by Islamic State of Iraq and the Levant to its territorial integrity.[58]
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Beginning on 17 August 2009, NATO deployed warships in an operation to protect maritime traffic in the Gulf of Aden and the Indian Ocean from Somali pirates, and help strengthen the navies and coast guards of regional states. The operation was approved by the North Atlantic Council and involves warships primarily from the United States though vessels from many other nations are also included. Operation Ocean Shield focuses on protecting the ships of Operation Allied Provider which are distributing aid as part of the World Food Programme mission in Somalia. Russia, China and South Korea have sent warships to participate in the activities as well.[59][60] The operation seeks to dissuade and interrupt pirate attacks, protect vessels, and abetting to increase the general level of security in the region.[61]
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During the Libyan Civil War, violence between protesters and the Libyan government under Colonel Muammar Gaddafi escalated, and on 17 March 2011 led to the passage of United Nations Security Council Resolution 1973, which called for a ceasefire, and authorized military action to protect civilians. A coalition that included several NATO members began enforcing a no-fly zone over Libya shortly afterwards, beginning with Opération Harmattan by the French Air Force on 19 March.
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On 20 March 2011, NATO states agreed on enforcing an arms embargo against Libya with Operation Unified Protector using ships from NATO Standing Maritime Group 1 and Standing Mine Countermeasures Group 1,[62] and additional ships and submarines from NATO members.[63] They would "monitor, report and, if needed, interdict vessels suspected of carrying illegal arms or mercenaries".[62]
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On 24 March, NATO agreed to take control of the no-fly zone from the initial coalition, while command of targeting ground units remained with the coalition's forces.[64][65] NATO began officially enforcing the UN resolution on 27 March 2011 with assistance from Qatar and the United Arab Emirates.[66] By June, reports of divisions within the alliance surfaced as only eight of the 28 member nations were participating in combat operations,[67] resulting in a confrontation between US Defense Secretary Robert Gates and countries such as Poland, Spain, the Netherlands, Turkey, and Germany to contribute more, the latter believing the organization has overstepped its mandate in the conflict.[68][69][70] In his final policy speech in Brussels on 10 June, Gates further criticized allied countries in suggesting their actions could cause the demise of NATO.[71] The German foreign ministry pointed to "a considerable [German] contribution to NATO and NATO-led operations" and to the fact that this engagement was highly valued by President Obama.[72]
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While the mission was extended into September, Norway that day announced it would begin scaling down contributions and complete withdrawal by 1 August.[73] Earlier that week it was reported Danish air fighters were running out of bombs.[74][75] The following week, the head of the Royal Navy said the country's operations in the conflict were not sustainable.[76] By the end of the mission in October 2011, after the death of Colonel Gaddafi, NATO planes had flown about 9,500 strike sorties against pro-Gaddafi targets.[77][78] A report from the organization Human Rights Watch in May 2012 identified at least 72 civilians killed in the campaign.[79]
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Following a coup d'état attempt in October 2013, Libyan Prime Minister Ali Zeidan requested technical advice and trainers from NATO to assist with ongoing security issues.[80]
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NATO has thirty members, mainly in Europe and North America. Some of these countries also have territory on multiple continents, which can be covered only as far south as the Tropic of Cancer in the Atlantic Ocean, which defines NATO's "area of responsibility" under Article 6 of the North Atlantic Treaty. During the original treaty negotiations, the United States insisted that colonies such as the Belgian Congo be excluded from the treaty.[81][82] French Algeria was however covered until their independence on 3 July 1962.[83] Twelve of these thirty are original members who joined in 1949, while the other eighteen joined in one of eight enlargement rounds.
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From the mid-1960s to the mid-1990s, France pursued a military strategy of independence from NATO under a policy dubbed "Gaullo-Mitterrandism".[84] Nicolas Sarkozy negotiated the return of France to the integrated military command and the Defence Planning Committee in 2009, the latter being disbanded the following year. France remains the only NATO member outside the Nuclear Planning Group and unlike the United States and the United Kingdom, will not commit its nuclear-armed submarines to the alliance.[15][23] Few members spend more than two percent of their gross domestic product on defence,[85] with the United States accounting for three quarters of NATO defence spending.[86]
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New membership in the alliance has been largely from Central and Eastern Europe, including former members of the Warsaw Pact. Accession to the alliance is governed with individual Membership Action Plans, and requires approval by each current member. NATO currently has one candidate country that is in the process of joining the alliance: Bosnia and Herzegovina. North Macedonia signed an accession protocol to become a NATO member state in February 2019, and became a member state on 27 March 2020.[87][88] Its accession had been blocked by Greece for many years due to the Macedonia naming dispute, which was resolved in 2018 by the Prespa agreement.[89] In order to support each other in the process, new and potential members in the region formed the Adriatic Charter in 2003.[90] Georgia was also named as an aspiring member, and was promised "future membership" during the 2008 summit in Bucharest,[91] though in 2014, US President Barack Obama said the country was not "currently on a path" to membership.[92]
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Russia continues to politically oppose further expansion, seeing it as inconsistent with informal understandings between Soviet leader Mikhail Gorbachev and European and US negotiators that allowed for a peaceful German reunification.[93] NATO's expansion efforts are often seen by Moscow leaders as a continuation of a Cold War attempt to surround and isolate Russia,[94] though they have also been criticized in the West.[95] A June 2016 Levada poll found that 68% of Russians think that deploying NATO troops in the Baltic states and Poland—former Eastern bloc countries bordering Russia—is a threat to Russia.[96] In contrast 65% of Poles surveyed in a 2017 Pew Research Center report identified Russia as a "major threat", with an average of 31% saying so across all NATO countries,[97] and 67% of Poles surveyed in 2018 favour US forces being based in Poland.[98] Of non-CIS Eastern European countries surveyed by Gallup in 2016, all but Serbia and Montenegro were more likely than not to view NATO as a protective alliance rather than a threat.[99] A 2006 study in the journal Security Studies argued that NATO enlargement contributed to democratic consolidation in Central and Eastern Europe.[100]
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Ukraine's relationship with NATO and Europe has been politically controversial, and improvement of these relations was one of the goals of the "Euromaidan" protests that saw the ousting of pro-Russian President Viktor Yanukovych in 2014. In March 2014, Prime Minister Arseniy Yatsenyuk reiterated the government's stance that Ukraine is not seeking NATO membership.[101] Ukraine's president subsequently signed a bill dropping his nation's nonaligned status in order to pursue NATO membership, but signalled that it would hold a referendum before seeking to join.[102] Ukraine is one of eight countries in Eastern Europe with an Individual Partnership Action Plan. IPAPs began in 2002, and are open to countries that have the political will and ability to deepen their relationship with NATO.[103]
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The Partnership for Peace (PfP) programme was established in 1994 and is based on individual bilateral relations between each partner country and NATO: each country may choose the extent of its participation.[105] Members include all current and former members of the Commonwealth of Independent States.[106] The Euro-Atlantic Partnership Council (EAPC) was first established on 29 May 1997, and is a forum for regular coordination, consultation and dialogue between all fifty participants.[107] The PfP programme is considered the operational wing of the Euro-Atlantic Partnership.[105] Other third countries also have been contacted for participation in some activities of the PfP framework such as Afghanistan.[108]
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The European Union (EU) signed a comprehensive package of arrangements with NATO under the Berlin Plus agreement on 16 December 2002. With this agreement, the EU was given the possibility of using NATO assets in case it wanted to act independently in an international crisis, on the condition that NATO itself did not want to act – the so-called "right of first refusal".[109] For example, Article 42(7) of the 1982 Treaty of Lisbon specifies that "If a Member State is the victim of armed aggression on its territory, the other Member States shall have towards it an obligation of aid and assistance by all the means in their power". The treaty applies globally to specified territories whereas NATO is restricted under its Article 6 to operations north of the Tropic of Cancer. It provides a "double framework" for the EU countries that are also linked with the PfP programme.[citation needed]
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Additionally, NATO cooperates and discusses its activities with numerous other non-NATO members. The Mediterranean Dialogue was established in 1994 to coordinate in a similar way with Israel and countries in North Africa. The Istanbul Cooperation Initiative was announced in 2004 as a dialogue forum for the Middle East along the same lines as the Mediterranean Dialogue. The four participants are also linked through the Gulf Cooperation Council.[110] In June 2018, Qatar expressed its wish to join NATO.[111] However, NATO declined membership, stating that only additional European countries could join according to Article 10 of NATO's founding treaty.[112] Qatar and NATO have previously signed a security agreement together in January 2018.[113]
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Political dialogue with Japan began in 1990, and since then, the Alliance has gradually increased its contact with countries that do not form part of any of these cooperation initiatives.[114] In 1998, NATO established a set of general guidelines that do not allow for a formal institutionalization of relations, but reflect the Allies' desire to increase cooperation. Following extensive debate, the term "Contact Countries" was agreed by the Allies in 2000. By 2012, the Alliance had broadened this group, which meets to discuss issues such as counter-piracy and technology exchange, under the names "partners across the globe" or "global partners".[115][116] Australia and New Zealand, both contact countries, are also members of the AUSCANNZUKUS strategic alliance, and similar regional or bilateral agreements between contact countries and NATO members also aid cooperation. NATO Secretary General Jens Stoltenberg stated that NATO needs to "address the rise of China," by closely cooperating with Australia, New Zealand, Japan and South Korea.[117] Colombia is the NATO's latest partner and Colombia has access to the full range of cooperative activities NATO offers to partners; Colombia became the first and only Latin American country to cooperate with NATO.[118]
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All agencies and organizations of NATO are integrated into either the civilian administrative or military executive roles. For the most part they perform roles and functions that directly or indirectly support the security role of the alliance as a whole.
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The civilian structure includes:
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The military structure includes:
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The organizations and agencies of NATO include:
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The NATO Parliamentary Assembly (NATO PA) is a body that sets broad strategic goals for NATO, which meets at two session per year. NATO PA interacts directly with the parliamentary structures of the national governments of the member states which appoint Permanent Members, or ambassadors to NATO. The NATO Parliamentary Assembly is made up of legislators from the member countries of the North Atlantic Alliance as well as thirteen associate members. It is however officially a different structure from NATO, and has as aim to join together deputies of NATO countries in order to discuss security policies on the NATO Council.
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The North Atlantic Treaty Organization (NATO, /ˈneɪtoʊ/; French: Organisation du traité de l'Atlantique nord, OTAN), also called the North Atlantic Alliance, is an intergovernmental military alliance between 30 North American and European countries. The organization implements the North Atlantic Treaty that was signed on 4 April 1949.[3][4] NATO constitutes a system of collective defence whereby its independent member states agree to mutual defence in response to an attack by any external party. NATO's Headquarters are located in Evere, Brussels, Belgium, while the headquarters of Allied Command Operations is near Mons, Belgium.
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Since its founding, the admission of new member states has increased the alliance from the original 12 countries to 30. The most recent member state to be added to NATO was North Macedonia on 27 March 2020. NATO currently recognizes Bosnia and Herzegovina, Georgia, and Ukraine as aspiring members.[5] An additional 20 countries participate in NATO's Partnership for Peace program, with 15 other countries involved in institutionalized dialogue programs. The combined military spending of all NATO members constitutes over 70% of the global total.[6] Members agreed that their aim is to reach or maintain the target defense spending of at least 2% of GDP by 2024.[7][8]
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On 4 March 1947, the Treaty of Dunkirk was signed by France and the United Kingdom as a Treaty of Alliance and Mutual Assistance in the event of a possible attack by Germany or the Soviet Union in the aftermath of World War II. In 1948, this alliance was expanded to include the Benelux countries, in the form of the Western Union, also referred to as the Brussels Treaty Organization (BTO), established by the Treaty of Brussels.[9] Talks for a new military alliance which could also include North America resulted in the signature of the North Atlantic Treaty on 4 April 1949 by the member states of the Western Union plus the United States, Canada, Portugal, Italy, Norway, Denmark and Iceland.[10]
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The North Atlantic Treaty was largely dormant until the Korean War initiated the establishment of NATO to implement it, by means of an integrated military structure: This included the formation of Supreme Headquarters Allied Powers Europe (SHAPE) in 1951, which adopted the Western Union's military structures and plans.[11] In 1952, the post of Secretary General of NATO was established as the organization's chief civilian. That year also saw the first major NATO maritime exercises, Exercise Mainbrace and the accession of Greece and Turkey to the organization.[12][13] Following the London and Paris Conferences, West Germany was permitted to rearm militarily, as they joined NATO in May 1955, which was, in turn, a major factor in the creation of the Soviet-dominated Warsaw Pact, delineating the two opposing sides of the Cold War.
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Doubts over the strength of the relationship between the European states and the United States ebbed and flowed, along with doubts over the credibility of the NATO defense against a prospective Soviet invasion – doubts that led to the development of the independent French nuclear deterrent and the withdrawal of France from NATO's military structure in 1966.[14][15] In 1982, the newly democratic Spain joined the alliance.
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The collapse of the Warsaw Pact in 1989–1991 removed the de facto main adversary of NATO and caused a strategic re-evaluation of NATO's purpose, nature, tasks, and focus on the continent of Europe. This shift started with the 1990 signing in Paris of the Treaty on Conventional Armed Forces in Europe between NATO and the Soviet Union, which mandated specific military reductions across the continent that continued after the dissolution of the Soviet Union in December 1991.[16] At that time, European countries accounted for 34 percent of NATO's military spending; by 2012, this had fallen to 21 percent.[17] NATO also began a gradual expansion to include newly autonomous Central and Eastern European nations, and extended its activities into political and humanitarian situations that had not formerly been NATO concerns.[18]
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After the fall of the Berlin Wall in Germany in 1989, the organization conducted its first military interventions in Bosnia from 1992 to 1995 and later Yugoslavia in 1999 during the break-up of Yugoslavia.[19] Politically, the organization sought better relations with former Warsaw Pact countries, most of which joined the alliance in 1999 and 2004. Article 5 of the North Atlantic treaty, requiring member states to come to the aid of any member state subject to an armed attack, was invoked for the first and only time after the September 11 attacks,[20] after which troops were deployed to Afghanistan under the NATO-led ISAF. The organization has operated a range of additional roles since then, including sending trainers to Iraq, assisting in counter-piracy operations[21] and in 2011 enforcing a no-fly zone over Libya in accordance with UN Security Council Resolution 1973. The less potent Article 4, which merely invokes consultation among NATO members, has been invoked five times following incidents in the Iraq War, Syrian Civil War, and annexation of Crimea.[22]
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The first post-Cold War expansion of NATO came with German reunification on 3 October 1990, when the former East Germany became part of the Federal Republic of Germany and the alliance. As part of post-Cold War restructuring, NATO's military structure was cut back and reorganized, with new forces such as the Headquarters Allied Command Europe Rapid Reaction Corps established. The changes brought about by the collapse of the Soviet Union on the military balance in Europe were recognized in the Adapted Conventional Armed Forces in Europe Treaty, which was signed in 1999. The policies of French President Nicolas Sarkozy resulted in a major reform of France's military position, culminating with the return to full membership on 4 April 2009, which also included France rejoining the NATO Military Command Structure, while maintaining an independent nuclear deterrent.[15][23][24]
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Between 1994 and 1997, wider forums for regional cooperation between NATO and its neighbors were set up, like the Partnership for Peace, the Mediterranean Dialogue initiative and the Euro-Atlantic Partnership Council. In 1998, the NATO–Russia Permanent Joint Council was established. Between 1999 and 2020 NATO incorporated the following Central and Eastern European countries, including several former communist states: the Czech Republic, Hungary, Poland, Bulgaria, Estonia, Latvia, Lithuania, Romania, Slovakia, Slovenia, Albania, Croatia, Montenegro, and North Macedonia.[18]
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The Russian intervention in Crimea in 2014 led to strong condemnation by NATO nations and the creation of a new "spearhead" force of 5,000 troops at bases in Estonia, Lithuania, Latvia, Poland, Romania, and Bulgaria.[25] At the subsequent 2014 Wales summit, the leaders of NATO's member states formally committed for the first time to spend the equivalent of at least 2% of their gross domestic products on defence by 2024, which had previously been only an informal guideline.[26] NATO did not condemn the 2016–present purges in Turkey.[27] NATO members have resisted the UN's Nuclear Weapon Ban Treaty, a binding agreement for negotiations for the total elimination of nuclear weapons, supported by more than 120 nations.[28]
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No military operations were conducted by NATO during the Cold War. Following the end of the Cold War, the first operations, Anchor Guard in 1990 and Ace Guard in 1991, were prompted by the Iraqi invasion of Kuwait. Airborne early warning aircraft were sent to provide coverage of southeastern Turkey, and later a quick-reaction force was deployed to the area.[29]
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The Bosnian War began in 1992, as a result of the break-up of Yugoslavia. The deteriorating situation led to United Nations Security Council Resolution 816 on 9 October 1992, ordering a no-fly zone over central Bosnia and Herzegovina, which NATO began enforcing on 12 April 1993 with Operation Deny Flight. From June 1993 until October 1996, Operation Sharp Guard added maritime enforcement of the arms embargo and economic sanctions against the Federal Republic of Yugoslavia. On 28 February 1994, NATO took its first wartime action by shooting down four Bosnian Serb aircraft violating the no-fly zone.[30]
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On 10 and 11 April 1994, the United Nations Protection Force called in air strikes to protect the Goražde safe area, resulting in the bombing of a Bosnian Serb military command outpost near Goražde by two US F-16 jets acting under NATO direction.[31] In retaliation, Serbs took 150 U.N. personnel hostage on 14 April.[32][33] On 16 April a British Sea Harrier was shot down over Goražde by Serb forces.[34]
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In August 1995, a two-week NATO bombing campaign, Operation Deliberate Force, began against the Army of the Republika Srpska, after the Srebrenica massacre.[35] Further NATO air strikes helped bring the Yugoslav wars to an end, resulting in the Dayton Agreement in November 1995.[35] As part of this agreement, NATO deployed a UN-mandated peacekeeping force, under Operation Joint Endeavor, named IFOR. Almost 60,000 NATO troops were joined by forces from non-NATO nations in this peacekeeping mission. This transitioned into the smaller SFOR, which started with 32,000 troops initially and ran from December 1996 until December 2004, when operations were then passed onto European Union Force Althea.[36] Following the lead of its member nations, NATO began to award a service medal, the NATO Medal, for these operations.[37]
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In an effort to stop Slobodan Milošević's Serbian-led crackdown on KLA separatists and Albanian civilians in Kosovo, the United Nations Security Council passed Resolution 1199 on 23 September 1998 to demand a ceasefire. Negotiations under US Special Envoy Richard Holbrooke broke down on 23 March 1999, and he handed the matter to NATO,[38] which started a 78-day bombing campaign on 24 March 1999.[39] Operation Allied Force targeted the military capabilities of what was then the Federal Republic of Yugoslavia. During the crisis, NATO also deployed one of its international reaction forces, the ACE Mobile Force (Land), to Albania as the Albania Force (AFOR), to deliver humanitarian aid to refugees from Kosovo.[40]
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Though the campaign was criticized for high civilian casualties, including bombing of the Chinese embassy in Belgrade, Milošević finally accepted the terms of an international peace plan on 3 June 1999, ending the Kosovo War. On 11 June, Milošević further accepted UN resolution 1244, under the mandate of which NATO then helped establish the KFOR peacekeeping force. Nearly one million refugees had fled Kosovo, and part of KFOR's mandate was to protect the humanitarian missions, in addition to deterring violence.[40][41] In August–September 2001, the alliance also mounted Operation Essential Harvest, a mission disarming ethnic Albanian militias in the Republic of Macedonia.[42] As of 1 December 2013[update], 4,882 KFOR soldiers, representing 31 countries, continue to operate in the area.[43]
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The US, the UK, and most other NATO countries opposed efforts to require the UN Security Council to approve NATO military strikes, such as the action against Serbia in 1999, while France and some others claimed that the alliance needed UN approval.[44] The US/UK side claimed that this would undermine the authority of the alliance, and they noted that Russia and China would have exercised their Security Council vetoes to block the strike on Yugoslavia, and could do the same in future conflicts where NATO intervention was required, thus nullifying the entire potency and purpose of the organization. Recognizing the post-Cold War military environment, NATO adopted the Alliance Strategic Concept during its Washington summit in April 1999 that emphasized conflict prevention and crisis management.[45]
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The September 11 attacks in the United States caused NATO to invoke Article 5 of the NATO Charter for the first time in the organization's history. The Article states that an attack on any member shall be considered to be an attack on all. The invocation was confirmed on 4 October 2001 when NATO determined that the attacks were indeed eligible under the terms of the North Atlantic Treaty.[46] The eight official actions taken by NATO in response to the attacks included Operation Eagle Assist and Operation Active Endeavour, a naval operation in the Mediterranean Sea designed to prevent the movement of terrorists or weapons of mass destruction, and to enhance the security of shipping in general, which began on 4 October 2001.[47]
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The alliance showed unity: On 16 April 2003, NATO agreed to take command of the International Security Assistance Force (ISAF), which included troops from 42 countries. The decision came at the request of Germany and the Netherlands, the two nations leading ISAF at the time of the agreement, and all nineteen NATO ambassadors approved it unanimously. The handover of control to NATO took place on 11 August, and marked the first time in NATO's history that it took charge of a mission outside the north Atlantic area.[48]
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ISAF was initially charged with securing Kabul and surrounding areas from the Taliban, al Qaeda and factional warlords, so as to allow for the establishment of the Afghan Transitional Administration headed by Hamid Karzai. In October 2003, the UN Security Council authorized the expansion of the ISAF mission throughout Afghanistan,[49] and ISAF subsequently expanded the mission in four main stages over the whole of the country.[50]
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On 31 July 2006, the ISAF additionally took over military operations in the south of Afghanistan from a US-led anti-terrorism coalition.[51] Due to the intensity of the fighting in the south, in 2011 France allowed a squadron of Mirage 2000 fighter/attack aircraft to be moved into the area, to Kandahar, in order to reinforce the alliance's efforts.[52] During its 2012 Chicago Summit, NATO endorsed a plan to end the Afghanistan war and to remove the NATO-led ISAF Forces by the end of December 2014.[53] ISAF was disestablished in December 2014 and replaced by the follow-on training Resolute Support Mission.[54]
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In August 2004, during the Iraq War, NATO formed the NATO Training Mission – Iraq, a training mission to assist the Iraqi security forces in conjunction with the US led MNF-I.[55] The NATO Training Mission-Iraq (NTM-I) was established at the request of the Iraqi Interim Government under the provisions of United Nations Security Council Resolution 1546. The aim of NTM-I was to assist in the development of Iraqi security forces training structures and institutions so that Iraq can build an effective and sustainable capability that addresses the needs of the nation. NTM-I was not a combat mission but is a distinct mission, under the political control of NATO's North Atlantic Council. Its operational emphasis was on training and mentoring. The activities of the mission were coordinated with Iraqi authorities and the US-led Deputy Commanding General Advising and Training, who was also dual-hatted as the Commander of NTM-I. The mission officially concluded on 17 December 2011.[56]
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Turkey invoked the first Article 4 meetings in 2003 at the start of the Iraq War. Turkey also invoked this article twice in 2012 during the Syrian Civil War, after the downing of an unarmed Turkish F-4 reconnaissance jet, and after a mortar was fired at Turkey from Syria,[57] and again in 2015 after threats by Islamic State of Iraq and the Levant to its territorial integrity.[58]
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Beginning on 17 August 2009, NATO deployed warships in an operation to protect maritime traffic in the Gulf of Aden and the Indian Ocean from Somali pirates, and help strengthen the navies and coast guards of regional states. The operation was approved by the North Atlantic Council and involves warships primarily from the United States though vessels from many other nations are also included. Operation Ocean Shield focuses on protecting the ships of Operation Allied Provider which are distributing aid as part of the World Food Programme mission in Somalia. Russia, China and South Korea have sent warships to participate in the activities as well.[59][60] The operation seeks to dissuade and interrupt pirate attacks, protect vessels, and abetting to increase the general level of security in the region.[61]
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During the Libyan Civil War, violence between protesters and the Libyan government under Colonel Muammar Gaddafi escalated, and on 17 March 2011 led to the passage of United Nations Security Council Resolution 1973, which called for a ceasefire, and authorized military action to protect civilians. A coalition that included several NATO members began enforcing a no-fly zone over Libya shortly afterwards, beginning with Opération Harmattan by the French Air Force on 19 March.
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On 20 March 2011, NATO states agreed on enforcing an arms embargo against Libya with Operation Unified Protector using ships from NATO Standing Maritime Group 1 and Standing Mine Countermeasures Group 1,[62] and additional ships and submarines from NATO members.[63] They would "monitor, report and, if needed, interdict vessels suspected of carrying illegal arms or mercenaries".[62]
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On 24 March, NATO agreed to take control of the no-fly zone from the initial coalition, while command of targeting ground units remained with the coalition's forces.[64][65] NATO began officially enforcing the UN resolution on 27 March 2011 with assistance from Qatar and the United Arab Emirates.[66] By June, reports of divisions within the alliance surfaced as only eight of the 28 member nations were participating in combat operations,[67] resulting in a confrontation between US Defense Secretary Robert Gates and countries such as Poland, Spain, the Netherlands, Turkey, and Germany to contribute more, the latter believing the organization has overstepped its mandate in the conflict.[68][69][70] In his final policy speech in Brussels on 10 June, Gates further criticized allied countries in suggesting their actions could cause the demise of NATO.[71] The German foreign ministry pointed to "a considerable [German] contribution to NATO and NATO-led operations" and to the fact that this engagement was highly valued by President Obama.[72]
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While the mission was extended into September, Norway that day announced it would begin scaling down contributions and complete withdrawal by 1 August.[73] Earlier that week it was reported Danish air fighters were running out of bombs.[74][75] The following week, the head of the Royal Navy said the country's operations in the conflict were not sustainable.[76] By the end of the mission in October 2011, after the death of Colonel Gaddafi, NATO planes had flown about 9,500 strike sorties against pro-Gaddafi targets.[77][78] A report from the organization Human Rights Watch in May 2012 identified at least 72 civilians killed in the campaign.[79]
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Following a coup d'état attempt in October 2013, Libyan Prime Minister Ali Zeidan requested technical advice and trainers from NATO to assist with ongoing security issues.[80]
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NATO has thirty members, mainly in Europe and North America. Some of these countries also have territory on multiple continents, which can be covered only as far south as the Tropic of Cancer in the Atlantic Ocean, which defines NATO's "area of responsibility" under Article 6 of the North Atlantic Treaty. During the original treaty negotiations, the United States insisted that colonies such as the Belgian Congo be excluded from the treaty.[81][82] French Algeria was however covered until their independence on 3 July 1962.[83] Twelve of these thirty are original members who joined in 1949, while the other eighteen joined in one of eight enlargement rounds.
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From the mid-1960s to the mid-1990s, France pursued a military strategy of independence from NATO under a policy dubbed "Gaullo-Mitterrandism".[84] Nicolas Sarkozy negotiated the return of France to the integrated military command and the Defence Planning Committee in 2009, the latter being disbanded the following year. France remains the only NATO member outside the Nuclear Planning Group and unlike the United States and the United Kingdom, will not commit its nuclear-armed submarines to the alliance.[15][23] Few members spend more than two percent of their gross domestic product on defence,[85] with the United States accounting for three quarters of NATO defence spending.[86]
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New membership in the alliance has been largely from Central and Eastern Europe, including former members of the Warsaw Pact. Accession to the alliance is governed with individual Membership Action Plans, and requires approval by each current member. NATO currently has one candidate country that is in the process of joining the alliance: Bosnia and Herzegovina. North Macedonia signed an accession protocol to become a NATO member state in February 2019, and became a member state on 27 March 2020.[87][88] Its accession had been blocked by Greece for many years due to the Macedonia naming dispute, which was resolved in 2018 by the Prespa agreement.[89] In order to support each other in the process, new and potential members in the region formed the Adriatic Charter in 2003.[90] Georgia was also named as an aspiring member, and was promised "future membership" during the 2008 summit in Bucharest,[91] though in 2014, US President Barack Obama said the country was not "currently on a path" to membership.[92]
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Russia continues to politically oppose further expansion, seeing it as inconsistent with informal understandings between Soviet leader Mikhail Gorbachev and European and US negotiators that allowed for a peaceful German reunification.[93] NATO's expansion efforts are often seen by Moscow leaders as a continuation of a Cold War attempt to surround and isolate Russia,[94] though they have also been criticized in the West.[95] A June 2016 Levada poll found that 68% of Russians think that deploying NATO troops in the Baltic states and Poland—former Eastern bloc countries bordering Russia—is a threat to Russia.[96] In contrast 65% of Poles surveyed in a 2017 Pew Research Center report identified Russia as a "major threat", with an average of 31% saying so across all NATO countries,[97] and 67% of Poles surveyed in 2018 favour US forces being based in Poland.[98] Of non-CIS Eastern European countries surveyed by Gallup in 2016, all but Serbia and Montenegro were more likely than not to view NATO as a protective alliance rather than a threat.[99] A 2006 study in the journal Security Studies argued that NATO enlargement contributed to democratic consolidation in Central and Eastern Europe.[100]
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Ukraine's relationship with NATO and Europe has been politically controversial, and improvement of these relations was one of the goals of the "Euromaidan" protests that saw the ousting of pro-Russian President Viktor Yanukovych in 2014. In March 2014, Prime Minister Arseniy Yatsenyuk reiterated the government's stance that Ukraine is not seeking NATO membership.[101] Ukraine's president subsequently signed a bill dropping his nation's nonaligned status in order to pursue NATO membership, but signalled that it would hold a referendum before seeking to join.[102] Ukraine is one of eight countries in Eastern Europe with an Individual Partnership Action Plan. IPAPs began in 2002, and are open to countries that have the political will and ability to deepen their relationship with NATO.[103]
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The Partnership for Peace (PfP) programme was established in 1994 and is based on individual bilateral relations between each partner country and NATO: each country may choose the extent of its participation.[105] Members include all current and former members of the Commonwealth of Independent States.[106] The Euro-Atlantic Partnership Council (EAPC) was first established on 29 May 1997, and is a forum for regular coordination, consultation and dialogue between all fifty participants.[107] The PfP programme is considered the operational wing of the Euro-Atlantic Partnership.[105] Other third countries also have been contacted for participation in some activities of the PfP framework such as Afghanistan.[108]
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The European Union (EU) signed a comprehensive package of arrangements with NATO under the Berlin Plus agreement on 16 December 2002. With this agreement, the EU was given the possibility of using NATO assets in case it wanted to act independently in an international crisis, on the condition that NATO itself did not want to act – the so-called "right of first refusal".[109] For example, Article 42(7) of the 1982 Treaty of Lisbon specifies that "If a Member State is the victim of armed aggression on its territory, the other Member States shall have towards it an obligation of aid and assistance by all the means in their power". The treaty applies globally to specified territories whereas NATO is restricted under its Article 6 to operations north of the Tropic of Cancer. It provides a "double framework" for the EU countries that are also linked with the PfP programme.[citation needed]
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Additionally, NATO cooperates and discusses its activities with numerous other non-NATO members. The Mediterranean Dialogue was established in 1994 to coordinate in a similar way with Israel and countries in North Africa. The Istanbul Cooperation Initiative was announced in 2004 as a dialogue forum for the Middle East along the same lines as the Mediterranean Dialogue. The four participants are also linked through the Gulf Cooperation Council.[110] In June 2018, Qatar expressed its wish to join NATO.[111] However, NATO declined membership, stating that only additional European countries could join according to Article 10 of NATO's founding treaty.[112] Qatar and NATO have previously signed a security agreement together in January 2018.[113]
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Political dialogue with Japan began in 1990, and since then, the Alliance has gradually increased its contact with countries that do not form part of any of these cooperation initiatives.[114] In 1998, NATO established a set of general guidelines that do not allow for a formal institutionalization of relations, but reflect the Allies' desire to increase cooperation. Following extensive debate, the term "Contact Countries" was agreed by the Allies in 2000. By 2012, the Alliance had broadened this group, which meets to discuss issues such as counter-piracy and technology exchange, under the names "partners across the globe" or "global partners".[115][116] Australia and New Zealand, both contact countries, are also members of the AUSCANNZUKUS strategic alliance, and similar regional or bilateral agreements between contact countries and NATO members also aid cooperation. NATO Secretary General Jens Stoltenberg stated that NATO needs to "address the rise of China," by closely cooperating with Australia, New Zealand, Japan and South Korea.[117] Colombia is the NATO's latest partner and Colombia has access to the full range of cooperative activities NATO offers to partners; Colombia became the first and only Latin American country to cooperate with NATO.[118]
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All agencies and organizations of NATO are integrated into either the civilian administrative or military executive roles. For the most part they perform roles and functions that directly or indirectly support the security role of the alliance as a whole.
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The civilian structure includes:
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The military structure includes:
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The organizations and agencies of NATO include:
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+
The NATO Parliamentary Assembly (NATO PA) is a body that sets broad strategic goals for NATO, which meets at two session per year. NATO PA interacts directly with the parliamentary structures of the national governments of the member states which appoint Permanent Members, or ambassadors to NATO. The NATO Parliamentary Assembly is made up of legislators from the member countries of the North Atlantic Alliance as well as thirteen associate members. It is however officially a different structure from NATO, and has as aim to join together deputies of NATO countries in order to discuss security policies on the NATO Council.
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en/4309.html.txt
ADDED
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1 |
+
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In biology, an organism (from Greek: ὀργανισμός, organismos) is any individual entity that embodies the properties of life. It is a synonym for "life form".
|
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+
|
5 |
+
Organisms are classified by taxonomy into groups such as multicellular animals, plants, and fungi; or unicellular microorganisms such as protists, bacteria, and archaea.[1] All types of organisms are capable of reproduction, growth and development, maintenance, and some degree of response to stimuli. Humans, squids, mushrooms, and vascular plants are examples of multicellular organisms that differentiate specialized tissues and organs during development.
|
6 |
+
|
7 |
+
An organism may be either a prokaryote or a eukaryote. Prokaryotes are represented by two separate domains – bacteria and archaea. Eukaryotic organisms are characterized by the presence of a membrane-bound cell nucleus and contain additional membrane-bound compartments called organelles (such as mitochondria in animals and plants and plastids in plants and algae, all generally considered to be derived from endosymbiotic bacteria).[2] Fungi, animals and plants are examples of kingdoms of organisms within the eukaryotes.
|
8 |
+
|
9 |
+
Estimates on the number of Earth's current species range from 2 million to 1 trillion,[3] of which over 1.7 million have been documented.[4] More than 99% of all species, amounting to over five billion species,[5] that ever lived are estimated to be extinct.[6][7]
|
10 |
+
|
11 |
+
In 2016, a set of 355 genes from the last universal common ancestor (LUCA) of all organisms was identified.[8][9]
|
12 |
+
|
13 |
+
The term "organism" (from Greek ὀργανισμός, organismos, from ὄργανον, organon, i.e. "instrument, implement, tool, organ of sense or apprehension")[10][11] first appeared in the English language in 1703 and took on its current definition by 1834 (Oxford English Dictionary). It is directly related to the term "organization". There is a long tradition of defining organisms as self-organizing beings, going back at least to Immanuel Kant's 1790 Critique of Judgment.[12]
|
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+
|
15 |
+
An organism may be defined as an assembly of molecules functioning as a more or less stable whole that exhibits the properties of life. Dictionary definitions can be broad, using phrases such as "any living structure, such as a plant, animal, fungus or bacterium, capable of growth and reproduction".[13] Many definitions exclude viruses and possible man-made non-organic life forms, as viruses are dependent on the biochemical machinery of a host cell for reproduction.[14] A superorganism is an organism consisting of many individuals working together as a single functional or social unit.[15]
|
16 |
+
|
17 |
+
There has been controversy about the best way to define the organism[16][17][18][19][20][21][22][23][24][25] and indeed about whether or not such a definition is necessary.[26][27] Several contributions[28] are responses to the suggestion that the category of "organism" may well not be adequate in biology.[29][page needed]
|
18 |
+
|
19 |
+
Viruses are not typically considered to be organisms because they are incapable of autonomous reproduction, growth or metabolism. Although some organisms are also incapable of independent survival and live as obligatory intracellular parasites, they are capable of independent metabolism and procreation. Although viruses have a few enzymes and molecules characteristic of living organisms, they have no metabolism of their own; they cannot synthesize and organize the organic compounds from which they are formed. Naturally, this rules out autonomous reproduction: they can only be passively replicated by the machinery of the host cell. In this sense, they are similar to inanimate matter.
|
20 |
+
|
21 |
+
While viruses sustain no independent metabolism and thus are usually not classified as organisms, they do have their own genes, and they do evolve by mechanisms similar to the evolutionary mechanisms of organisms. Thus, an argument that viruses should be classed as living organisms is their ability to undergo evolution and replicate through self-assembly. However, some scientists argue that viruses neither evolve nor self-reproduce. Instead, viruses are evolved by their host cells, meaning that there was co-evolution of viruses and host cells. If host cells did not exist, viral evolution would be impossible. This is not true for cells. If viruses did not exist, the direction of cellular evolution could be different, but cells would nevertheless be able to evolve. As for the reproduction, viruses totally rely on hosts' machinery to replicate.[30] The discovery of viruses with genes coding for energy metabolism and protein synthesis fuelled the debate about whether viruses are living organisms. The presence of these genes suggested that viruses were once able to metabolize. However, it was found later that the genes coding for energy and protein metabolism have a cellular origin. Most likely, these genes were acquired through horizontal gene transfer from viral hosts.[30]
|
22 |
+
|
23 |
+
Organisms are complex chemical systems, organized in ways that promote reproduction and some measure of sustainability or survival. The same laws that govern non-living chemistry govern the chemical processes of life. It is generally the phenomena of entire organisms that determine their fitness to an environment and therefore the survivability of their DNA-based genes.
|
24 |
+
|
25 |
+
Organisms clearly owe their origin, metabolism, and many other internal functions to chemical phenomena, especially the chemistry of large organic molecules. Organisms are complex systems of chemical compounds that, through interaction and environment, play a wide variety of roles.
|
26 |
+
|
27 |
+
Organisms are semi-closed chemical systems. Although they are individual units of life (as the definition requires), they are not closed to the environment around them. To operate they constantly take in and release energy. Autotrophs produce usable energy (in the form of organic compounds) using light from the sun or inorganic compounds while heterotrophs take in organic compounds from the environment.
|
28 |
+
|
29 |
+
The primary chemical element in these compounds is carbon. The chemical properties of this element such as its great affinity for bonding with other small atoms, including other carbon atoms, and its small size making it capable of forming multiple bonds, make it ideal as the basis of organic life. It is able to form small three-atom compounds (such as carbon dioxide), as well as large chains of many thousands of atoms that can store data (nucleic acids), hold cells together, and transmit information (protein).
|
30 |
+
|
31 |
+
Compounds that make up organisms may be divided into macromolecules and other, smaller molecules. The four groups of macromolecule are nucleic acids, proteins, carbohydrates and lipids. Nucleic acids (specifically deoxyribonucleic acid, or DNA) store genetic data as a sequence of nucleotides. The particular sequence of the four different types of nucleotides (adenine, cytosine, guanine, and thymine) dictate many characteristics that constitute the organism. The sequence is divided up into codons, each of which is a particular sequence of three nucleotides and corresponds to a particular amino acid. Thus a sequence of DNA codes for a particular protein that, due to the chemical properties of the amino acids it is made from, folds in a particular manner and so performs a particular function.
|
32 |
+
|
33 |
+
These protein functions have been recognized:
|
34 |
+
|
35 |
+
A bilayer of phospholipids makes up the membrane of cells that constitutes a barrier, containing everything within the cell and preventing compounds from freely passing into, and out of, the cell. Due to the selective permeability of the phospholipid membrane, only specific compounds can pass through it.
|
36 |
+
|
37 |
+
All organisms consist of structural units called cells; some contain a single cell (unicellular) and others contain many units (multicellular). Multicellular organisms are able to specialize cells to perform specific functions. A group of such cells is a tissue, and in animals these occur as four basic types, namely epithelium, nervous tissue, muscle tissue, and connective tissue. Several types of tissue work together in the form of an organ to produce a particular function (such as the pumping of the blood by the heart, or as a barrier to the environment as the skin). This pattern continues to a higher level with several organs functioning as an organ system such as the reproductive system, and digestive system. Many multicellular organisms consist of several organ systems, which coordinate to allow for life.
|
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+
|
39 |
+
The cell theory, first developed in 1839 by Schleiden and Schwann, states that all organisms are composed of one or more cells; all cells come from preexisting cells; and cells contain the hereditary information necessary for regulating cell functions and for transmitting information to the next generation of cells.
|
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+
|
41 |
+
There are two types of cells, eukaryotic and prokaryotic. Prokaryotic cells are usually singletons, while eukaryotic cells are usually found in multicellular organisms. Prokaryotic cells lack a nuclear membrane so DNA is unbound within the cell; eukaryotic cells have nuclear membranes.
|
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+
|
43 |
+
All cells, whether prokaryotic or eukaryotic, have a membrane, which envelops the cell, separates its interior from its environment, regulates what moves in and out, and maintains the electric potential of the cell. Inside the membrane, a salty cytoplasm takes up most of the cell volume. All cells possess DNA, the hereditary material of genes, and RNA, containing the information necessary to build various proteins such as enzymes, the cell's primary machinery. There are also other kinds of biomolecules in cells.
|
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+
|
45 |
+
All cells share several similar characteristics of:[31]
|
46 |
+
|
47 |
+
The last universal common ancestor (LUCA) is the most recent organism from which all organisms now living on Earth descend.[32] Thus it is the most recent common ancestor of all current life on Earth. The LUCA is estimated to have lived some 3.5 to 3.8 billion years ago (sometime in the Paleoarchean era).[33][34] The earliest evidence for life on Earth is graphite found to be biogenic in 3.7 billion-year-old metasedimentary rocks discovered in Western Greenland[35] and microbial mat fossils found in 3.48 billion-year-old sandstone discovered in Western Australia.[36][37] Although more than 99 percent of all species that ever lived on the planet are estimated to be extinct,[6][7] there are currently 2 million to 1 trillion species of life on Earth.[3]
|
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+
|
49 |
+
Information about the early development of life includes input from many different fields, including geology and planetary science. These sciences provide information about the history of the Earth and the changes produced by life. However, a great deal of information about the early Earth has been destroyed by geological processes over the course of time.
|
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+
|
51 |
+
All organisms are descended from a common ancestor or ancestral gene pool. Evidence for common descent may be found in traits shared between all living organisms. In Darwin's day, the evidence of shared traits was based solely on visible observation of morphologic similarities, such as the fact that all birds have wings, even those that do not fly.
|
52 |
+
|
53 |
+
There is strong evidence from genetics that all organisms have a common ancestor. For example, every living cell makes use of nucleic acids as its genetic material, and uses the same twenty amino acids as the building blocks for proteins. All organisms use the same genetic code (with some extremely rare and minor deviations) to translate nucleic acid sequences into proteins. The universality of these traits strongly suggests common ancestry, because the selection of many of these traits seems arbitrary. Horizontal gene transfer makes it more difficult to study the last universal ancestor.[38] However, the universal use of the same genetic code, same nucleotides, and same amino acids makes the existence of such an ancestor overwhelmingly likely.[39]
|
54 |
+
|
55 |
+
Chlorobacteria (accepted name = Chloroflexi)
|
56 |
+
|
57 |
+
Hadobacteria (= Deinococcus-Thermus group)
|
58 |
+
|
59 |
+
Cyanobacteria
|
60 |
+
|
61 |
+
Spirochaetae
|
62 |
+
|
63 |
+
Fibrobacteres
|
64 |
+
|
65 |
+
Chlorobi
|
66 |
+
|
67 |
+
Bacteroidetes
|
68 |
+
|
69 |
+
Planctomycetes
|
70 |
+
|
71 |
+
Chlamydiae
|
72 |
+
|
73 |
+
Lentisphaerae
|
74 |
+
|
75 |
+
Verrucomicrobia
|
76 |
+
|
77 |
+
Deferribacteres
|
78 |
+
|
79 |
+
Acidobacteria
|
80 |
+
|
81 |
+
Deltaproteobacteria
|
82 |
+
|
83 |
+
Epsilonproteobacteria
|
84 |
+
|
85 |
+
Alphaproteobacteria
|
86 |
+
|
87 |
+
Betaproteobacteria
|
88 |
+
|
89 |
+
Gammaproteobacteria
|
90 |
+
|
91 |
+
Thermotogae
|
92 |
+
|
93 |
+
Fusobacteria
|
94 |
+
|
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+
Negativicutes
|
96 |
+
|
97 |
+
Endobacteria (=Firmicutes, Mollicutes)
|
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|
99 |
+
Actinobacteria
|
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|
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+
Archaea
|
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|
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+
Eukarya
|
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|
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+
The most commonly accepted location of the root of the tree of life is between a monophyletic domain Bacteria and a clade formed by Archaea and Eukaryota of what is referred to as the "traditional tree of life" based on several molecular studies.[40][41][42][43][44][45] A very small minority of studies have concluded differently, namely that the root is in the domain Bacteria, either in the phylum Firmicutes[46] or that the phylum Chloroflexi is basal to a clade with Archaea and Eukaryotes and the rest of Bacteria as proposed by Thomas Cavalier-Smith.[47]
|
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+
|
107 |
+
Research published in 2016, by William F. Martin, by genetically analyzing 6.1 million protein-coding genes from sequenced prokaryotic genomes of various phylogenetic trees, identified 355 protein clusters from amongst 286,514 protein clusters that were probably common to the LUCA. The results "depict LUCA as anaerobic, CO2-fixing, H2-dependent with a Wood–Ljungdahl pathway (the reductive acetyl-coenzyme A pathway), N2-fixing and thermophilic. LUCA's biochemistry was replete with FeS clusters and radical reaction mechanisms. Its cofactors reveal dependence upon transition metals, flavins, S-adenosyl methionine, coenzyme A, ferredoxin, molybdopterin, corrins and selenium. Its genetic code required nucleoside modifications and S-adenosylmethionine-dependent methylations." The results depict methanogenic clostria as a basal clade in the 355 lineages examined, and suggest that the LUCA inhabited an anaerobic hydrothermal vent setting in a geochemically active environment rich in H2, CO2, and iron.[8] However, the identification of these genes as being present in LUCA was criticized, suggesting that many of the proteins assumed to be present in LUCA represent later horizontal gene transfers between archaea and bacteria.[48]
|
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+
|
109 |
+
Sexual reproduction is widespread among current eukaryotes, and was likely present in the last common ancestor.[49] This is suggested by the finding of a core set of genes for meiosis in the descendants of lineages that diverged early from the eukaryotic evolutionary tree.[50] and Malik et al.[51] It is further supported by evidence that eukaryotes previously regarded as "ancient asexuals", such as Amoeba, were likely sexual in the past, and that most present day asexual amoeboid lineages likely arose recently and independently.[52]
|
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|
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+
In prokaryotes, natural bacterial transformation involves the transfer of DNA from one bacterium to another and integration of the donor DNA into the recipient chromosome by recombination. Natural bacterial transformation is considered to be a primitive sexual process and occurs in both bacteria and archaea, although it has been studied mainly in bacteria. Transformation is clearly a bacterial adaptation and not an accidental occurrence, because it depends on numerous gene products that specifically interact with each other to enter a state of natural competence to perform this complex process.[53] Transformation is a common mode of DNA transfer among prokaryotes.[54]
|
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|
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+
The ancestry of living organisms has traditionally been reconstructed from morphology, but is increasingly supplemented with phylogenetics – the reconstruction of phylogenies by the comparison of genetic (DNA) sequence.
|
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+
|
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+
Sequence comparisons suggest recent horizontal transfer of many genes among diverse species including across the boundaries of phylogenetic "domains". Thus determining the phylogenetic history of a species can not be done conclusively by determining evolutionary trees for single genes.[55]
|
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|
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+
Biologist Peter Gogarten suggests "the original metaphor of a tree no longer fits the data from recent genome research", therefore "biologists (should) use the metaphor of a mosaic to describe the different histories combined in individual genomes and use (the) metaphor of a net to visualize the rich exchange and cooperative effects of HGT among microbes."[56]
|
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+
|
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+
Modern biotechnology is challenging traditional concepts of organism and species. Cloning is the process of creating a new multicellular organism, genetically identical to another, with the potential of creating entirely new species of organisms. Cloning is the subject of much ethical debate.
|
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|
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+
In 2008, the J. Craig Venter Institute assembled a synthetic bacterial genome, Mycoplasma genitalium, by using recombination in yeast of 25 overlapping DNA fragments in a single step. The use of yeast recombination greatly simplifies the assembly of large DNA molecules from both synthetic and natural fragments.[57] Other companies, such as Synthetic Genomics, have already been formed to take advantage of the many commercial uses of custom designed genomes.
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Athens (/ˈæθɪnz/ ATH-inz;[3] Greek: Αθήνα, romanized: Athína [aˈθina] (listen); Ancient Greek: Ἀθῆναι, romanized: Athênai (pl.) [atʰɛ̂ːnai̯]), is the capital and largest city of Greece. Athens dominates the Attica region and is one of the world's oldest cities, with its recorded history spanning over 3,400 years[citation needed] and its earliest human presence started somewhere between the 11th and 7th millennium BC.[4]
|
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|
5 |
+
Classical Athens was a powerful city-state that emerged in conjunction with the seagoing development of the port of Piraeus. A center for the arts, learning and philosophy, home of Plato's Academy and Aristotle's Lyceum,[5][6] it is widely referred to as the cradle of Western civilization and the birthplace of democracy,[7][8] largely because of its cultural and political impact on the European continent, and in particular the Romans.[9] In modern times, Athens is a large cosmopolitan metropolis and central to economic, financial, industrial, maritime, political and cultural life in Greece.
|
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|
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Athens is a global city and one of the biggest economic centers in southeastern Europe. It has a large financial sector, and its port Piraeus is both the largest passenger port in Europe,[10][11] and the second largest in the world.[12]
|
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|
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The Municipality of Athens (also City of Athens), which actually constitutes a small administrative unit of the entire city, had a population of 664,046 (in 2011)[2] within its official limits, and a land area of 38.96 km2 (15.04 sq mi).[13][14] The Athens Urban Area (Greater Athens and Greater Piraeus) extends beyond its administrative municipal city limits, with a population of 3,090,508 (in 2011)[15] over an area of 412 km2 (159 sq mi).[14] According to Eurostat[16] in 2011, the functional urban area (FUA) of Athens was the 9th most populous FUA in the European Union (the 6th most populous capital city of the EU), with a population of 3.8 million people. Athens is also the southernmost capital on the European mainland and the warmest major city in Europe.
|
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|
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The heritage of the classical era is still evident in the city, represented by ancient monuments and works of art, the most famous of all being the Parthenon, considered a key landmark of early Western civilization. The city also retains Roman and Byzantine monuments, as well as a smaller number of Ottoman monuments. Athens is home to two UNESCO World Heritage Sites, the Acropolis of Athens and the medieval Daphni Monastery. Landmarks of the modern era, dating back to the establishment of Athens as the capital of the independent Greek state in 1834, include the Hellenic Parliament and the so-called "architectural trilogy of Athens", consisting of the National Library of Greece, the National and Kapodistrian University of Athens and the Academy of Athens. Athens is also home to several museums and cultural institutions, such as the National Archeological Museum, featuring the world's largest collection of ancient Greek antiquities, the Acropolis Museum, the Museum of Cycladic Art, the Benaki Museum and the Byzantine and Christian Museum. Athens was the host city of the first modern-day Olympic Games in 1896, and 108 years later it hosted the 2004 Summer Olympics, making it one of the few cities to have hosted the Olympics more than once.[17]
|
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|
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In Ancient Greek, the name of the city was Ἀθῆναι (Athênai, pronounced [atʰɛ̂ːnai̯] in Classical Attic) a plural. In earlier Greek, such as Homeric Greek, the name had been current in the singular form though, as Ἀθήνη (Athḗnē).[18] It was possibly rendered in the plural later on, like those of Θῆβαι (Thêbai) and Μυκῆναι (Μukênai). The root of the word is probably not of Greek or Indo-European origin,[19] and is possibly a remnant of the Pre-Greek substrate of Attica.[19] In antiquity, it was debated whether Athens took its name from its patron goddess Athena (Attic Ἀθηνᾶ, Athēnâ, Ionic Ἀθήνη, Athḗnē, and Doric Ἀθάνα, Athā́nā) or Athena took her name from the city.[20] Modern scholars now generally agree that the goddess takes her name from the city,[20] because the ending -ene is common in names of locations, but rare for personal names.[20] During the medieval period, the name of the city was rendered once again in the singular as Ἀθήνα. However, after the establishment of the modern Greek state, and partly due to the conservatism of the written language, Ἀθῆναι [aˈθine] became again the official name of the city and remained so until the abandonment of Katharevousa in the 1970s, when Ἀθήνα, Athína, became the official name.
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According to the ancient Athenian founding myth, Athena, the goddess of wisdom, competed against Poseidon, the God of the Seas, for patronage of the yet-unnamed city;[21] they agreed that whoever gave the Athenians the better gift would become their patron[21] and appointed Cecrops, the king of Athens, as the judge.[21] According to the account given by Pseudo-Apollodorus, Poseidon struck the ground with his trident and a salt water spring welled up.[21] In an alternative version of the myth from Vergil's Georgics, Poseidon instead gave the Athenians the first horse.[21] In both versions, Athena offered the Athenians the first domesticated olive tree.[21][22] Cecrops accepted this gift[21] and declared Athena the patron goddess of Athens.[21][22] Eight different etymologies, now commonly rejected, were proposed during the 17th century.[citation needed] Christian Lobeck proposed as the root of the name the word ἄθος (áthos) or ἄνθος (ánthos) meaning "flower", to denote Athens as the "flowering city". Ludwig von Döderlein proposed the stem of the verb θάω, stem θη- (tháō, thē-, "to suck") to denote Athens as having fertile soil.[23]
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In classical literature, the city was sometimes referred to as the City of the Violet Crown, first documented in Pindar's ἰοστέφανοι Ἀθᾶναι (iostéphanoi Athânai), or as τὸ κλεινὸν ἄστυ (tò kleinòn ásty, "the glorious city"). In medieval texts, variant names include Setines, Satine, and Astines, all derivations involving false splitting of prepositional phrases.[24] Today the caption η πρωτεύουσα (ī protévousa), "the capital", has become somewhat common. King Alphonse X of Castile in his 'General Storia' (Global History) cites etymology of Athens as: 'Athinai', 'the one without death', death here meaning ignorance.
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Athenians were called cicada-wearers (Ancient Greek: Τεττιγοφόροι) because they used to wear pins of golden cicadas. A symbol of being autochthon (earth-born), because the legendary founder of Athens, Erechtheus was an autochthon or of being musicians, because the cicada is a "musician" insect.[25]
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The oldest known human presence in Athens is the Cave of Schist, which has been dated to between the 11th and 7th millennia BC.[4] Athens has been continuously inhabited for at least 5,000 years.[26][27] By 1400 BC the settlement had become an important centre of the Mycenaean civilization and the Acropolis was the site of a major Mycenaean fortress, whose remains can be recognised from sections of the characteristic Cyclopean walls.[28] Unlike other Mycenaean centers, such as Mycenae and Pylos, it is not known whether Athens suffered destruction in about 1200 BC, an event often attributed to a Dorian invasion, and the Athenians always maintained that they were pure Ionians with no Dorian element. However, Athens, like many other Bronze Age settlements, went into economic decline for around 150 years afterwards.
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Iron Age burials, in the Kerameikos and other locations, are often richly provided for and demonstrate that from 900 BC onwards Athens was one of the leading centres of trade and prosperity in the region.[29] The leading position of Athens may well have resulted from its central location in the Greek world, its secure stronghold on the Acropolis and its access to the sea, which gave it a natural advantage over inland rivals such as Thebes and Sparta.
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By the 6th century BC, widespread social unrest led to the reforms of Solon. These would pave the way for the eventual introduction of democracy by Cleisthenes in 508 BC. Athens had by this time become a significant naval power with a large fleet, and helped the rebellion of the Ionian cities against Persian rule. In the ensuing Greco-Persian Wars Athens, together with Sparta, led the coalition of Greek states that would eventually repel the Persians, defeating them decisively at Marathon in 490 BC, and crucially at Salamis in 480 BC. However, this did not prevent Athens from being captured and sacked twice by the Persians within one year, after a heroic but ultimately failed resistance at Thermopylae by Spartans and other Greeks led by King Leonidas,[30] after both Boeotia and Attica fell to the Persians.
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The decades that followed became known as the Golden Age of Athenian democracy, during which time Athens became the leading city of Ancient Greece, with its cultural achievements laying the foundations for Western civilization. The playwrights Aeschylus, Sophocles and Euripides flourished in Athens during this time, as did the historians Herodotus and Thucydides, the physician Hippocrates, and the philosopher Socrates. Guided by Pericles, who promoted the arts and fostered democracy, Athens embarked on an ambitious building program that saw the construction of the Acropolis of Athens (including the Parthenon), as well as empire-building via the Delian League. Originally intended as an association of Greek city-states to continue the fight against the Persians, the league soon turned into a vehicle for Athens's own imperial ambitions. The resulting tensions brought about the Peloponnesian War (431–404 BC), in which Athens was defeated by its rival Sparta.
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By the mid-4th century BC, the northern Greek kingdom of Macedon was becoming dominant in Athenian affairs. In 338 BC the armies of Philip II defeated an alliance of some of the Greek city-states including Athens and Thebes at the Battle of Chaeronea, effectively ending Athenian independence. Later, under Rome, Athens was given the status of a free city because of its widely admired schools. The Roman emperor Hadrian, in the 2nd century CE, ordered the construction of a library, a gymnasium, an aqueduct which is still in use, several temples and sanctuaries, a bridge and financed the completion of the Temple of Olympian Zeus.
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By the end of Late Antiquity, Athens had shrunk due to sacks by the Herulians, Visigoths, and Early Slavs which caused massive destruction in the city. In this era, the first Christian churches were built in Athens, and the Parthenon and other temples were converted into churches. Athens expanded its settlement in the second half of the Middle Byzantine Period, in the 9th to 10th centuries CE, and was relatively prosperous during the Crusades, benefiting from Italian trade. After the Fourth Crusade the Duchy of Athens was established. In 1458 it was conquered by the Ottoman Empire and entered a long period of decline.
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Following the Greek War of Independence and the establishment of the Greek Kingdom, Athens was chosen as the capital of the newly independent Greek state in 1834, largely because of historical and sentimental reasons. At the time, it was reduced to a town of about 4,000 people in a loose swarm of houses along the foot of the Acropolis. The first King of Greece, Otto of Bavaria, commissioned the architects Stamatios Kleanthis and Eduard Schaubert to design a modern city plan fit for the capital of a state.
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The first modern city plan consisted of a triangle defined by the Acropolis, the ancient cemetery of Kerameikos and the new palace of the Bavarian king (now housing the Greek Parliament), so as to highlight the continuity between modern and ancient Athens. Neoclassicism, the international style of this epoch, was the architectural style through which Bavarian, French and Greek architects such as Hansen, Klenze, Boulanger or Kaftantzoglou designed the first important public buildings of the new capital. In 1896, Athens hosted the first modern Olympic Games. During the 1920s a number of Greek refugees, expelled from Asia Minor after the Greco-Turkish War, swelled Athens's population; nevertheless it was most particularly following World War II, and from the 1950s and 1960s, that the population of the city exploded, and Athens experienced a gradual expansion.
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In the 1980s it became evident that smog from factories and an ever-increasing fleet of automobiles, as well as a lack of adequate free space due to congestion, had evolved into the city's most important challenge. A series of anti-pollution measures taken by the city's authorities in the 1990s, combined with a substantial improvement of the city's infrastructure (including the Attiki Odos motorway, the expansion of the Athens Metro, and the new Athens International Airport), considerably alleviated pollution and transformed Athens into a much more functional city. In 2004 Athens hosted the 2004 Summer Olympics.
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Tondo of the Aison Cup, showing the victory of Theseus over the Minotaur in the presence of Athena. Theseus was responsible, according to the myth, for the synoikismos ("dwelling together")—the political unification of Attica under Athens.
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The earliest coinage of Athens, c. 545–525/15 BC
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Coat of Arms of the Duchy of Athens during the rule of the de la Roche family (13th century)
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The Roman Agora and the Gate of Athena in Plaka district.
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The Temple of Olympian Zeus with river Ilisos by Edward Dodwell, 1821
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The Entry of King Otto in Athens, Peter von Hess, 1839.
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The Stadiou Street in Central Athens in 1908.
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Athens sprawls across the central plain of Attica that is often referred to as the Athens Basin or the Attica Basin (Greek: Λεκανοπέδιο Αθηνών/Αττικής). The basin is bounded by four large mountains: Mount Aigaleo to the west, Mount Parnitha to the north, Mount Pentelicus to the northeast and Mount Hymettus to the east.[31] Beyond Mount Aegaleo lies the Thriasian plain, which forms an extension of the central plain to the west. The Saronic Gulf lies to the southwest. Mount Parnitha is the tallest of the four mountains (1,413 m (4,636 ft)),[32] and has been declared a national park.
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Athens is built around a number of hills. Lycabettus is one of the tallest hills of the city proper and provides a view of the entire Attica Basin. The meteorology of Athens is deemed to be one of the most complex in the world because its mountains cause a temperature inversion phenomenon which, along with the Greek Government's difficulties controlling industrial pollution, was responsible for the air pollution problems the city has faced.[27] This issue is not unique to Athens; for instance, Los Angeles and Mexico City also suffer from similar atmospheric inversion problems.[27]
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The Cephissus river, the Ilisos and the Eridanos stream are the historical rivers of Athens.
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By the late 1970s, the pollution of Athens had become so destructive that according to the then Greek Minister of Culture, Constantine Trypanis, "...the carved details on the five the caryatids of the Erechtheum had seriously degenerated, while the face of the horseman on the Parthenon's west side was all but obliterated."[33] A series of measures taken by the authorities of the city throughout the 1990s resulted in the improvement of air quality; the appearance of smog (or nefos as the Athenians used to call it) has become less common.
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Measures taken by the Greek authorities throughout the 1990s have improved the quality of air over the Attica Basin. Nevertheless, air pollution still remains an issue for Athens, particularly during the hottest summer days. In late June 2007,[34] the Attica region experienced a number of brush fires,[34] including a blaze that burned a significant portion of a large forested national park in Mount Parnitha,[35] considered critical to maintaining a better air quality in Athens all year round.[34] Damage to the park has led to worries over a stalling in the improvement of air quality in the city.[34]
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The major waste management efforts undertaken in the last decade (particularly the plant built on the small island of Psytalia) have greatly improved water quality in the Saronic Gulf, and the coastal waters of Athens are now accessible again to swimmers. In January 2007, Athens faced a waste management problem when its landfill near Ano Liosia, an Athenian suburb, reached capacity.[36] The crisis eased by mid-January when authorities began taking the garbage to a temporary landfill.[36]
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Athens ranks in the lowest percentage for the risk on frequency and severity of terrorist attacks according to the EU Global Terrorism Database (EIU 2007-2016 calculations). The city also ranked 35th in Digital Security, 21st on Health Security, 29th on Infrastructure Security and 41st on Personal Security globally in a 2017 The Economist Intelligence Unit report.[37] It also ranks as a very safe city (39th globally out of 162 cities overall) on the ranking of the safest and most dangerous countries.[38] A 2019 crime index from Numbeo places Athens at 130th position, rating safer than Tampa, Florida or Dublin, Ireland.[39] According to a Mercer 2019 Quality of Living Survey, Athens ranks 89th on the Mercer Quality of Living Survey ranking.[40]
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Athens has a hot-summer Mediterranean climate (Köppen climate classification: Csa). The dominant feature of Athens' climate is alternation between prolonged hot and dry summers and mild winters with moderate rainfall.[41] With an average of 455 millimetres (17.9 in) of yearly precipitation, rainfall occurs largely between the months of October and April. July and August are the driest months when thunderstorms occur sparsely. Furthermore, some coastal areas of Athens, known as the Athens Riviera have a hot semi-arid climate (BSh) according to the climate atlas published by the Hellenic National Meteorological Service (H.N.M.S).[42]
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Owing to the rain shadow of the Pindus Mountains, annual precipitation of Athens is lower than most other parts of Greece, especially western Greece. As an example, Ioannina receives around 1,300 mm (51 in) per year, and Agrinio around 800 mm (31 in) per year. Daily average highs for July have been measured around 34 °C or 93 °F in downtown Athens, but some parts of the city may be even hotter for the higher density of buildings, such as the center,[43] in particular, western areas due to a combination of industrialization and a number of natural factors, knowledge of which has existed since the mid-19th century.[44][45][46]
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Athens is affected by the urban heat island effect in some areas which is caused by human activity,[47][48] altering its temperatures compared to the surrounding rural areas,[49][50][51][52] and leaving detrimental effects on energy usage, expenditure for cooling,[53][54] and health.[48] The urban heat island of the city has also been found to be partially responsible for alterations of the climatological temperature time-series of specific Athens meteorological stations, because of its impact on the temperatures and the temperature trends recorded by some meteorological stations.[55][56][57][58][59] On the other hand, specific meteorological stations, such as the National Garden station and Thiseio meteorological station, are less affected or do not experience the urban heat island.[49][60]
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Athens holds the World Meteorological Organization record for the highest temperature ever recorded in Europe, at 48 °C (118.4 °F), which was recorded in the Elefsina and Tatoi suburbs of Athens on 10 July 1977.[61]
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The Municipality of Athens, the City Centre of the Athens Urban Area, is divided into several districts: Omonoia, Syntagma, Exarcheia, Agios Nikolaos, Neapolis, Lykavittos, Lofos Strefi, Lofos Finopoulou, Lofos Filopappou, Pedion Areos, Metaxourgeio, Aghios Kostantinos, Larissa Station, Kerameikos, Psiri, Monastiraki, Gazi, Thission, Kapnikarea, Aghia Irini, Aerides, Anafiotika, Plaka, Acropolis, Pnyka, Makrygianni, Lofos Ardittou, Zappeion, Aghios Spyridon, Pangrati, Kolonaki, Dexameni, Evaggelismos, Gouva, Aghios Ioannis, Neos Kosmos, Koukaki, Kynosargous, Fix, Ano Petralona, Kato Petralona, Rouf, Votanikos, Profitis Daniil, Akadimia Platonos, Kolonos, Kolokynthou, Attikis Square, Lofos Skouze, Sepolia, Kypseli, Aghios Meletios, Nea Kypseli, Gyzi, Polygono, Ampelokipoi, Panormou-Gerokomeio, Pentagono, Ellinorosson, Nea Filothei, Ano Kypseli, Tourkovounia-Lofos Patatsou, Lofos Elikonos, Koliatsou, Thymarakia, Kato Patisia, Treis Gefyres, Aghios Eleftherios, Ano Patisia, Kypriadou, Menidi, Prompona, Aghios Panteleimonas, Pangrati, Goudi, Vyronas and Ilisia.
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Parnitha National Park is punctuated by well-marked paths, gorges, springs, torrents and caves dotting the protected area. Hiking and mountain-biking in all four mountains are popular outdoor activities for residents of the city. The National Garden of Athens was completed in 1840 and is a green refuge of 15.5 hectares in the centre of the Greek capital. It is to be found between the Parliament and Zappeion buildings, the latter of which maintains its own garden of seven hectares.
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Parts of the City Centre have been redeveloped under a masterplan called the Unification of Archeological Sites of Athens, which has also gathered funding from the EU to help enhance the project.[68][69] The landmark Dionysiou Areopagitou Street has been pedestrianised, forming a scenic route. The route starts from the Temple of Olympian Zeus at Vasilissis Olgas Avenue, continues under the southern slopes of the Acropolis near Plaka, and finishes just beyond the Temple of Hephaestus in Thiseio. The route in its entirety provides visitors with views of the Parthenon and the Agora (the meeting point of ancient Athenians), away from the busy City Centre.
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The hills of Athens also provide green space. Lycabettus, Philopappos hill and the area around it, including Pnyx and Ardettos hill, are planted with pines and other trees, with the character of a small forest rather than typical metropolitan parkland. Also to be found is the Pedion tou Areos (Field of Mars) of 27.7 hectares, near the National Archaeological Museum.
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Athens' largest zoo is the Attica Zoological Park, a 20-hectare (49-acre) private zoo located in the suburb of Spata. The zoo is home to around 2000 animals representing 400 species, and is open 365 days a year. Smaller zoos exist within public gardens or parks, such as the zoo within the National Garden of Athens.
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The Athens Metropolitan Area consists of 58[15] densely populated municipalities, sprawling around the Municipality of Athens (the City Centre) in virtually all directions. For the Athenians, all the urban municipalities surrounding the City Centre are called suburbs. According to their geographic location in relation to the City of Athens, the suburbs are divided into four zones; the northern suburbs (including Agios Stefanos, Dionysos, Ekali, Nea Erythraia, Kifissia, Kryoneri, Maroussi, Pefki, Lykovrysi, Metamorfosi, Nea Ionia, Nea Filadelfeia, Irakleio, Vrilissia, Melissia, Penteli, Chalandri, Agia Paraskevi, Gerakas, Pallini, Galatsi, Psychiko and Filothei); the southern suburbs (including Alimos, Nea Smyrni, Moschato, Tavros, Agios Ioannis Rentis, Kallithea, Piraeus, Agios Dimitrios, Palaio Faliro, Elliniko, Glyfada, Lagonisi, Saronida, Argyroupoli, Ilioupoli, Varkiza, Voula, Vari and Vouliagmeni); the eastern suburbs (including Zografou, Dafni, Vyronas, Kaisariani, Cholargos and Papagou); and the western suburbs (including Peristeri, Ilion, Egaleo, Koridallos, Agia Varvara, Keratsini, Perama, Nikaia, Drapetsona, Chaidari, Petroupoli, Agioi Anargyroi, Ano Liosia, Aspropyrgos, Eleusina, Acharnes and Kamatero).
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The Athens city coastline, extending from the major commercial port of Piraeus to the southernmost suburb of Varkiza for some 25 km (20 mi),[70] is also connected to the City Centre by tram.
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In the northern suburb of Maroussi, the upgraded main Olympic Complex (known by its Greek acronym OAKA) dominates the skyline. The area has been redeveloped according to a design by the Spanish architect Santiago Calatrava, with steel arches, landscaped gardens, fountains, futuristic glass, and a landmark new blue glass roof which was added to the main stadium. A second Olympic complex, next to the sea at the beach of Palaio Faliro, also features modern stadia, shops and an elevated esplanade. Work is underway to transform the grounds of the old Athens Airport – named Elliniko – in the southern suburbs, into one of the largest landscaped parks in Europe, to be named the Hellenikon Metropolitan Park.[71]
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Many of the southern suburbs (such as Alimos, Palaio Faliro, Elliniko, Glyfada, Voula, Vouliagmeni and Varkiza) known as the Athens Riviera, host a number of sandy beaches, most of which are operated by the Greek National Tourism Organisation and require an entrance fee. Casinos operate on both Mount Parnitha, some 25 km (16 mi)[72] from downtown Athens (accessible by car or cable car), and the nearby town of Loutraki (accessible by car via the Athens – Corinth National Highway, or the suburban rail service Proastiakos).
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The large City Centre (Greek: Κέντρο της Αθήνας) of the Greek capital falls directly within the Municipality of Athens or Athens Municipality (Greek: Δήμος Αθηναίων)—also City of Athens. Athens Municipality is the largest in population size in Greece. Piraeus also forms a significant city centre on its own,[73] within the Athens Urban Area and it is the second largest in population size within it, with Peristeri, Kallithea and Kypseli following.
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The Athens Urban Area (Greek: Πολεοδομικό Συγκρότημα Αθηνών), also known as Urban Area of the Capital (Greek: Πολεοδομικό Συγκρότημα Πρωτεύουσας), today consists of 40 municipalities, 35 of which make up what is referred to as the Greater Athens municipalities, located within 4 regional units (North Athens, West Athens, Central Athens, South Athens); and a further 5, which make up the Greater Piraeus municipalities, located within the regional unit of Piraeus as mentioned above. The densely built up urban area of the Greek capital sprawls across 412 km2 (159 sq mi)[14] throughout the Attica Basin and has a total population of 3,074,160 (in 2011).
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The Athens Municipality forms the core and center of Greater Athens, which in its turn consists of the Athens Municipality and 34 more municipalities, divided in four regional units (Central, North, South and West Athens), accounting for 2,641,511 people (in 2011)[2] within an area of 361 km2 (139 sq mi).[14] Until 2010, these four regional units made up the abolished Athens Prefecture. The municipality of Piraeus, the historic Athenian port, with its 4 suburban municipalities make up the regional unit of Piraeus, which in turn forms Greater Piraeus.
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Greater Athens and Greater Piraeus with part of East[74] and West Attica[75] regional units combined make up the continuous Athens Urban Area,[75][76][77] also called the "Urban Area of the Capital" or simply "Athens" (the most common use of the term), spanning over 412 km2 (159 sq mi),[78] with a population of 3,090,508 people as of 2011. The Athens Urban Area is considered to form the city of Athens as a whole, despite its administrative divisions, which is the largest in Greece and one of the most populated urban areas in Europe.
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The Athens Metropolitan Area (Greek: Μητροπολιτική Περιοχή της Αθήνας), also commonly known in English as Greater Athens,[79] spans 2,928.717 km2 (1,131 sq mi) within the Attica region and includes a total of 58 municipalities, which are organized in seven regional units (those outlined above, along with East Attica and West Attica), having reached a population of 3,737,550 based on the preliminary results of the 2011 census. Athens and Piraeus municipalities serve as the two metropolitan centres of the Athens Metropolitan Area.[80] There are also some inter-municipal centres serving specific areas. For example, Kifissia and Glyfada serve as inter-municipal centres for northern and southern suburbs respectively.
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The Municipality of Athens has an official population of 664,046 people.[2] The four regional units that make up what is referred to as Greater Athens have a combined population of 2,640,701. They together with the regional unit of Piraeus (Greater Piraeus) make up the dense Athens Urban Area which reaches a total population of 3,090,508 inhabitants (in 2011).[15] As Eurostat the FUA of Athens had in 2013 3,828,434 inhabitants, being apparently decreasing compared with the pre-economic crisis date of 2009 (4,164,175)[16]
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The municipality (Center) of Athens is the most populous in Greece, with a population of 664,046 people (in 2011)[2] and an area of 38.96 km2 (15.04 sq mi),[13] forming the core of the Athens Urban Area within the Attica Basin. The incumbent Mayor of Athens is Kostas Bakoyannis of New Democracy. The municipality is divided into seven municipal districts which are mainly used for administrative purposes.
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As of the 2011 census, the population for each of the seven municipal districts of Athens is as follows:[81]
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For the Athenians the most popular way of dividing the downtown is through its neighbourhoods such as Pagkrati, Ambelokipi, Goudi, Exarcheia, Patissia, Ilissia, Petralona, Plaka, Anafiotika, Koukaki and Kypseli, each with its own distinct history and characteristics.
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The Athens Metropolitan Area, with an area of 2,928.717 km2 (1,131 sq mi) and inhabited by 3,753,783 people in 2011,[2] consists of the Athens Urban Area with the addition of the towns and villages of East and West Attica, which surround the dense urban area of the Greek capital. It actually sprawls over the whole peninsula of Attica, which is the best part of the region of Attica, excluding the islands.
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Mycenean Athens in 1600–1100 BC could have reached the size of Tiryns; that would put the population at the range of 10,000–15,000.[82] During the Greek Dark Ages the population of Athens was around 4,000 people. In 700 BC the population grew to 10,000. In 500 BC the area probably contained 200,000 people. During the classical period the city's population is estimated from 150,000–350,000 and up to 610,000 according to Thucydides. When Demetrius of Phalerum conducted a population census in 317 BC the population was 21,000 free citizens, plus 10,000 resident aliens and 400,000 slaves. This suggests a total population of 431,000.[83][84][85][86][87][88][89][90][91] This figure is highly suspect because of the lopsided number of slaves and does not include free women and children and resident foreigners: an estimated based on Thucydides is: 40,000 male citizens, 100,000 family members, 70,000 metics (resident foreigners) and 150,000-400,000 slaves. However the numbers would include all of Athenian territory which consisted almost entirely of modern Attica region expect the territory of the city-state of Megaris and the island section.
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The ancient site of the main city is centred on the rocky hill of the acropolis. In the whole of Athenian territory they existed many towns. Acharnae, Afidnes, Cytherus, Colonus, Corydallus, Cropia, Decelea, Euonymos, Vravron among others was important towns in Athenian countryside. The new port of Piraeus was a prototype harbour with the infrastructure and housing located in the site between modern passenger section of the port (named Kantharos in ancient times) and the Pasalimani harbour (named Zea in ancient times). The old one Phaliro was in the site of modern Palaio Faliro and gradually declined after the construction of the new prototype port but remained as a minor port and important settlement with historic significance in late Classical times. The rapid expansion of the modern city, which continues to this day, was initiated in the 1950s and 1960s, because of Greece's transition from an agricultural to an industrial nation.[92] The expansion is now particularly toward the East and North East (a tendency greatly related to the new Eleftherios Venizelos International Airport and the Attiki Odos, the freeway that cuts across Attica). By this process Athens has engulfed many former suburbs and villages in Attica, and continues to do so. The table below shows the historical population of Athens in recent times.
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Athens became the capital of Greece in 1834, following Nafplion, which was the provisional capital from 1829. The municipality (City) of Athens is also the capital of the Attica region. The term Athens can refer either to the Municipality of Athens, to Greater Athens, or to the entire Athens Urban Area.
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The Hellenic Parliament
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The Presidential Mansion, formerly the Crown Prince Palace, in Herodou Attikou Street.
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The Maximos Mansion, official office of the Prime Minister of the Hellenic Republic, in Herodou Attikou Street.
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The Athens City Hall in Kotzia Square was designed by Panagiotis Kolkas and completed in 1874.[98]
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The Embassy of France in Vasilissis Sofias Avenue.
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The Italian Embassy in Vasilissis Sofias Avenue.
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Athens is twinned with:[99]
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Athens is the financial capital of Greece. According to data from 2014, Athens as a metropolitan economic area produced 130 billion US-dollars as GDP in PPP, which consists nearly a half of the production for the whole country. In the list with the strongest economic metropoles of the world Athens was ranked that year 102nd, while the GDP per capita for the same year was 32,000 US-dollars.[112]
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Athens is one of the major economic centres in south-eastern Europe and is considered as a regional economic power in Europe generally. The port of Piraeus, where big investments by COSCO have already been delivered during the recent decade, the completion of the new Cargo Centre in Thriasion,[113] the expansion of the Athens Metro and the Athens Tram, as well as the projected metropolitan park in Elliniko and other economic projects are the economic landmarks of the upcoming years.
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Important Greek companies such as Hellenic Aerospace Industry, Mytilineos Holdings, Titan Cement, Hellenic Petroleum, Papadopoulos E.J., Folli Follie, Jumbo S.A., OPAP, and COSMOTE have their headquarters in the metropolitan area of Athens. Multinational companies such as Ericsson, Sony, Siemens, Motorola, Samsung, Microsoft, Novartis, Mondelez, Coca-Cola, etc. have their regional research and development headquarters also there.
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The banking sector is represented by National Bank of Greece, Alpha Bank, Eurobank, and Piraeus Bank, while the Bank of Greece is also situated in the City Centre. The Athens Stock Exchange, the only in Greece, has been severely hit by the Greek government-debt crisis and the decision of the government to proceed into capital controls during summer 2015. As a whole the economy of Athens and Greece has been severely hit with today's data showing a change from long recession to growth of 1.4% in 2017.[114]
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Tourism is also a great contributor for the economy of the city, which is considered as one of the top destinations in Europe for city-break tourism and is also the gateway for excursions to the islands or the mainland. Greece attracted 26.5 million visitors in 2015, 30.1 million visitors in 2017 and over 33 million in 2018, making Greece one of the most visited countries in Europe and the world, and contributing 18% to the nation's Gross Domestic Product. Athens welcomed more than 5 million tourists in 2018 and 1,4 million of them were "city-breakers" (in 2013 the city-breakers were only 220.000).[115]
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Athens is the major transportation hub of Greece. They combines the largest Airport in Greece as also the largest port in Greece which also the largest port in Mediterranean in containers transport and the largest passenger port in Europe. Is also major national hub for Intercity (Ktel) and international buses as well as for domestic and international rail transport. Public transport is serviced by a variety of transportation means, forming the largest mass transit system of Greece. The Athens Mass Transit System consists of a large bus fleet, a trolleybus fleet that mainly serves Athens's city center, the city's Metro, a Proastiakos service[116] and a tram network, connecting the southern suburbs to the City Centre.[117]
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OSY (Greek: ΟΣΥ) (Odikes Sygkoinonies S.A.) which is subsidiary company of OASA (Athens urban transport organisation), is the main operator of buses and trolleybusses in Athens. Its network consists of about 300 bus lines and 22 trolleybus lines which span the Athens Metropolitan Area,[118] with a fleet of 1,839 buses and 366 trolleybuses.[119] Of those 1,839 buses 416 run on compressed natural gas,[119] making up the largest fleet of natural gas-powered buses in Europe[120] and all trolleybusses are equipped to enable them to run on diesel in case of power failure.[121]
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International links are provided by a number of private companies. National and regional bus links are provided by KTEL from two InterCity Bus Terminals, Kifissos Bus Terminal A and Liosion Bus Terminal B, both located in the north-western part of the city. Kifissos provides connections towards Peloponnese, North Greece, West Greece and some Ionian Islands, whereas Liosion is used for most of Central Greece.
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The Athens Metro is operated by STASY S.A (Greek: ΣΤΑΣΥ) (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation) and provides public transport throughout the Athens Urban Area. While its main purpose is transport, it also houses Greek artifacts found during construction of the system.[122] The Athens Metro has an operating staff of 387 and runs three metro lines, namely the line 1 (Green Line), line 2 (Red Line) and line 3 (Blue Line) lines, of which the first was constructed
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in 1869, and the other two largely during the 1990s, with the initial new sections opened in January 2000. Line 1 mostly runs at ground level and the other two (Line 2 & 3) routes run entirely underground. A fleet of 42 trains, using 252 carriages, operates on the network,[123] with a daily occupancy of 1,353,000 passengers.[124]
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The line 1 (Green Line) serves 24 stations, and is the oldest line of the Athens metro network. It runs from Piraeus station to Kifissia station and covers a distance of 25.6-kilometre (15.9 mi). There are transfer connections with the Blue Line 3 at Monastiraki station and with the Red Line 2 at Omonoia and Attiki stations.
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The line 2 (Red Line) runs from Anthoupoli station to Elliniko station and covers a distance of 17.5 km (10.9 mi).[123] The line connects the western suburbs of Athens with the southeast suburbs, passing through the center of Athens. The Red Line has transfer connections with the Green Line 1 at Attiki and Omonoia stations. There are also transfer connections with the Blue Line 3 at Syntagma station and with the tram at Syntagma, Syngrou Fix and Neos Kosmos stations.
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The line 3 (Blue Line) runs from Agia Marina and Egaleo stations, through the central Monastiraki and Syntagma stations to Doukissis Plakentias avenue in the northeastern suburb of Halandri, covering a distance of 16 km (10 mi).[123] It then ascends to ground level and continues to Eleftherios Venizelos International Airport, using the Suburban Railway infrastructure, extending its total length to 39 km (24 mi).[123] The spring 2007 extension from Monastiraki westwards to Egaleo connected some of the main night life hubs of the city, namely those of Gazi (Kerameikos station) with Psirri (Monastiraki station) and the City Centre (Syntagma station). Extensions are under construction to the west southwest suburbs of Athens, as far as the port of Piraeus. The new stations will be Agia Barvara, Koridallos, Nikaia, Maniatika, Piraeus and Dimotiko Theatro. The completed extension will be ready in 2022 (the first three stations will open in summer 2020), connecting the biggest port of Greece, Piraeus Port, with the biggest airport of Greece the Athens International Airport.
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The Athens commuter rail service, referred to as the "Proastiakós", connects Eleftherios Venizelos International Airport to the city of Kiato, 106 km (66 mi)[125] west of Athens, via Larissa station, the city's central rail station and the port of Piraeus. The length of Athens's commuter rail network extends to 120 km (75 mi),[125] and is expected to stretch to 281 km (175 mi) by 2010.[125]
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Athens tram is operated by STASY S.A (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation). It has a fleet of 35 Sirio type vehicles[126] which serve 48 stations,[126] employ 345 people with an average daily occupancy of 65,000 passengers.[126] The tram network spans a total length of 27 km (17 mi) and covers ten Athenian suburbs.[126] The network runs from Syntagma Square to the southwestern suburb of Palaio Faliro, where the line splits in two branches; the first runs along the Athens coastline toward the southern suburb of Voula, while the other heads toward Neo Faliro. The network covers the majority of the Athens coastline.[127] Further extensions are planned towards the major commercial port of Piraeus.[126] The expansion to Piraeus will include 12 new stations, increase the overall length of tram route by 5.4 km (3 mi), and increase the overall transportation network.[128]
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Athens is served by the Athens International Airport (ATH), located near the town of Spata, in the eastern Messoghia plain, some 35 km (22 mi) east of center of Athens.[129] The airport, awarded the "European Airport of the Year 2004" Award,[130] is intended as an expandable hub for air travel in southeastern Europe and was constructed in 51 months, costing 2.2 billion euros. It employs a staff of 14,000.[130]
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The airport is served by the Metro, the suburban rail, buses to Piraeus port, Athens' City Centre, Liosion and Kifisos Intercity bus stations and Elliniko metro's line 2 southern terminal, and also taxis. The airport accommodates 65 landings and take-offs per hour,[129] with its 24-passenger boarding bridges,[129] 144 check-in counters and broader 150,000 m2 (1,614,587 sq ft) main terminal;[129] and a commercial area of 7,000 m2 (75,347 sq ft) which includes cafés, duty-free shops,[130] and a small museum.
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In 2018, the airport handled 24,135,736 a huge increase over the last 4 years. In 2014, the airport handled 15,196,369 passengers, an increase of 21.2% over the previous year of 2013.[131] Of those 15,196,369 passengers, 5,267,593 passed through the airport for domestic flights,[132] and 9,970,006 passengers travelled through for international flights.[132] Beyond the dimensions of its passenger capacity, ATH handled 205,294 total flights in 2007, or approximately 562 flights per day.[133]
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Athens is the hub of the country's national railway system (OSE), connecting the capital with major cities across Greece and abroad (Istanbul, Sofia, Belgrade and Bucharest). The Port of Piraeus is the largest port in Greece and one of the largest in Europe. It is the largest container port in East Mediterranean Sea Basin. It is also the busiest passenger port in Europe and one of the largest passenger ports in the world. Connects Athens to the numerous Greek islands of the Aegean Sea, with ferries departing, while also serving the cruise ships that arrive.[134][135][136] Rafina and Lavrio act as alternative ports of Athens, connects the city with numerous Greek islands of the Aegean Sea, Evia and Cesme in Turkey,[137][138] while also serving the cruise ships that arrive.
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Two main motorways of Greece begin in Athens, namely the A1/E75, heading north towards Greece's second largest city, Thessaloniki; and the border crossing of Evzones and the A8/E94 heading west, towards Greece's third largest city, Patras, which incorporated the GR-8A. Before their completion much of the road traffic used the GR-1 and the GR-8.
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Athens' Metropolitan Area is served by the motorway network of the Attiki Odos toll-motorway (code: A6). Its main section extends from the western industrial suburb of Elefsina to Athens International Airport; while two beltways, namely the Aigaleo Beltway (A65) and the Hymettus Beltway (A64) serve parts of western and eastern Athens respectively. The span of the Attiki Odos in all its length is 65 km (40 mi),[139] making it the largest metropolitan motorway network in all of Greece.
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Located on Panepistimiou Street, the old campus of the University of Athens, the National Library, and the Athens Academy form the "Athens Trilogy" built in the mid-19th century. Most of the university's workings have been moved to a much larger, modern campus located in the eastern suburb of Zografou. The second higher education institution in the city is the Athens Polytechnic School, found in Patission Street. This was the location where on 17 November 1973, more than 13 students were killed and hundreds injured inside the university during the Athens Polytechnic uprising,[140] against the military junta that ruled the nation from 21 April 1967 until 23 July 1974.
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Other universities that lie within Athens are the Athens University of Economics and Business, the Panteion University, the Agricultural University of Athens and the University of Piraeus. There are overall ten state-supported Institutions of Higher (or Tertiary) education located in the Athens Urban Area, these are by chronological order: Athens School of Fine Arts (1837), National Technical University of Athens (1837), National and Kapodistrian University of Athens (1837), Agricultural University of Athens (1920), Athens University of Economics and Business (1920), Panteion University of Social and Political Sciences (1927), University of Piraeus (1938), Harokopio University of Athens (1990), School of Pedagogical and Technological Education (2002), University of West Attica (2018). There are also several other private colleges, as they called formally in Greece, as the establishment of private universities is prohibited by the constitution. Many of them are accredited by a foreign state or university such as the American College of Greece and the Athens Campus of the University of Indianapolis.[141]
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The city is a world centre of archaeological research. Along with national institutions, such as the Athens University and the Archaeological Society, there are multiple archaeological Museums including the National Archaeological Museum, the Cycladic Museum, the Epigraphic Museum, the Byzantine & Christian Museum, as well as museums at the ancient Agora, Acropolis, Kerameikos, and the Kerameikos Archaeological Museum. The city is also home to the Demokritos laboratory for Archaeometry, alongside regional and national archaeological authorities that form part of the Greek Department of Culture.
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Athens hosts 17 Foreign Archaeological Institutes which promote and facilitate research by scholars from their home countries. As a result, Athens has more than a dozen archaeological libraries and three specialized archaeological laboratories, and is the venue of several hundred specialized lectures, conferences and seminars, as well as dozens of archaeological exhibitions, each year. At any given time, hundreds of international scholars and researchers in all disciplines of archaeology are to be found in the city.
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Athens incorporates architectural styles ranging from Greco-Roman and Neoclassical to modern times. They are often to be found in the same areas, as Athens is not marked by a uniformity of architectural style. A visitor will quickly notice the absence of tall buildings: Athens has very strict height restriction laws in order to ensure the Acropolis hill is visible throughout the city. There is evidence of continuity in elements of the architectural environment through the city's history, although this is still a contested subject [142].
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For the greatest part of the 19th century Neoclassicism dominated Athens, as well as some deviations from it such as Eclecticism, especially in the early 20th century. Thus, the Old Royal Palace was the first important public building to be built, between 1836 and 1843. Later in the mid and late 19th century, Theophil Freiherr von Hansen and Ernst Ziller took part in the construction of many neoclassical buildings such as the Athens Academy and the Zappeion Hall. Ziller also designed many private mansions in the centre of Athens which gradually became public, usually through donations, such as Schliemann's Iliou Melathron.
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Beginning in the 1920s, modern architecture including Bauhaus and Art Deco began to exert an influence on almost all Greek architects, and buildings both public and private were constructed in accordance with these styles. Localities with a great number of such buildings include Kolonaki, and some areas of the centre of the city; neighbourhoods developed in this period include Kypseli.[143]
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In the 1950s and 1960s during the extension and development of Athens, other modern movements such as the International style played an important role. The centre of Athens was largely rebuilt, leading to the demolition of a number of neoclassical buildings. The architects of this era employed materials such as glass, marble and aluminium, and many blended modern and classical elements.[144] After World War II, internationally known architects to have designed and built in the city included Walter Gropius, with his design for the US Embassy, and, among others, Eero Saarinen, in his postwar design for the east terminal of the Ellinikon Airport.
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All over the city can be found several statues or busts. Apart from the neoclassicals by Leonidas Drosis at the Academy of Athens (Plato, Socrates, Apollo, Athena), other notable include the statue of Theseus by Georgios Fytalis at Thiseion, of philhellenes like Lord Byron, George Canning and William Gladstone, the equestrian statue of Theodoros Kolokotronis by Lazaros Sochos in front of the Old Parliament, statues of Ioannis Kapodistrias, Rigas Feraios and Adamantios Korais at the University, of Evangelos Zappas and Konstantinos Zappas at Zappeion, of Ioannis Varvakis at the National Garden, the "woodbreaker" by Dimitrios Filippotis, the equestrian statue of Alexandros Papagos at Papagou district and various busts of fighters of Greek independence at the Pedion tou Areos. A significant landmark is also the Tomb of the Unknown Soldier in Syntagma.
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Athens' most important museums include:
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Athens has been a destination for travellers since antiquity. Over the past decade, the city's infrastructure and social amenities have improved, in part because of its successful bid to stage the 2004 Olympic Games. The Greek Government, aided by the EU, has funded major infrastructure projects such as the state-of-the-art Eleftherios Venizelos International Airport,[145] the expansion of the Athens Metro system,[68] and the new Attiki Odos Motorway.[68]
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Athens was voted as the third best European city to visit in 2015 by European Best Destination. More than 240,000 people voted.
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Athens is home to 148 theatrical stages, more than any other city in the world, including the ancient Odeon of Herodes Atticus, home to the Athens Festival, which runs from May to October each year.[146][147] In addition to a large number of multiplexes, Athens plays host to open air garden cinemas. The city also supports music venues, including the Athens Concert Hall (Megaro Moussikis), which attracts world class artists.[148] The Athens Planetarium,[149] located in Andrea Syngrou Avenue, in Palaio Faliro[150] is one of the largest and best equipped digital planetaria in the world.[151] The Stavros Niarchos Foundation Cultural Center, inaugurated in 2016, will house the National Library of Greece and the Greek National Opera.[152]
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The most successful songs during the period 1870–1930 were the so-called Athenian serenades (Αθηναϊκές καντάδες), based on the Heptanesean kantádhes (καντάδες 'serenades'; sing.: καντάδα) and the songs performed on stage (επιθεωρησιακά τραγούδια 'theatrical revue songs') in revues, musical comedies, operettas and nocturnes that were dominating Athens' theatre scene.
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Notable composers of operettas or nocturnes were Kostas Giannidis, Dionysios Lavrangas, Nikos Hatziapostolou, while Theophrastos Sakellaridis' The Godson remains probably the most popular operetta. Despite the fact that the Athenian songs were not autonomous artistic creations (in contrast with the serenades) and despite their original connection with mainly dramatic forms of Art, they eventually became hits as independent songs. Notable actors of Greek operettas, who made also a series of melodies and songs popular at that time, include Orestis Makris, Kalouta sisters, Vasilis Avlonitis, Afroditi Laoutari, Eleni Papadaki, Marika Nezer, Marika Krevata and others. After 1930, wavering among American and European musical influences as well as the Greek musical tradition. Greek composers begin to write music using the tunes of the tango, waltz, swing, foxtrot, some times combined with melodies in the style of Athenian serenades' repertory. Nikos Gounaris was probably the most renowned composer and singer of the time.
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In 1923, after the population exchange between Greece and Turkey, many ethnic Greeks from Asia Minor fled to Athens as a result of the Greco-Turkish War. They settled in poor neighborhoods and brought with them Rebetiko music, making it popular also in Greece, which became later the base for the Laïko music. Other forms of song popular today in Greece are elafrolaika, entechno, dimotika, and skyladika.[153] Greece's most notable, and internationally famous, composers of Greek song, mainly of the entechno form, are Manos Hadjidakis and Mikis Theodorakis. Both composers have achieved fame abroad for their composition of film scores.[153]
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Athens has a long tradition in sports and sporting events, serving as home to the most important clubs in Greek sport and housing a large number of sports facilities. The city has also been host to sports events of international importance.
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Athens has hosted the Summer Olympic Games twice, in 1896 and 2004. The 2004 Summer Olympics required the development of the Athens Olympic Stadium, which has since gained a reputation as one of the most beautiful stadiums in the world, and one of its most interesting modern monuments.[154] The biggest stadium in the country, it hosted two finals of the UEFA Champions League, in 1994 and 2007. Athens' other major stadium, located in the Piraeus area, is the Karaiskakis Stadium, a sports and entertainment complex, host of the 1971 UEFA Cup Winners' Cup Final.
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Athens has hosted the EuroLeague final three times, the first in 1985 and second in 1993, both at the Peace and Friendship Stadium, most known as SEF, a large indoor arena,[155] and the third time in 2007 at the Olympic Indoor Hall. Events in other sports such as athletics, volleyball, water polo etc., have been hosted in the capital's venues.
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Athens is home to three European multi-sport clubs: Olympiacos, Panathinaikos, AEK Athens. In football, Olympiacos have dominated the domestic competitions, Panathinaikos made it to the 1971 European Cup Final, while AEK Athens is the other member of the big three. These clubs also have basketball teams; Panathinaikos and Olympiacos are among the top powers in European basketball, having won the Euroleague six times and three respectively, whilst AEK Athens was the first Greek team to win a European trophy in any team sport.
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Other notable clubs within Athens are Athinaikos, Panionios, Atromitos, Apollon, Panellinios, Ethnikos Piraeus, Maroussi BCE and Peristeri B.C.. Athenian clubs have also had domestic and international success in other sports.
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The Athens area encompasses a variety of terrain, notably hills and mountains rising around the city, and the capital is the only major city in Europe to be bisected by a mountain range. Four mountain ranges extend into city boundaries and thousands of miles of trails criss-cross the city and neighbouring areas, providing exercise and wilderness access on foot and bike.
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Beyond Athens and across the prefecture of Attica, outdoor activities include skiing, rock climbing, hang gliding and windsurfing. Numerous outdoor clubs serve these sports, including the Athens Chapter of the Sierra Club, which leads over 4,000 outings annually in the area.
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Beside the above clubs, inside the boundaries of Athens Municipality there are some more clubs with presence in national divisions or notable action for short periods. Some of them are PAO Rouf (Rouf) with earlier presence in Gamma Ethniki, Petralona F.C.(el) (Petralona), football club founded in 1963, with earlier presence in Beta Ethniki, Attikos F.C.(el) (Kolonos), football club founded in 1919 with short presence in Gamma Ethniki, Athinais Kypselis [es] (Kypseli), football club founded in 1938 with short presence in Gamma Ethniki, Gyziakos (Gyzi), basketball club founded in 1937 with short presence in Beta Ethniki basketball and Aetos B.C.(el) (Agios Panteleimonas), basketball club founded in 1992 with earlier presence in A2 Ethniki Basketball. Another important Athenian sport club is the Athens Tennis Club founded in 1895 with important offer for the Greek tennis.[156]
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The revival of the modern Olympic Games was brought forth in 1896, by Frenchman Pierre de Coubertin. Thanks to his efforts, Athens was awarded the first modern Olympic Games. In 1896, the city had a population of 123,000[93] and the event helped boost the city's international profile. Of the venues used for these Olympics, the Kallimarmaro Stadium, and Zappeion were most crucial. The Kallimarmaro is a replica of the ancient Athenian stadiums, and the only major stadium (in its capacity of 60,000) to be made entirely of white marble from Mount Penteli, the same material used for construction of the Parthenon.
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Fencing before the king of Greece at the 1896 Summer Olympics.
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The Panathenaic Stadium of Athens (Kallimarmaron) dates back to the 4th century BC and has hosted the first modern Olympic Games in 1896.
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The 1906 Summer Olympics, or the 1906 Intercalated games, were held in Athens. The intercalated competitions were intermediate games to the internationally organized Olympics, and were meant to be organized in Greece every four years, between the main Olympics. This idea later lost support from the IOC and these games were discontinued.
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Athens was awarded the 2004 Summer Olympics on 5 September 1997 in Lausanne, Switzerland, after having lost a previous bid to host the 1996 Summer Olympics, to Atlanta, United States.[17] It was to be the second time Athens would host the games, following the inaugural event of 1896. After an unsuccessful bid in 1990, the 1997 bid was radically improved, including an appeal to Greece's Olympic history. In the last round of voting, Athens defeated Rome with 66 votes to 41.[17] Prior to this round, the cities of Buenos Aires, Stockholm and Cape Town had been eliminated from competition, having received fewer votes.[17]
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During the first three years of preparations, the International Olympic Committee had expressed concern over the speed of construction progress for some of the new Olympic venues. In 2000 the Organising Committee's president was replaced by Gianna Angelopoulos-Daskalaki, who was the president of the original Bidding Committee in 1997. From that point forward, preparations continued at a highly accelerated, almost frenzied pace.
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Although the heavy cost was criticized, estimated at $1.5 billion, Athens was transformed into a more functional city that enjoys modern technology both in transportation and in modern urban development.[157] Some of the finest sporting venues in the world were created in the city, all of which were fully ready for the games. The games welcomed over 10,000 athletes from all 202 countries.[157]
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The 2004 Games were judged a success, as both security and organization worked well, and only a few visitors reported minor problems mainly concerning accommodation issues. The 2004 Olympic Games were described as Unforgettable, dream Games, by IOC President Jacques Rogge for their return to the birthplace of the Olympics, and for meeting the challenges of holding the Olympic Games.[157] The only observable problem was a somewhat sparse attendance of some early events. Eventually, however, a total of more than 3.5 million tickets were sold, which was higher than any other Olympics with the exception of Sydney (more than 5 million tickets were sold there in 2000).[158]
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In 2008 it was reported that most of the Olympic venues had fallen into disrepair: according to those reports, 21 of the 22 facilities built for the games had either been left abandoned or are in a state of dereliction, with several squatter camps having sprung up around certain facilities, and a number of venues afflicted by vandalism, graffiti or strewn with rubbish.[159][160] These claims, however, are disputed and likely to be inaccurate, as most of the facilities used for the Athens Olympics are either in use or in the process of being converted for post-Olympics use. The Greek Government has created a corporation, Olympic Properties SA, which is overseeing the post-Olympics management, development and conversion of these facilities, some of which will be sold off (or have already been sold off) to the private sector,[161] while other facilities are still in use just as during the Olympics, or have been converted for commercial use or modified for other sports.[162] Concerts and theatrical shows, such as those by the troupe Cirque du Soleil, have recently been held in the complex.[153]
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1 |
+
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In biology, an organism (from Greek: ὀργανισμός, organismos) is any individual entity that embodies the properties of life. It is a synonym for "life form".
|
4 |
+
|
5 |
+
Organisms are classified by taxonomy into groups such as multicellular animals, plants, and fungi; or unicellular microorganisms such as protists, bacteria, and archaea.[1] All types of organisms are capable of reproduction, growth and development, maintenance, and some degree of response to stimuli. Humans, squids, mushrooms, and vascular plants are examples of multicellular organisms that differentiate specialized tissues and organs during development.
|
6 |
+
|
7 |
+
An organism may be either a prokaryote or a eukaryote. Prokaryotes are represented by two separate domains – bacteria and archaea. Eukaryotic organisms are characterized by the presence of a membrane-bound cell nucleus and contain additional membrane-bound compartments called organelles (such as mitochondria in animals and plants and plastids in plants and algae, all generally considered to be derived from endosymbiotic bacteria).[2] Fungi, animals and plants are examples of kingdoms of organisms within the eukaryotes.
|
8 |
+
|
9 |
+
Estimates on the number of Earth's current species range from 2 million to 1 trillion,[3] of which over 1.7 million have been documented.[4] More than 99% of all species, amounting to over five billion species,[5] that ever lived are estimated to be extinct.[6][7]
|
10 |
+
|
11 |
+
In 2016, a set of 355 genes from the last universal common ancestor (LUCA) of all organisms was identified.[8][9]
|
12 |
+
|
13 |
+
The term "organism" (from Greek ὀργανισμός, organismos, from ὄργανον, organon, i.e. "instrument, implement, tool, organ of sense or apprehension")[10][11] first appeared in the English language in 1703 and took on its current definition by 1834 (Oxford English Dictionary). It is directly related to the term "organization". There is a long tradition of defining organisms as self-organizing beings, going back at least to Immanuel Kant's 1790 Critique of Judgment.[12]
|
14 |
+
|
15 |
+
An organism may be defined as an assembly of molecules functioning as a more or less stable whole that exhibits the properties of life. Dictionary definitions can be broad, using phrases such as "any living structure, such as a plant, animal, fungus or bacterium, capable of growth and reproduction".[13] Many definitions exclude viruses and possible man-made non-organic life forms, as viruses are dependent on the biochemical machinery of a host cell for reproduction.[14] A superorganism is an organism consisting of many individuals working together as a single functional or social unit.[15]
|
16 |
+
|
17 |
+
There has been controversy about the best way to define the organism[16][17][18][19][20][21][22][23][24][25] and indeed about whether or not such a definition is necessary.[26][27] Several contributions[28] are responses to the suggestion that the category of "organism" may well not be adequate in biology.[29][page needed]
|
18 |
+
|
19 |
+
Viruses are not typically considered to be organisms because they are incapable of autonomous reproduction, growth or metabolism. Although some organisms are also incapable of independent survival and live as obligatory intracellular parasites, they are capable of independent metabolism and procreation. Although viruses have a few enzymes and molecules characteristic of living organisms, they have no metabolism of their own; they cannot synthesize and organize the organic compounds from which they are formed. Naturally, this rules out autonomous reproduction: they can only be passively replicated by the machinery of the host cell. In this sense, they are similar to inanimate matter.
|
20 |
+
|
21 |
+
While viruses sustain no independent metabolism and thus are usually not classified as organisms, they do have their own genes, and they do evolve by mechanisms similar to the evolutionary mechanisms of organisms. Thus, an argument that viruses should be classed as living organisms is their ability to undergo evolution and replicate through self-assembly. However, some scientists argue that viruses neither evolve nor self-reproduce. Instead, viruses are evolved by their host cells, meaning that there was co-evolution of viruses and host cells. If host cells did not exist, viral evolution would be impossible. This is not true for cells. If viruses did not exist, the direction of cellular evolution could be different, but cells would nevertheless be able to evolve. As for the reproduction, viruses totally rely on hosts' machinery to replicate.[30] The discovery of viruses with genes coding for energy metabolism and protein synthesis fuelled the debate about whether viruses are living organisms. The presence of these genes suggested that viruses were once able to metabolize. However, it was found later that the genes coding for energy and protein metabolism have a cellular origin. Most likely, these genes were acquired through horizontal gene transfer from viral hosts.[30]
|
22 |
+
|
23 |
+
Organisms are complex chemical systems, organized in ways that promote reproduction and some measure of sustainability or survival. The same laws that govern non-living chemistry govern the chemical processes of life. It is generally the phenomena of entire organisms that determine their fitness to an environment and therefore the survivability of their DNA-based genes.
|
24 |
+
|
25 |
+
Organisms clearly owe their origin, metabolism, and many other internal functions to chemical phenomena, especially the chemistry of large organic molecules. Organisms are complex systems of chemical compounds that, through interaction and environment, play a wide variety of roles.
|
26 |
+
|
27 |
+
Organisms are semi-closed chemical systems. Although they are individual units of life (as the definition requires), they are not closed to the environment around them. To operate they constantly take in and release energy. Autotrophs produce usable energy (in the form of organic compounds) using light from the sun or inorganic compounds while heterotrophs take in organic compounds from the environment.
|
28 |
+
|
29 |
+
The primary chemical element in these compounds is carbon. The chemical properties of this element such as its great affinity for bonding with other small atoms, including other carbon atoms, and its small size making it capable of forming multiple bonds, make it ideal as the basis of organic life. It is able to form small three-atom compounds (such as carbon dioxide), as well as large chains of many thousands of atoms that can store data (nucleic acids), hold cells together, and transmit information (protein).
|
30 |
+
|
31 |
+
Compounds that make up organisms may be divided into macromolecules and other, smaller molecules. The four groups of macromolecule are nucleic acids, proteins, carbohydrates and lipids. Nucleic acids (specifically deoxyribonucleic acid, or DNA) store genetic data as a sequence of nucleotides. The particular sequence of the four different types of nucleotides (adenine, cytosine, guanine, and thymine) dictate many characteristics that constitute the organism. The sequence is divided up into codons, each of which is a particular sequence of three nucleotides and corresponds to a particular amino acid. Thus a sequence of DNA codes for a particular protein that, due to the chemical properties of the amino acids it is made from, folds in a particular manner and so performs a particular function.
|
32 |
+
|
33 |
+
These protein functions have been recognized:
|
34 |
+
|
35 |
+
A bilayer of phospholipids makes up the membrane of cells that constitutes a barrier, containing everything within the cell and preventing compounds from freely passing into, and out of, the cell. Due to the selective permeability of the phospholipid membrane, only specific compounds can pass through it.
|
36 |
+
|
37 |
+
All organisms consist of structural units called cells; some contain a single cell (unicellular) and others contain many units (multicellular). Multicellular organisms are able to specialize cells to perform specific functions. A group of such cells is a tissue, and in animals these occur as four basic types, namely epithelium, nervous tissue, muscle tissue, and connective tissue. Several types of tissue work together in the form of an organ to produce a particular function (such as the pumping of the blood by the heart, or as a barrier to the environment as the skin). This pattern continues to a higher level with several organs functioning as an organ system such as the reproductive system, and digestive system. Many multicellular organisms consist of several organ systems, which coordinate to allow for life.
|
38 |
+
|
39 |
+
The cell theory, first developed in 1839 by Schleiden and Schwann, states that all organisms are composed of one or more cells; all cells come from preexisting cells; and cells contain the hereditary information necessary for regulating cell functions and for transmitting information to the next generation of cells.
|
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+
|
41 |
+
There are two types of cells, eukaryotic and prokaryotic. Prokaryotic cells are usually singletons, while eukaryotic cells are usually found in multicellular organisms. Prokaryotic cells lack a nuclear membrane so DNA is unbound within the cell; eukaryotic cells have nuclear membranes.
|
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+
|
43 |
+
All cells, whether prokaryotic or eukaryotic, have a membrane, which envelops the cell, separates its interior from its environment, regulates what moves in and out, and maintains the electric potential of the cell. Inside the membrane, a salty cytoplasm takes up most of the cell volume. All cells possess DNA, the hereditary material of genes, and RNA, containing the information necessary to build various proteins such as enzymes, the cell's primary machinery. There are also other kinds of biomolecules in cells.
|
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+
|
45 |
+
All cells share several similar characteristics of:[31]
|
46 |
+
|
47 |
+
The last universal common ancestor (LUCA) is the most recent organism from which all organisms now living on Earth descend.[32] Thus it is the most recent common ancestor of all current life on Earth. The LUCA is estimated to have lived some 3.5 to 3.8 billion years ago (sometime in the Paleoarchean era).[33][34] The earliest evidence for life on Earth is graphite found to be biogenic in 3.7 billion-year-old metasedimentary rocks discovered in Western Greenland[35] and microbial mat fossils found in 3.48 billion-year-old sandstone discovered in Western Australia.[36][37] Although more than 99 percent of all species that ever lived on the planet are estimated to be extinct,[6][7] there are currently 2 million to 1 trillion species of life on Earth.[3]
|
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+
|
49 |
+
Information about the early development of life includes input from many different fields, including geology and planetary science. These sciences provide information about the history of the Earth and the changes produced by life. However, a great deal of information about the early Earth has been destroyed by geological processes over the course of time.
|
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+
|
51 |
+
All organisms are descended from a common ancestor or ancestral gene pool. Evidence for common descent may be found in traits shared between all living organisms. In Darwin's day, the evidence of shared traits was based solely on visible observation of morphologic similarities, such as the fact that all birds have wings, even those that do not fly.
|
52 |
+
|
53 |
+
There is strong evidence from genetics that all organisms have a common ancestor. For example, every living cell makes use of nucleic acids as its genetic material, and uses the same twenty amino acids as the building blocks for proteins. All organisms use the same genetic code (with some extremely rare and minor deviations) to translate nucleic acid sequences into proteins. The universality of these traits strongly suggests common ancestry, because the selection of many of these traits seems arbitrary. Horizontal gene transfer makes it more difficult to study the last universal ancestor.[38] However, the universal use of the same genetic code, same nucleotides, and same amino acids makes the existence of such an ancestor overwhelmingly likely.[39]
|
54 |
+
|
55 |
+
Chlorobacteria (accepted name = Chloroflexi)
|
56 |
+
|
57 |
+
Hadobacteria (= Deinococcus-Thermus group)
|
58 |
+
|
59 |
+
Cyanobacteria
|
60 |
+
|
61 |
+
Spirochaetae
|
62 |
+
|
63 |
+
Fibrobacteres
|
64 |
+
|
65 |
+
Chlorobi
|
66 |
+
|
67 |
+
Bacteroidetes
|
68 |
+
|
69 |
+
Planctomycetes
|
70 |
+
|
71 |
+
Chlamydiae
|
72 |
+
|
73 |
+
Lentisphaerae
|
74 |
+
|
75 |
+
Verrucomicrobia
|
76 |
+
|
77 |
+
Deferribacteres
|
78 |
+
|
79 |
+
Acidobacteria
|
80 |
+
|
81 |
+
Deltaproteobacteria
|
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+
|
83 |
+
Epsilonproteobacteria
|
84 |
+
|
85 |
+
Alphaproteobacteria
|
86 |
+
|
87 |
+
Betaproteobacteria
|
88 |
+
|
89 |
+
Gammaproteobacteria
|
90 |
+
|
91 |
+
Thermotogae
|
92 |
+
|
93 |
+
Fusobacteria
|
94 |
+
|
95 |
+
Negativicutes
|
96 |
+
|
97 |
+
Endobacteria (=Firmicutes, Mollicutes)
|
98 |
+
|
99 |
+
Actinobacteria
|
100 |
+
|
101 |
+
Archaea
|
102 |
+
|
103 |
+
Eukarya
|
104 |
+
|
105 |
+
The most commonly accepted location of the root of the tree of life is between a monophyletic domain Bacteria and a clade formed by Archaea and Eukaryota of what is referred to as the "traditional tree of life" based on several molecular studies.[40][41][42][43][44][45] A very small minority of studies have concluded differently, namely that the root is in the domain Bacteria, either in the phylum Firmicutes[46] or that the phylum Chloroflexi is basal to a clade with Archaea and Eukaryotes and the rest of Bacteria as proposed by Thomas Cavalier-Smith.[47]
|
106 |
+
|
107 |
+
Research published in 2016, by William F. Martin, by genetically analyzing 6.1 million protein-coding genes from sequenced prokaryotic genomes of various phylogenetic trees, identified 355 protein clusters from amongst 286,514 protein clusters that were probably common to the LUCA. The results "depict LUCA as anaerobic, CO2-fixing, H2-dependent with a Wood–Ljungdahl pathway (the reductive acetyl-coenzyme A pathway), N2-fixing and thermophilic. LUCA's biochemistry was replete with FeS clusters and radical reaction mechanisms. Its cofactors reveal dependence upon transition metals, flavins, S-adenosyl methionine, coenzyme A, ferredoxin, molybdopterin, corrins and selenium. Its genetic code required nucleoside modifications and S-adenosylmethionine-dependent methylations." The results depict methanogenic clostria as a basal clade in the 355 lineages examined, and suggest that the LUCA inhabited an anaerobic hydrothermal vent setting in a geochemically active environment rich in H2, CO2, and iron.[8] However, the identification of these genes as being present in LUCA was criticized, suggesting that many of the proteins assumed to be present in LUCA represent later horizontal gene transfers between archaea and bacteria.[48]
|
108 |
+
|
109 |
+
Sexual reproduction is widespread among current eukaryotes, and was likely present in the last common ancestor.[49] This is suggested by the finding of a core set of genes for meiosis in the descendants of lineages that diverged early from the eukaryotic evolutionary tree.[50] and Malik et al.[51] It is further supported by evidence that eukaryotes previously regarded as "ancient asexuals", such as Amoeba, were likely sexual in the past, and that most present day asexual amoeboid lineages likely arose recently and independently.[52]
|
110 |
+
|
111 |
+
In prokaryotes, natural bacterial transformation involves the transfer of DNA from one bacterium to another and integration of the donor DNA into the recipient chromosome by recombination. Natural bacterial transformation is considered to be a primitive sexual process and occurs in both bacteria and archaea, although it has been studied mainly in bacteria. Transformation is clearly a bacterial adaptation and not an accidental occurrence, because it depends on numerous gene products that specifically interact with each other to enter a state of natural competence to perform this complex process.[53] Transformation is a common mode of DNA transfer among prokaryotes.[54]
|
112 |
+
|
113 |
+
The ancestry of living organisms has traditionally been reconstructed from morphology, but is increasingly supplemented with phylogenetics – the reconstruction of phylogenies by the comparison of genetic (DNA) sequence.
|
114 |
+
|
115 |
+
Sequence comparisons suggest recent horizontal transfer of many genes among diverse species including across the boundaries of phylogenetic "domains". Thus determining the phylogenetic history of a species can not be done conclusively by determining evolutionary trees for single genes.[55]
|
116 |
+
|
117 |
+
Biologist Peter Gogarten suggests "the original metaphor of a tree no longer fits the data from recent genome research", therefore "biologists (should) use the metaphor of a mosaic to describe the different histories combined in individual genomes and use (the) metaphor of a net to visualize the rich exchange and cooperative effects of HGT among microbes."[56]
|
118 |
+
|
119 |
+
Modern biotechnology is challenging traditional concepts of organism and species. Cloning is the process of creating a new multicellular organism, genetically identical to another, with the potential of creating entirely new species of organisms. Cloning is the subject of much ethical debate.
|
120 |
+
|
121 |
+
In 2008, the J. Craig Venter Institute assembled a synthetic bacterial genome, Mycoplasma genitalium, by using recombination in yeast of 25 overlapping DNA fragments in a single step. The use of yeast recombination greatly simplifies the assembly of large DNA molecules from both synthetic and natural fragments.[57] Other companies, such as Synthetic Genomics, have already been formed to take advantage of the many commercial uses of custom designed genomes.
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|
1 |
+
|
2 |
+
|
3 |
+
In biology, an organism (from Greek: ὀργανισμός, organismos) is any individual entity that embodies the properties of life. It is a synonym for "life form".
|
4 |
+
|
5 |
+
Organisms are classified by taxonomy into groups such as multicellular animals, plants, and fungi; or unicellular microorganisms such as protists, bacteria, and archaea.[1] All types of organisms are capable of reproduction, growth and development, maintenance, and some degree of response to stimuli. Humans, squids, mushrooms, and vascular plants are examples of multicellular organisms that differentiate specialized tissues and organs during development.
|
6 |
+
|
7 |
+
An organism may be either a prokaryote or a eukaryote. Prokaryotes are represented by two separate domains – bacteria and archaea. Eukaryotic organisms are characterized by the presence of a membrane-bound cell nucleus and contain additional membrane-bound compartments called organelles (such as mitochondria in animals and plants and plastids in plants and algae, all generally considered to be derived from endosymbiotic bacteria).[2] Fungi, animals and plants are examples of kingdoms of organisms within the eukaryotes.
|
8 |
+
|
9 |
+
Estimates on the number of Earth's current species range from 2 million to 1 trillion,[3] of which over 1.7 million have been documented.[4] More than 99% of all species, amounting to over five billion species,[5] that ever lived are estimated to be extinct.[6][7]
|
10 |
+
|
11 |
+
In 2016, a set of 355 genes from the last universal common ancestor (LUCA) of all organisms was identified.[8][9]
|
12 |
+
|
13 |
+
The term "organism" (from Greek ὀργανισμός, organismos, from ὄργανον, organon, i.e. "instrument, implement, tool, organ of sense or apprehension")[10][11] first appeared in the English language in 1703 and took on its current definition by 1834 (Oxford English Dictionary). It is directly related to the term "organization". There is a long tradition of defining organisms as self-organizing beings, going back at least to Immanuel Kant's 1790 Critique of Judgment.[12]
|
14 |
+
|
15 |
+
An organism may be defined as an assembly of molecules functioning as a more or less stable whole that exhibits the properties of life. Dictionary definitions can be broad, using phrases such as "any living structure, such as a plant, animal, fungus or bacterium, capable of growth and reproduction".[13] Many definitions exclude viruses and possible man-made non-organic life forms, as viruses are dependent on the biochemical machinery of a host cell for reproduction.[14] A superorganism is an organism consisting of many individuals working together as a single functional or social unit.[15]
|
16 |
+
|
17 |
+
There has been controversy about the best way to define the organism[16][17][18][19][20][21][22][23][24][25] and indeed about whether or not such a definition is necessary.[26][27] Several contributions[28] are responses to the suggestion that the category of "organism" may well not be adequate in biology.[29][page needed]
|
18 |
+
|
19 |
+
Viruses are not typically considered to be organisms because they are incapable of autonomous reproduction, growth or metabolism. Although some organisms are also incapable of independent survival and live as obligatory intracellular parasites, they are capable of independent metabolism and procreation. Although viruses have a few enzymes and molecules characteristic of living organisms, they have no metabolism of their own; they cannot synthesize and organize the organic compounds from which they are formed. Naturally, this rules out autonomous reproduction: they can only be passively replicated by the machinery of the host cell. In this sense, they are similar to inanimate matter.
|
20 |
+
|
21 |
+
While viruses sustain no independent metabolism and thus are usually not classified as organisms, they do have their own genes, and they do evolve by mechanisms similar to the evolutionary mechanisms of organisms. Thus, an argument that viruses should be classed as living organisms is their ability to undergo evolution and replicate through self-assembly. However, some scientists argue that viruses neither evolve nor self-reproduce. Instead, viruses are evolved by their host cells, meaning that there was co-evolution of viruses and host cells. If host cells did not exist, viral evolution would be impossible. This is not true for cells. If viruses did not exist, the direction of cellular evolution could be different, but cells would nevertheless be able to evolve. As for the reproduction, viruses totally rely on hosts' machinery to replicate.[30] The discovery of viruses with genes coding for energy metabolism and protein synthesis fuelled the debate about whether viruses are living organisms. The presence of these genes suggested that viruses were once able to metabolize. However, it was found later that the genes coding for energy and protein metabolism have a cellular origin. Most likely, these genes were acquired through horizontal gene transfer from viral hosts.[30]
|
22 |
+
|
23 |
+
Organisms are complex chemical systems, organized in ways that promote reproduction and some measure of sustainability or survival. The same laws that govern non-living chemistry govern the chemical processes of life. It is generally the phenomena of entire organisms that determine their fitness to an environment and therefore the survivability of their DNA-based genes.
|
24 |
+
|
25 |
+
Organisms clearly owe their origin, metabolism, and many other internal functions to chemical phenomena, especially the chemistry of large organic molecules. Organisms are complex systems of chemical compounds that, through interaction and environment, play a wide variety of roles.
|
26 |
+
|
27 |
+
Organisms are semi-closed chemical systems. Although they are individual units of life (as the definition requires), they are not closed to the environment around them. To operate they constantly take in and release energy. Autotrophs produce usable energy (in the form of organic compounds) using light from the sun or inorganic compounds while heterotrophs take in organic compounds from the environment.
|
28 |
+
|
29 |
+
The primary chemical element in these compounds is carbon. The chemical properties of this element such as its great affinity for bonding with other small atoms, including other carbon atoms, and its small size making it capable of forming multiple bonds, make it ideal as the basis of organic life. It is able to form small three-atom compounds (such as carbon dioxide), as well as large chains of many thousands of atoms that can store data (nucleic acids), hold cells together, and transmit information (protein).
|
30 |
+
|
31 |
+
Compounds that make up organisms may be divided into macromolecules and other, smaller molecules. The four groups of macromolecule are nucleic acids, proteins, carbohydrates and lipids. Nucleic acids (specifically deoxyribonucleic acid, or DNA) store genetic data as a sequence of nucleotides. The particular sequence of the four different types of nucleotides (adenine, cytosine, guanine, and thymine) dictate many characteristics that constitute the organism. The sequence is divided up into codons, each of which is a particular sequence of three nucleotides and corresponds to a particular amino acid. Thus a sequence of DNA codes for a particular protein that, due to the chemical properties of the amino acids it is made from, folds in a particular manner and so performs a particular function.
|
32 |
+
|
33 |
+
These protein functions have been recognized:
|
34 |
+
|
35 |
+
A bilayer of phospholipids makes up the membrane of cells that constitutes a barrier, containing everything within the cell and preventing compounds from freely passing into, and out of, the cell. Due to the selective permeability of the phospholipid membrane, only specific compounds can pass through it.
|
36 |
+
|
37 |
+
All organisms consist of structural units called cells; some contain a single cell (unicellular) and others contain many units (multicellular). Multicellular organisms are able to specialize cells to perform specific functions. A group of such cells is a tissue, and in animals these occur as four basic types, namely epithelium, nervous tissue, muscle tissue, and connective tissue. Several types of tissue work together in the form of an organ to produce a particular function (such as the pumping of the blood by the heart, or as a barrier to the environment as the skin). This pattern continues to a higher level with several organs functioning as an organ system such as the reproductive system, and digestive system. Many multicellular organisms consist of several organ systems, which coordinate to allow for life.
|
38 |
+
|
39 |
+
The cell theory, first developed in 1839 by Schleiden and Schwann, states that all organisms are composed of one or more cells; all cells come from preexisting cells; and cells contain the hereditary information necessary for regulating cell functions and for transmitting information to the next generation of cells.
|
40 |
+
|
41 |
+
There are two types of cells, eukaryotic and prokaryotic. Prokaryotic cells are usually singletons, while eukaryotic cells are usually found in multicellular organisms. Prokaryotic cells lack a nuclear membrane so DNA is unbound within the cell; eukaryotic cells have nuclear membranes.
|
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|
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All cells, whether prokaryotic or eukaryotic, have a membrane, which envelops the cell, separates its interior from its environment, regulates what moves in and out, and maintains the electric potential of the cell. Inside the membrane, a salty cytoplasm takes up most of the cell volume. All cells possess DNA, the hereditary material of genes, and RNA, containing the information necessary to build various proteins such as enzymes, the cell's primary machinery. There are also other kinds of biomolecules in cells.
|
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|
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All cells share several similar characteristics of:[31]
|
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|
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The last universal common ancestor (LUCA) is the most recent organism from which all organisms now living on Earth descend.[32] Thus it is the most recent common ancestor of all current life on Earth. The LUCA is estimated to have lived some 3.5 to 3.8 billion years ago (sometime in the Paleoarchean era).[33][34] The earliest evidence for life on Earth is graphite found to be biogenic in 3.7 billion-year-old metasedimentary rocks discovered in Western Greenland[35] and microbial mat fossils found in 3.48 billion-year-old sandstone discovered in Western Australia.[36][37] Although more than 99 percent of all species that ever lived on the planet are estimated to be extinct,[6][7] there are currently 2 million to 1 trillion species of life on Earth.[3]
|
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|
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Information about the early development of life includes input from many different fields, including geology and planetary science. These sciences provide information about the history of the Earth and the changes produced by life. However, a great deal of information about the early Earth has been destroyed by geological processes over the course of time.
|
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|
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All organisms are descended from a common ancestor or ancestral gene pool. Evidence for common descent may be found in traits shared between all living organisms. In Darwin's day, the evidence of shared traits was based solely on visible observation of morphologic similarities, such as the fact that all birds have wings, even those that do not fly.
|
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|
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There is strong evidence from genetics that all organisms have a common ancestor. For example, every living cell makes use of nucleic acids as its genetic material, and uses the same twenty amino acids as the building blocks for proteins. All organisms use the same genetic code (with some extremely rare and minor deviations) to translate nucleic acid sequences into proteins. The universality of these traits strongly suggests common ancestry, because the selection of many of these traits seems arbitrary. Horizontal gene transfer makes it more difficult to study the last universal ancestor.[38] However, the universal use of the same genetic code, same nucleotides, and same amino acids makes the existence of such an ancestor overwhelmingly likely.[39]
|
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|
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Chlorobacteria (accepted name = Chloroflexi)
|
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|
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Hadobacteria (= Deinococcus-Thermus group)
|
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|
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Cyanobacteria
|
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Spirochaetae
|
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Fibrobacteres
|
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Chlorobi
|
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|
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Bacteroidetes
|
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|
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Planctomycetes
|
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|
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Chlamydiae
|
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Lentisphaerae
|
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Verrucomicrobia
|
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Deferribacteres
|
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Acidobacteria
|
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Deltaproteobacteria
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Epsilonproteobacteria
|
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|
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Alphaproteobacteria
|
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|
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Betaproteobacteria
|
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|
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Gammaproteobacteria
|
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|
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Thermotogae
|
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|
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Fusobacteria
|
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|
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Negativicutes
|
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|
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Endobacteria (=Firmicutes, Mollicutes)
|
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|
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Actinobacteria
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Archaea
|
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Eukarya
|
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The most commonly accepted location of the root of the tree of life is between a monophyletic domain Bacteria and a clade formed by Archaea and Eukaryota of what is referred to as the "traditional tree of life" based on several molecular studies.[40][41][42][43][44][45] A very small minority of studies have concluded differently, namely that the root is in the domain Bacteria, either in the phylum Firmicutes[46] or that the phylum Chloroflexi is basal to a clade with Archaea and Eukaryotes and the rest of Bacteria as proposed by Thomas Cavalier-Smith.[47]
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Research published in 2016, by William F. Martin, by genetically analyzing 6.1 million protein-coding genes from sequenced prokaryotic genomes of various phylogenetic trees, identified 355 protein clusters from amongst 286,514 protein clusters that were probably common to the LUCA. The results "depict LUCA as anaerobic, CO2-fixing, H2-dependent with a Wood–Ljungdahl pathway (the reductive acetyl-coenzyme A pathway), N2-fixing and thermophilic. LUCA's biochemistry was replete with FeS clusters and radical reaction mechanisms. Its cofactors reveal dependence upon transition metals, flavins, S-adenosyl methionine, coenzyme A, ferredoxin, molybdopterin, corrins and selenium. Its genetic code required nucleoside modifications and S-adenosylmethionine-dependent methylations." The results depict methanogenic clostria as a basal clade in the 355 lineages examined, and suggest that the LUCA inhabited an anaerobic hydrothermal vent setting in a geochemically active environment rich in H2, CO2, and iron.[8] However, the identification of these genes as being present in LUCA was criticized, suggesting that many of the proteins assumed to be present in LUCA represent later horizontal gene transfers between archaea and bacteria.[48]
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Sexual reproduction is widespread among current eukaryotes, and was likely present in the last common ancestor.[49] This is suggested by the finding of a core set of genes for meiosis in the descendants of lineages that diverged early from the eukaryotic evolutionary tree.[50] and Malik et al.[51] It is further supported by evidence that eukaryotes previously regarded as "ancient asexuals", such as Amoeba, were likely sexual in the past, and that most present day asexual amoeboid lineages likely arose recently and independently.[52]
|
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|
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In prokaryotes, natural bacterial transformation involves the transfer of DNA from one bacterium to another and integration of the donor DNA into the recipient chromosome by recombination. Natural bacterial transformation is considered to be a primitive sexual process and occurs in both bacteria and archaea, although it has been studied mainly in bacteria. Transformation is clearly a bacterial adaptation and not an accidental occurrence, because it depends on numerous gene products that specifically interact with each other to enter a state of natural competence to perform this complex process.[53] Transformation is a common mode of DNA transfer among prokaryotes.[54]
|
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|
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The ancestry of living organisms has traditionally been reconstructed from morphology, but is increasingly supplemented with phylogenetics – the reconstruction of phylogenies by the comparison of genetic (DNA) sequence.
|
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Sequence comparisons suggest recent horizontal transfer of many genes among diverse species including across the boundaries of phylogenetic "domains". Thus determining the phylogenetic history of a species can not be done conclusively by determining evolutionary trees for single genes.[55]
|
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|
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Biologist Peter Gogarten suggests "the original metaphor of a tree no longer fits the data from recent genome research", therefore "biologists (should) use the metaphor of a mosaic to describe the different histories combined in individual genomes and use (the) metaphor of a net to visualize the rich exchange and cooperative effects of HGT among microbes."[56]
|
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|
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Modern biotechnology is challenging traditional concepts of organism and species. Cloning is the process of creating a new multicellular organism, genetically identical to another, with the potential of creating entirely new species of organisms. Cloning is the subject of much ethical debate.
|
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|
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In 2008, the J. Craig Venter Institute assembled a synthetic bacterial genome, Mycoplasma genitalium, by using recombination in yeast of 25 overlapping DNA fragments in a single step. The use of yeast recombination greatly simplifies the assembly of large DNA molecules from both synthetic and natural fragments.[57] Other companies, such as Synthetic Genomics, have already been formed to take advantage of the many commercial uses of custom designed genomes.
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In music, the organ is a keyboard instrument of one or more pipe divisions or other means for producing tones, each played with its own keyboard, played either with the hands on a keyboard or with the feet using pedals.
|
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+
|
5 |
+
The organ is a relatively old musical instrument,[1] dating from the time of Ctesibius of Alexandria (285–222 BC), who invented the water organ. It was played throughout the Ancient Greek and Ancient Roman world, particularly during races and games.[2] During the early medieval period it spread from the Byzantine Empire, where it continued to be used in secular (non-religious) and imperial court music, to Western Europe, where it gradually assumed a prominent place in the liturgy of the Catholic Church.[2] Subsequently, it re-emerged as a secular and recital instrument in the Classical music tradition.
|
6 |
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|
7 |
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Pipe organs use air moving through pipes to produce sounds. Since the 16th century, pipe organs have used various materials for pipes, which can vary widely in timbre and volume. Increasingly hybrid organs are appearing in which pipes are augmented with electric additions. Great economies of space and cost are possible especially when the lowest (and largest) of the pipes can be replaced.
|
8 |
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|
9 |
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Non-piped organs include:
|
10 |
+
|
11 |
+
Mechanical organs include the barrel organ, water organ, and Orchestrion. These are controlled by mechanical means such as pinned barrels or book music. Little barrel organs dispense with the hands of an organist and bigger organs are powered in most cases by an organ grinder or today by other means such as an electric motor.
|
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|
13 |
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The pipe organ is the largest musical instrument. These instruments vary greatly in size, ranging from a cubic meter to a height reaching five floors,[4] and are built in churches, synagogues, concert halls, and homes. Small organs are called "positive" (easily placed in different locations) or "portative" (small enough to carry while playing).
|
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|
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The pipes are divided into ranks and controlled by the use of hand stops and combination pistons. Although the keyboard is not expressive as on a piano and does not affect dynamics (it is binary; pressing a key only turns the sound on or off), some divisions may be enclosed in a swell box, allowing the dynamics to be controlled by shutters. Some organs are totally enclosed, meaning that all the divisions can be controlled by one set of shutters. Some special registers with free reed pipes are expressive.
|
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|
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It has existed in its current form since the 14th century, though similar designs were common in the Eastern Mediterranean from the early Byzantine period (from the 4th century AD) and precursors, such as the hydraulic organ, have been found dating to the late Hellenistic period (1st century BC). Along with the clock, it was considered one of the most complex human-made mechanical creations before the Industrial Revolution. Pipe organs range in size from a single short keyboard to huge instruments with over 10,000 pipes. A large modern organ typically has three or four keyboards (manuals) with five octaves (61 notes) each, and a two-and-a-half octave (32-note) pedal board.
|
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|
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Wolfgang Amadeus Mozart called the organ the "King of instruments".[5] Some of the biggest instruments have 64-foot pipes (a foot here means "sonic-foot", a measure quite close to the English measurement unit)[citation needed] and it sounds to an 8 Hz frequency fundamental tone. Perhaps the most distinctive feature is the ability to range from the slightest sound to the most powerful, plein-jeu impressive sonic discharge, which can be sustained in time indefinitely by the organist. For instance, the Wanamaker organ, located in Philadelphia, USA, has sonic resources comparable with three simultaneous symphony orchestras. Another interesting feature lies in its intrinsic "polyphony" approach: each set of pipes can be played simultaneously with others, and the sounds mixed and interspersed in the environment, not in the instrument itself.
|
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|
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+
Most organs in Europe, the Americas, and Australasia can be found in Christian churches.
|
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|
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The introduction of church organs is traditionally attributed to Pope Vitalian in the 7th century.[citation needed] Due to its simultaneous ability to provide a musical foundation below the vocal register, support in the vocal register, and increased brightness above the vocal register, the organ is ideally suited to accompany human voices, whether a congregation, a choir, or a cantor or soloist.
|
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|
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Most services also include solo organ repertoire for independent performance rather than by way of accompaniment, often as a prelude at the beginning the service and a postlude at the conclusion of the service.
|
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|
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Today this organ may be a pipe organ (see above), a digital or electronic organ that generates the sound with digital signal processing (DSP) chips, or a combination of pipes and electronics. It may be called a church organ or classical organ to differentiate it from the theatre organ, which is a different style of instrument. However, as classical organ repertoire was developed for the pipe organ and in turn influenced its development, the line between a church and a concert organ became harder to draw.
|
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In the late 19th century and early 20th century, symphonic organs flourished in secular venues in the United States and the United Kingdom, designed to replace symphony orchestras by playing transcriptions of orchestral pieces. Symphonic and orchestral organs largely fell out of favor as the orgelbewegung (organ reform movement) took hold in the middle of the 20th century, and organ builders began to look to historical models for inspiration in constructing new instruments. Today, modern builders construct organs in a variety of styles for both secular and sacred applications.
|
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The theatre organ or cinema organ was designed to accompany silent movies. Like a symphonic organ, it is made to replace an orchestra. However, it includes many more gadgets, such as mechanical percussion accessories and other imitative sounds useful in creating movie sound accompaniments such as auto horns, doorbells, and bird whistles. It typically features the Tibia pipe family as its foundation stops and the regular use of a tremulant possessing a depth greater than that on a classical organ.
|
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Theatre organs tend not to take nearly as much space as standard organs, relying on extension (sometimes called unification) and higher wind pressures to produce a greater variety of tone and larger volume of sound from fewer pipes. Unification gives a smaller instrument the capability of a much larger one, and works well for monophonic styles of playing (chordal, or chords with solo voice). The sound is, however, thicker and more homogeneous than a classically designed organ.
|
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|
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+
In the USA the American Theater Organ Society (ATOS) has been instrumental in programs to preserve examples of such instruments.
|
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|
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A chamber organ is a small pipe organ, often with only one manual, and sometimes without separate pedal pipes that is placed in a small room, that this diminutive organ can fill with sound. It is often confined to chamber organ repertoire, as often the organs have too few voice capabilities to rival the grand pipe organs in the performance of the classics. The sound and touch are unique to the instrument, sounding nothing like a large organ with few stops drawn out, but rather much more intimate. They are usually tracker instruments, although the modern builders are often building electropneumatic chamber organs.
|
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|
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Pre-Beethoven keyboard music may usually be as easily played on a chamber organ as on a piano or harpsichord, and a chamber organ is sometimes preferable to a harpsichord for continuo playing as it is more suitable for producing a sustained tone.
|
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|
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The pump organ, reed organ or harmonium, was the other main type of organ before the development of the electronic organ. It generated its sounds using reeds similar to those of an accordion. Smaller, cheaper and more portable than the corresponding pipe instrument, these were widely used in smaller churches and in private homes, but their volume and tonal range was extremely limited. They were generally limited to one or two manuals; they seldom had a pedalboard.
|
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|
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The chord organ was invented by Laurens Hammond in 1950.[6] It provided chord buttons for the left hand, similar to an accordion. Other reed organ manufacturers have also produced chord organs, most notably Magnus from 1958 to the late 1970s.[7]
|
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+
|
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+
Since the 1930s, pipeless electric instruments have been available to produce similar sounds and perform similar roles to pipe organs. Many of these have been bought both by houses of worship and other potential pipe organ customers, and also by many musicians both professional and amateur for whom a pipe organ would not be a possibility. Far smaller and cheaper to buy than a corresponding pipe instrument, and in many cases portable, they have taken organ music into private homes and into dance bands and other new environments, and have almost completely replaced the reed organ.
|
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|
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+
The Hammond organ was the first successful electric organ, released in the 1930s. It used mechanical, rotating tonewheels to produce the sound waveforms. Its system of drawbars allowed for setting volumes for specific sounds, and it provided vibrato-like effects. The drawbars allow the player to choose volume levels. By emphasizing certain harmonics from the overtone series, desired sounds (such as 'brass' or 'string') can be imitated. Generally, the older Hammond drawbar organs had only preamplifiers and were connected to an external, amplified speaker. The Leslie speaker, which rotates to create a distinctive tremolo, became the most popular.
|
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|
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+
Though originally produced to replace organs in the church, the Hammond organ, especially the model B-3, became popular in jazz, particularly soul jazz, and in gospel music. Since these were the roots of rock and roll, the Hammond organ became a part of the rock and roll sound. It was widely used in rock and popular music during the 1960s and 1970s by bands like Emerson, Lake and Palmer, Procol Harum, Santana and Deep Purple. Its popularity resurged in pop music around 2000, in part due to the availability of clonewheel organs that were light enough for one person to carry.
|
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|
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+
In contrast to Hammond's electro-mechanical design, Allen Organ Company introduced the first totally electronic organ in 1938, based on the stable oscillator designed and patented by the Company's founder, Jerome Markowitz.[8] Allen continued to advance analog tone generation through the 1960s with additional patents.[9] In 1971, in collaboration with North American Rockwell,[10] Allen introduced the world's first commercially available digital musical instrument. The first Allen Digital Organ is now in the Smithsonian Institution.[11]
|
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+
|
53 |
+
Frequency divider organs used oscillators instead of mechanical parts to make sound. These were even cheaper and more portable than the Hammond. They featured an ability to bend pitches.
|
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+
|
55 |
+
In the 1940s until the 1970s, small organs were sold that simplified traditional organ stops. These instruments can be considered the predecessor to modern portable keyboards, as they included one-touch chords, rhythm and accompaniment devices, and other electronically assisted gadgets. Lowrey was the leading manufacturer of this type of organs in the smaller (spinet) instruments.
|
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|
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+
In the '60s and '70s, a type of simple, portable electronic organ called the combo organ was popular, especially with pop, Ska (in the late 1970s and early 1980s) and rock bands, and was a signature sound in the pop music of the period, such as The Doors and Iron Butterfly. The most popular combo organs were manufactured by Farfisa and Vox.
|
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|
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+
Conn-Selmer and Rodgers, dominant in the market for larger instruments, also made electronic organs that used separate oscillators for each note rather than frequency dividers, giving them a richer sound, closer to a pipe organ, due to the slight imperfections in tuning.
|
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+
|
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+
Hybrids, starting in the early 20th century,[12] incorporate a few ranks of pipes to produce some sounds, and use electronic circuits or digital samples for other sounds and to resolve borrowing collisions. Major manufacturers include Allen, Walker, Compton, Wicks, Marshall & Ogletree, Phoenix, Makin Organs, Wyvern Organs and Rodgers.
|
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+
|
63 |
+
The development of the integrated circuit enabled another revolution in electronic keyboard instruments.
|
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+
Digital organs sold since the 1970s utilize additive synthesis, then sampling technology (1980s) and physical modelling synthesis (1990s) are also utilized to produce the sound.
|
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+
|
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Virtual pipe organs use MIDI to access samples of real pipe organs stored on a computer, as opposed to digital organs that use DSP and processor hardware inside a console to produce the sounds or deliver the sound samples. Touch screen monitors allows the user to control the virtual organ console; a traditional console and its physical stop and coupler controls is not required. In such a basic form, a virtual organ can be obtained at a much lower cost than other digital classical organs.
|
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|
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The wind can also be created by using pressurized steam instead of air. The steam organ, or calliope, was invented in the United States in the 19th century. Calliopes usually have very loud and clean sound. Calliopes are used as outdoors instruments, and many have been built on wheeled platforms.
|
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|
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The organ has had an important place in classical music, particularly since the 16th century. Spain's Antonio de Cabezón, the Netherlands' Jan Pieterszoon Sweelinck, and Italy's Girolamo Frescobaldi were three of the most important organist-composers before 1650. Influenced in part by Sweelinck and Frescobaldi, the North German school rose from the mid-17th century onwards to great prominence, with leading members of this school having included Buxtehude, Franz Tunder, Georg Böhm, Georg Philipp Telemann, and above all Johann Sebastian Bach, whose contributions to organ music continue to reign supreme.
|
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|
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During this time, the French Classical school also flourished. François Couperin, Nicolas Lebègue, André Raison, and Nicolas de Grigny were French organist-composers of the period. Bach knew Grigny's organ output well, and admired it. In England, Handel was famous for his organ-playing no less than for his composing; several of his organ concertos, intended for his own use, are still frequently performed.
|
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|
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After Bach's death in 1750, the organ's prominence gradually shrank, as the instrument itself increasingly lost ground to the piano. Nevertheless, Felix Mendelssohn, César Franck, and the less famous A.P.F. Boëly (all of whom were themselves expert organists) led, independently of one another, a resurgence of valuable organ writing during the 19th century. This resurgence, much of it informed by Bach's example, achieved particularly impressive things in France (even though Franck himself was of Belgian birth). Major names in French Romantic organ composition are Charles-Marie Widor, Louis Vierne, Alexandre Guilmant, Charles Tournemire, and Eugène Gigout. Of these, Vierne and Tournemire were Franck pupils.
|
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|
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In Germany, Max Reger (late 19th century) owes much to the harmonic daring of Liszt (himself an organ composer) and of Wagner. Paul Hindemith produced three organ sonatas and several works combining organ with chamber groups. Sigfrid Karg-Elert specialized in smaller organ pieces, mostly chorale-preludes.
|
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|
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Among French organist-composers, Marcel Dupré, Maurice Duruflé, Olivier Messiaen and Jean Langlais made significant contributions to the 20th-century organ repertoire. Organ was also used a lot for improvisation, with organists such as Charles Tournemire, Marcel Dupré, Pierre Cochereau, Pierre Pincemaille and Thierry Escaich.
|
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|
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Some composers incorporated the instrument in symphonic works for its dramatic effect, notably Mahler, Holst, Elgar, Scriabin, Respighi, and Richard Strauss. Saint-Saëns's Organ Symphony employs the organ more as an equitable orchestral instrument than for purely dramatic effect. Poulenc wrote the sole organ concerto since Handel's to have achieved mainstream popularity.
|
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|
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Because the organ has both manuals and pedals, organ music has come to be notated on three staves. The music played on the manuals is laid out like music for other keyboard instruments on the top two staves, and the music for the pedals is notated on the third stave or sometimes, to save space, added to the bottom of the second stave as was the early practice. To aid the eye in reading three staves at once, the bar lines are broken between the lowest two staves; the brace surrounds only the upper two staves. Because music racks are often built quite low to preserve sightlines over the console, organ music is usually published in oblong or landscape format.
|
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|
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Electronic organs and electromechanical organs such as the Hammond organ have an established role in a number of popular-music genres, such as blues, jazz, gospel, and 1960s and 1970s rock music. Electronic and electromechanical organs were originally designed as lower-cost substitutes for pipe organs. Despite this intended role as a sacred music instrument, electronic and electromechanical organs' distinctive tone-often modified with electronic effects such as vibrato, rotating Leslie speakers, and overdrive-became an important part of the sound of popular music.
|
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|
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The electric organ, especially the Hammond B-3, has occupied a significant role in jazz ever since Jimmy Smith made it popular in the 1950s. It can function as a replacement for both piano and bass in the standard jazz combo. The Hammond organ is the centrepiece of the organ trio, a small ensemble which typically includes an organist (playing melodies, chords and basslines), a drummer and a third instrumentalist (either jazz guitar or saxophone). In the 2000s, many performers use electronic or digital organs, called clonewheel organs, as they are much lighter and easier to transport than the heavy, bulky B-3.
|
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|
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+
Performers of 20th century popular organ music include William Rowland who composed "Piano Rags"; George Wright (1920–1998) and Virgil Fox (1912–1980), who bridged both the classical and religious areas of music.
|
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|
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Church-style pipe organs are sometimes used in rock music. Examples include Tangerine Dream, Rick Wakeman (with Yes and solo), Keith Emerson (with The Nice and Emerson, Lake and Palmer), George Duke (with Frank Zappa), Dennis DeYoung (with Styx), Arcade Fire, Muse, Roger Hodgson (formerly of Supertramp), Natalie Merchant (with 10,000 Maniacs), Billy Preston and Iron Butterfly.
|
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|
92 |
+
Artists using the Hammond organ include Bob Dylan, Counting Crows, Pink Floyd, Hootie & the Blowfish, Sheryl Crow, Sly Stone and Deep Purple.
|
93 |
+
|
94 |
+
From their creation on radio in the 1930s to the times of television in the early 1970s soap operas incorporated organ music in the background of scenes and in their opening and closing theme music. In the early 1970s the organ was phased out in favour of more dramatic, full-blown orchestras, which in turn were replaced with more modern pop-style compositions.
|
95 |
+
|
96 |
+
In the United States and Canada, organ music is commonly associated with several sports, most notably baseball, basketball, and ice hockey.
|
97 |
+
|
98 |
+
The baseball organ has been referred to as "an accessory to the overall auditory experience of the ballpark."[citation needed] The first team to introduce an organ was the Chicago Cubs, who put an organ in Wrigley Field as an experiment in 1941 for two games. Ebbets Field, home of the Brooklyn Dodgers, hired baseball's first full-time organist, Gladys Goodding. Over the years, many ballparks caught on to the trend, and many organists became well-known and associated with their parks or signature tunes.
|
99 |
+
|
100 |
+
List of NHL Organists
|
101 |
+
|
102 |
+
Panpipes
|
103 |
+
|
104 |
+
Hydraulis
|
105 |
+
|
106 |
+
Portative
|
107 |
+
|
108 |
+
Positive
|
109 |
+
|
110 |
+
Regal
|
111 |
+
|
112 |
+
(after the 16th century)[13]
|
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|
1 |
+
|
2 |
+
|
3 |
+
|
4 |
+
|
5 |
+
Gold is a chemical element with the symbol Au (from Latin: aurum) and atomic number 79, making it one of the higher atomic number elements that occur naturally. In a pure form, it is a bright, slightly reddish yellow, dense, soft, malleable, and ductile metal. Chemically, gold is a transition metal and a group 11 element. It is one of the least reactive chemical elements and is solid under standard conditions. Gold often occurs in free elemental (native) form, as nuggets or grains, in rocks, in veins, and in alluvial deposits. It occurs in a solid solution series with the native element silver (as electrum) and also naturally alloyed with copper and palladium. Less commonly, it occurs in minerals as gold compounds, often with tellurium (gold tellurides).
|
6 |
+
|
7 |
+
Gold is resistant to most acids, though it does dissolve in aqua regia (a mixture of nitric acid and hydrochloric acid), which forms a soluble tetrachloroaurate anion. Gold is insoluble in nitric acid, which dissolves silver and base metals, a property that has long been used to refine gold and to confirm the presence of gold in metallic objects, giving rise to the term acid test. Gold also dissolves in alkaline solutions of cyanide, which are used in mining and electroplating. Gold dissolves in mercury, forming amalgam alloys, but this is not a chemical reaction.
|
8 |
+
|
9 |
+
A relatively rare element,[6][7] gold is a precious metal that has been used as a neutron reflector in nuclear weapons (w71), and for coinage, jewelry, and other arts throughout recorded history. In the past, a gold standard was often implemented as a monetary policy, but gold coins ceased to be minted as a circulating currency in the 1930s, and the world gold standard was abandoned for a fiat currency system after 1971.
|
10 |
+
|
11 |
+
A total of 190,040 tonnes of gold exists above ground, as of 2017[update].[8] This is equal to a cube with each side measuring roughly 21.3 metres. The world consumption of new gold produced is about 50% in jewelry, 40% in investments, and 10% in industry.[9] Gold's high malleability, ductility, resistance to corrosion and most other chemical reactions, and conductivity of electricity have led to its continued use in corrosion resistant electrical connectors in all types of computerized devices (its chief industrial use). Gold is also used in infrared shielding, colored-glass production, gold leafing, and tooth restoration. Certain gold salts are still used as anti-inflammatories in medicine. As of 2017[update], the world's largest gold producer by far was China with 440 tonnes per year.[10]
|
12 |
+
|
13 |
+
Gold is the most malleable of all metals. It can be drawn into a wire of single-atom width, and then stretched considerably before it breaks.[12] Such nanowires distort via formation, reorientation and migration of dislocations and crystal twins without noticeable hardening.[13] A single gram of gold can be beaten into a sheet of 1 square meter, and an avoirdupois ounce into 300 square feet. Gold leaf can be beaten thin enough to become semi-transparent. The transmitted light appears greenish blue, because gold strongly reflects yellow and red.[14] Such semi-transparent sheets also strongly reflect infrared light, making them useful as infrared (radiant heat) shields in visors of heat-resistant suits, and in sun-visors for spacesuits.[15] Gold is a good conductor of heat and electricity.
|
14 |
+
|
15 |
+
Gold has a density of 19.3 g/cm3, almost identical to that of tungsten at 19.25 g/cm3; as such, tungsten has been used in counterfeiting of gold bars, such as by plating a tungsten bar with gold,[16][17][18][19] or taking an existing gold bar, drilling holes, and replacing the removed gold with tungsten rods.[20] By comparison, the density of lead is 11.34 g/cm3, and that of the densest element, osmium, is 22.588±0.015 g/cm3.[21]
|
16 |
+
|
17 |
+
Whereas most metals are gray or silvery white, gold is slightly reddish-yellow.[22] This color is determined by the frequency of plasma oscillations among the metal's valence electrons, in the ultraviolet range for most metals but in the visible range for gold due to relativistic effects affecting the orbitals around gold atoms.[23][24] Similar effects impart a golden hue to metallic caesium.
|
18 |
+
|
19 |
+
Common colored gold alloys include the distinctive eighteen-karat rose gold created by the addition of copper. Alloys containing palladium or nickel are also important in commercial jewelry as these produce white gold alloys. Fourteen-karat gold-copper alloy is nearly identical in color to certain bronze alloys, and both may be used to produce police and other badges. Fourteen- and eighteen-karat gold alloys with silver alone appear greenish-yellow and are referred to as green gold. Blue gold can be made by alloying with iron, and purple gold can be made by alloying with aluminium. Less commonly, addition of manganese, indium, and other elements can produce more unusual colors of gold for various applications.[25]
|
20 |
+
|
21 |
+
Colloidal gold, used by electron-microscopists, is red if the particles are small; larger particles of colloidal gold are blue.[26]
|
22 |
+
|
23 |
+
Gold has only one stable isotope, 197Au, which is also its only naturally occurring isotope, so gold is both a mononuclidic and monoisotopic element. Thirty-six radioisotopes have been synthesized, ranging in atomic mass from 169 to 205. The most stable of these is 195Au with a half-life of 186.1 days. The least stable is 171Au, which decays by proton emission with a half-life of 30 µs. Most of gold's radioisotopes with atomic masses below 197 decay by some combination of proton emission, α decay, and β+ decay. The exceptions are 195Au, which decays by electron capture, and 196Au, which decays most often by electron capture (93%) with a minor β− decay path (7%).[27] All of gold's radioisotopes with atomic masses above 197 decay by β− decay.[28]
|
24 |
+
|
25 |
+
At least 32 nuclear isomers have also been characterized, ranging in atomic mass from 170 to 200. Within that range, only 178Au, 180Au, 181Au, 182Au, and 188Au do not have isomers. Gold's most stable isomer is 198m2Au with a half-life of 2.27 days. Gold's least stable isomer is 177m2Au with a half-life of only 7 ns. 184m1Au has three decay paths: β+ decay, isomeric transition, and alpha decay. No other isomer or isotope of gold has three decay paths.[28]
|
26 |
+
|
27 |
+
The production of gold from a more common element, such as lead, has long been a subject of human inquiry, and the ancient and medieval discipline of alchemy often focused on it; however, the transmutation of the chemical elements did not become possible until the understanding of nuclear physics in the 20th century. The first synthesis of gold was conducted by Japanese physicist Hantaro Nagaoka, who synthesized gold from mercury in 1924 by neutron bombardment.[29] An American team, working without knowledge of Nagaoka's prior study, conducted the same experiment in 1941, achieving the same result and showing that the isotopes of gold produced by it were all radioactive.[30]
|
28 |
+
|
29 |
+
Gold can currently be manufactured in a nuclear reactor by irradiation either of platinum or mercury.
|
30 |
+
|
31 |
+
Only the mercury isotope 196Hg, which occurs with a frequency of 0.15% in natural mercury, can be converted to gold by neutron capture, and following electron capture-decay into 197Au with slow neutrons. Other mercury isotopes are converted when irradiated with slow neutrons into one another, or formed mercury isotopes which beta decay into thallium.
|
32 |
+
|
33 |
+
Using fast neutrons, the mercury isotope 198Hg, which composes 9.97% of natural mercury, can be converted by splitting off a neutron and becoming 197Hg, which then disintegrates to stable gold. This reaction, however, possesses a smaller activation cross-section and is feasible only with un-moderated reactors.
|
34 |
+
|
35 |
+
It is also possible to eject several neutrons with very high energy into the other mercury isotopes in order to form 197Hg. However, such high-energy neutrons can be produced only by particle accelerators.[clarification needed]
|
36 |
+
|
37 |
+
Although gold is the most noble of the noble metals,[31][32] it still forms many diverse compounds. The oxidation state of gold in its compounds ranges from −1 to +5, but Au(I) and Au(III) dominate its chemistry. Au(I), referred to as the aurous ion, is the most common oxidation state with soft ligands such as thioethers, thiolates, and tertiary phosphines. Au(I) compounds are typically linear. A good example is Au(CN)2−, which is the soluble form of gold encountered in mining. The binary gold halides, such as AuCl, form zigzag polymeric chains, again featuring linear coordination at Au. Most drugs based on gold are Au(I) derivatives.[33]
|
38 |
+
|
39 |
+
Au(III) (referred to as the auric) is a common oxidation state, and is illustrated by gold(III) chloride, Au2Cl6. The gold atom centers in Au(III) complexes, like other d8 compounds, are typically square planar, with chemical bonds that have both covalent and ionic character.
|
40 |
+
|
41 |
+
Gold does not react with oxygen at any temperature[34] and, up to 100 °C, is resistant to attack from ozone.[35]
|
42 |
+
|
43 |
+
Some free halogens react with gold.[36] Gold is strongly attacked by fluorine at dull-red heat[37] to form gold(III) fluoride. Powdered gold reacts with chlorine at 180 °C to form AuCl3.[38] Gold reacts with bromine at 140 °C to form gold(III) bromide, but reacts only very slowly with iodine to form the monoiodide.
|
44 |
+
|
45 |
+
Gold does not react with sulfur directly,[39] but gold(III) sulfide can be made by passing hydrogen sulfide through a dilute solution of gold(III) chloride or chlorauric acid.
|
46 |
+
|
47 |
+
Gold readily dissolves in mercury at room temperature to form an amalgam, and forms alloys with many other metals at higher temperatures. These alloys can be produced to modify the hardness and other metallurgical properties, to control melting point or to create exotic colors.[25]
|
48 |
+
|
49 |
+
Gold is unaffected by most acids. It does not react with hydrofluoric, hydrochloric, hydrobromic, hydriodic, sulfuric, or nitric acid. It does react with selenic acid, and is dissolved by aqua regia, a 1:3 mixture of nitric acid and hydrochloric acid. Nitric acid oxidizes the metal to +3 ions, but only in minute amounts, typically undetectable in the pure acid because of the chemical equilibrium of the reaction. However, the ions are removed from the equilibrium by hydrochloric acid, forming AuCl4− ions, or chloroauric acid, thereby enabling further oxidation.
|
50 |
+
|
51 |
+
Gold is similarly unaffected by most bases. It does not react with aqueous, solid, or molten sodium or potassium hydroxide. It does however, react with sodium or potassium cyanide under alkaline conditions when oxygen is present to form soluble complexes.[39]
|
52 |
+
|
53 |
+
Common oxidation states of gold include +1 (gold(I) or aurous compounds) and +3 (gold(III) or auric compounds). Gold ions in solution are readily reduced and precipitated as metal by adding any other metal as the reducing agent. The added metal is oxidized and dissolves, allowing the gold to be displaced from solution and be recovered as a solid precipitate.
|
54 |
+
|
55 |
+
Less common oxidation states of gold include −1, +2, and +5.
|
56 |
+
|
57 |
+
The −1 oxidation state occurs in aurides, compounds containing the Au− anion. Caesium auride (CsAu), for example, crystallizes in the caesium chloride motif;[40] rubidium, potassium, and tetramethylammonium aurides are also known.[41] Gold has the highest electron affinity of any metal, at 222.8 kJ/mol, making Au− a stable species.[42]
|
58 |
+
|
59 |
+
Gold(II) compounds are usually diamagnetic with Au–Au bonds such as [Au(CH2)2P(C6H5)2]2Cl2. The evaporation of a solution of Au(OH)3 in concentrated H2SO4 produces red crystals of gold(II) sulfate, Au2(SO4)2. Originally thought to be a mixed-valence compound, it has been shown to contain Au4+2 cations, analogous to the better-known mercury(I) ion, Hg2+2.[43][44] A gold(II) complex, the tetraxenonogold(II) cation, which contains xenon as a ligand, occurs in [AuXe4](Sb2F11)2.[45]
|
60 |
+
|
61 |
+
Gold pentafluoride, along with its derivative anion, AuF−6, and its difluorine complex, gold heptafluoride, is the sole example of gold(V), the highest verified oxidation state.[46]
|
62 |
+
|
63 |
+
Some gold compounds exhibit aurophilic bonding, which describes the tendency of gold ions to interact at distances that are too long to be a conventional Au–Au bond but shorter than van der Waals bonding. The interaction is estimated to be comparable in strength to that of a hydrogen bond.
|
64 |
+
|
65 |
+
Well-defined cluster compounds are numerous.[41] In such cases, gold has a fractional oxidation state. A representative example is the octahedral species {Au(P(C6H5)3)}62+. Gold chalcogenides, such as gold sulfide, feature equal amounts of Au(I) and Au(III).
|
66 |
+
|
67 |
+
Medicinal applications of gold and its complexes have a long history dating back thousands of years.[47] Several gold complexes have been applied to treat rheumatoid arthritis, the most frequently used being aurothiomalate, aurothioglucose, and auranofin. Both gold(I) and gold(III) compounds have been investigated as possible anti-cancer drugs. For gold(III) complexes, reduction to gold(0/I) under physiological conditions has to be considered. Stable complexes can be generated using different types of bi-, tri-, and tetradentate ligand systems, and their efficacy has been demonstrated in vitro and in vivo.[48]
|
68 |
+
|
69 |
+
Gold is thought to have been produced in supernova nucleosynthesis, and from the collision of neutron stars,[49] and to have been present in the dust from which the Solar System formed.[50]
|
70 |
+
|
71 |
+
Traditionally, gold in the universe is thought to have formed by the r-process (rapid neutron capture) in supernova nucleosynthesis,[51] but more recently it has been suggested that gold and other elements heavier than iron may also be produced in quantity by the r-process in the collision of neutron stars.[52] In both cases, satellite spectrometers at first only indirectly detected the resulting gold.[53] However, in August 2017, the spectrascopic signatures of heavy elements, including gold, were observed by electromagnetic observatories in the GW170817 neutron star merger event, after gravitational wave detectors confirmed the event as a neutron star merger.[54] Current astrophysical models suggest that this single neutron star merger event generated between 3 and 13 Earth masses of gold. This amount, along with estimations of the rate of occurrence of these neutron star merger events, suggests that such mergers may produce enough gold to account for most of the abundance of this element in the universe.[55]
|
72 |
+
|
73 |
+
Because the Earth was molten when it was formed, almost all of the gold present in the early Earth probably sank into the planetary core. Therefore, most of the gold that is in the Earth's crust and mantle has in one model thought to have been delivered to Earth later, by asteroid impacts during the Late Heavy Bombardment, about 4 billion years ago.[56][57]
|
74 |
+
|
75 |
+
Gold which is reachable by humans has, in one case, been associated with a particular asteroid impact. The asteroid that formed Vredefort crater 2.020 billion years ago is often credited with seeding the Witwatersrand basin in South Africa with the richest gold deposits on earth.[58][59][60][61] However, this scenario is now questioned. The gold-bearing Witwatersrand rocks were laid down between 700 and 950 million years before the Vredefort impact.[62][63] These gold-bearing rocks had furthermore been covered by a thick layer of Ventersdorp lavas and the Transvaal Supergroup of rocks before the meteor struck, and thus the gold did not actually arrive in the asteroid/meteorite. What the Vredefort impact achieved, however, was to distort the Witwatersrand basin in such a way that the gold-bearing rocks were brought to the present erosion surface in Johannesburg, on the Witwatersrand, just inside the rim of the original 300 km diameter crater caused by the meteor strike. The discovery of the deposit in 1886 launched the Witwatersrand Gold Rush. Some 22% of all the gold that is ascertained to exist today on Earth has been extracted from these Witwatersrand rocks.[63]
|
76 |
+
|
77 |
+
Notwithstanding the impact above, much of the rest of the gold on Earth is thought to have been incorporated into the planet since its very beginning, as planetesimals formed the planet's mantle, early in Earth's creation. In 2017, an international group of scientists, established that gold "came to the Earth's surface from the deepest regions of our planet",[64] the mantle, evidenced by their findings at Deseado Massif in the Argentinian Patagonia.[65][clarification needed]
|
78 |
+
|
79 |
+
On Earth, gold is found in ores in rock formed from the Precambrian time onward.[66] It most often occurs as a native metal, typically in a metal solid solution with silver (i.e. as a gold silver alloy). Such alloys usually have a silver content of 8–10%. Electrum is elemental gold with more than 20% silver. Electrum's color runs from golden-silvery to silvery, dependent upon the silver content. The more silver, the lower the specific gravity.
|
80 |
+
|
81 |
+
Native gold occurs as very small to microscopic particles embedded in rock, often together with quartz or sulfide minerals such as "fool's gold", which is a pyrite.[67] These are called lode deposits. The metal in a native state is also found in the form of free flakes, grains or larger nuggets[66] that have been eroded from rocks and end up in alluvial deposits called placer deposits. Such free gold is always richer at the surface of gold-bearing veins[clarification needed] owing to the oxidation of accompanying minerals followed by weathering, and washing of the dust into streams and rivers, where it collects and can be welded by water action to form nuggets.
|
82 |
+
|
83 |
+
Gold sometimes occurs combined with tellurium as the minerals calaverite, krennerite, nagyagite, petzite and sylvanite (see telluride minerals), and as the rare bismuthide maldonite (Au2Bi) and antimonide aurostibite (AuSb2). Gold also occurs in rare alloys with copper, lead, and mercury: the minerals auricupride (Cu3Au), novodneprite (AuPb3) and weishanite ((Au, Ag)3Hg2).
|
84 |
+
|
85 |
+
Recent research suggests that microbes can sometimes play an important role in forming gold deposits, transporting and precipitating gold to form grains and nuggets that collect in alluvial deposits.[68]
|
86 |
+
|
87 |
+
Another recent study has claimed water in faults vaporizes during an earthquake, depositing gold. When an earthquake strikes, it moves along a fault. Water often lubricates faults, filling in fractures and jogs. About 6 miles (10 kilometers) below the surface, under incredible temperatures and pressures, the water carries high concentrations of carbon dioxide, silica, and gold. During an earthquake, the fault jog suddenly opens wider. The water inside the void instantly vaporizes, flashing to steam and forcing silica, which forms the mineral quartz, and gold out of the fluids and onto nearby surfaces.[69]
|
88 |
+
|
89 |
+
The world's oceans contain gold. Measured concentrations of gold in the Atlantic and Northeast Pacific are 50–150 femtomol/L or 10–30 parts per quadrillion (about 10–30 g/km3). In general, gold concentrations for south Atlantic and central Pacific samples are the same (~50 femtomol/L) but less certain. Mediterranean deep waters contain slightly higher concentrations of gold (100–150 femtomol/L) attributed to wind-blown dust and/or rivers. At 10 parts per quadrillion the Earth's oceans would hold 15,000 tonnes of gold.[70] These figures are three orders of magnitude less than reported in the literature prior to 1988, indicating contamination problems with the earlier data.
|
90 |
+
|
91 |
+
A number of people have claimed to be able to economically recover gold from sea water, but they were either mistaken or acted in an intentional deception. Prescott Jernegan ran a gold-from-seawater swindle in the United States in the 1890s, as did an English fraudster in the early 1900s.[71] Fritz Haber did research on the extraction of gold from sea water in an effort to help pay Germany's reparations following World War I.[72] Based on the published values of 2 to 64 ppb of gold in seawater a commercially successful extraction seemed possible. After analysis of 4,000 water samples yielding an average of 0.004 ppb it became clear that extraction would not be possible and he stopped the project.[73]
|
92 |
+
|
93 |
+
The earliest recorded metal employed by humans appears to be gold, which can be found free or "native". Small amounts of natural gold have been found in Spanish caves used during the late Paleolithic period, c. 40,000 BC.[75] Gold artifacts made their first appearance at the very beginning of the pre-dynastic period in Egypt, at the end of the fifth millennium BC and the start of the fourth, and smelting was developed during the course of the 4th millennium; gold artifacts appear in the archeology of Lower Mesopotamia during the early 4th millennium.[76] Gold artifacts in the Balkans appear from the 4th millennium BC, such as those found in the Varna Necropolis near Lake Varna in Bulgaria, thought by one source (La Niece 2009) to be the earliest "well-dated" find of gold artifacts.[66] As of 1990, gold artifacts found at the Wadi Qana cave cemetery of the 4th millennium BC in West Bank were the earliest from the Levant.[77] Gold artifacts such as the golden hats and the Nebra disk appeared in Central Europe from the 2nd millennium BC Bronze Age.
|
94 |
+
|
95 |
+
The oldest known map of a gold mine was drawn in the 19th Dynasty of Ancient Egypt (1320–1200 BC), whereas the first written reference to gold was recorded in the 12th Dynasty around 1900 BC.[78] Egyptian hieroglyphs from as early as 2600 BC describe gold, which King Tushratta of the Mitanni claimed was "more plentiful than dirt" in Egypt.[79] Egypt and especially Nubia had the resources to make them major gold-producing areas for much of history. One of the earliest known maps, known as the Turin Papyrus Map, shows the plan of a gold mine in Nubia together with indications of the local geology. The primitive working methods are described by both Strabo and Diodorus Siculus, and included fire-setting. Large mines were also present across the Red Sea in what is now Saudi Arabia.
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Gold is mentioned in the Amarna letters numbered 19[80] and 26[81] from around the 14th century BC.[82][83]
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Gold is mentioned frequently in the Old Testament, starting with Genesis 2:11 (at Havilah), the story of the golden calf, and many parts of the temple including the Menorah and the golden altar. In the New Testament, it is included with the gifts of the magi in the first chapters of Matthew. The Book of Revelation 21:21 describes the city of New Jerusalem as having streets "made of pure gold, clear as crystal". Exploitation of gold in the south-east corner of the Black Sea is said to date from the time of Midas, and this gold was important in the establishment of what is probably the world's earliest coinage in Lydia around 610 BC.[84] The legend of the golden fleece dating from eighth century BCE may refer to the use of fleeces to trap gold dust from placer deposits in the ancient world. From the 6th or 5th century BC, the Chu (state) circulated the Ying Yuan, one kind of square gold coin.
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In Roman metallurgy, new methods for extracting gold on a large scale were developed by introducing hydraulic mining methods, especially in Hispania from 25 BC onwards and in Dacia from 106 AD onwards. One of their largest mines was at Las Medulas in León, where seven long aqueducts enabled them to sluice most of a large alluvial deposit. The mines at Roşia Montană in Transylvania were also very large, and until very recently, still mined by opencast methods. They also exploited smaller deposits in Britain, such as placer and hard-rock deposits at Dolaucothi. The various methods they used are well described by Pliny the Elder in his encyclopedia Naturalis Historia written towards the end of the first century AD.
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During Mansa Musa's (ruler of the Mali Empire from 1312 to 1337) hajj to Mecca in 1324, he passed through Cairo in July 1324, and was reportedly accompanied by a camel train that included thousands of people and nearly a hundred camels where he gave away so much gold that it depressed the price in Egypt for over a decade, causing high inflation.[85] A contemporary Arab historian remarked:
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Gold was at a high price in Egypt until they came in that year. The mithqal did not go below 25 dirhams and was generally above, but from that time its value fell and it cheapened in price and has remained cheap till now. The mithqal does not exceed 22 dirhams or less. This has been the state of affairs for about twelve years until this day by reason of the large amount of gold which they brought into Egypt and spent there [...].
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The European exploration of the Americas was fueled in no small part by reports of the gold ornaments displayed in great profusion by Native American peoples, especially in Mesoamerica, Peru, Ecuador and Colombia. The Aztecs regarded gold as the product of the gods, calling it literally "god excrement" (teocuitlatl in Nahuatl), and after Moctezuma II was killed, most of this gold was shipped to Spain.[88] However, for the indigenous peoples of North America gold was considered useless and they saw much greater value in other minerals which were directly related to their utility, such as obsidian, flint, and slate.[89] El Dorado is applied to a legendary story in which precious stones were found in fabulous abundance along with gold coins. The concept of El Dorado underwent several transformations, and eventually accounts of the previous myth were also combined with those of a legendary lost city. El Dorado, was the term used by the Spanish Empire to describe a mythical tribal chief (zipa) of the Muisca native people in Colombia, who, as an initiation rite, covered himself with gold dust and submerged in Lake Guatavita. The legends surrounding El Dorado changed over time, as it went from being a man, to a city, to a kingdom, and then finally to an empire.
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Gold played a role in western culture, as a cause for desire and of corruption, as told in children's fables such as Rumpelstiltskin—where Rumpelstiltskin turns hay into gold for the peasant's daughter in return for her child when she becomes a princess—and the stealing of the hen that lays golden eggs in Jack and the Beanstalk.
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The top prize at the Olympic Games and many other sports competitions is the gold medal.
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75% of the presently accounted for gold has been extracted since 1910. It has been estimated that the currently known amount of gold internationally would form a single cube 20 m (66 ft) on a side (equivalent to 8,000 m3).[90]
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One main goal of the alchemists was to produce gold from other substances, such as lead — presumably by the interaction with a mythical substance called the philosopher's stone. Although they never succeeded in this attempt, the alchemists did promote an interest in systematically finding out what can be done with substances, and this laid the foundation for today's chemistry. Their symbol for gold was the circle with a point at its center (☉), which was also the astrological symbol and the ancient Chinese character for the Sun.
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The Dome of the Rock is covered with an ultra-thin golden glassier. The Sikh Golden temple, the Harmandir Sahib, is a building covered with gold. Similarly the Wat Phra Kaew emerald Buddhist temple (wat) in Thailand has ornamental gold-leafed statues and roofs. Some European king and queen's crowns were made of gold, and gold was used for the bridal crown since antiquity. An ancient Talmudic text circa 100 AD describes Rachel, wife of Rabbi Akiva, receiving a "Jerusalem of Gold" (diadem). A Greek burial crown made of gold was found in a grave circa 370 BC.
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Minoan jewellery; 2300-2100 BC; various sizes; Metropolitan Museum of Art (New York City)
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Pair of Sumerian earrings with cuneiform inscriptions; 2093–2046 BC; Sulaymaniyah Museum (Sulaymaniyah, Iraq)
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Ancient Egyptian statuette of Amun; 945–715 BC; gold; 17.5 × 4.7 cm; Metropolitan Museum of Art
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Ancient Egyptian signet ring; 664–525 BC; gold; diameter: 3 × 3.4 cm; British Museum (London)
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Ancient Greek stater; 323-315 BC; 18 mm; Metropolitan Museum of Art
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Etruscan funerary wreath; 4th–3rd century BC; length: 33.3 cm; Metropolitan Museum of Art
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Roman aureus of Hadrian; 134–138 AD; 7.4 g; Metropolitan Museum of Art
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Quimbaya lime container; 5th–9th century; gold; height: 23 cm; Metropolitan Museum of Art
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Byzantine scyphate; 1059-1067; diameter: 2.5 cm; Cleveland Museum of Art (Cleveland, Ohio, USA)
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Pre-Columbian pendant with two bat-head worriors who carry spears; 11th–16th century; gold; overall: 7.62 cm; from the Chiriqui Province (Panama); Metropolitan Museum of Art
|
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English Neoclassical box; 1741; overall: 4.4 × 11.6 × 9.2 cm; Metropolitan Museum of Art
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French Rococo glass bottle mounted in gold; circa 1775; overall: 7 x 2.9 cm; Cleveland Museum of Art
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"Gold" is cognate with similar words in many Germanic languages, deriving via Proto-Germanic *gulþą from Proto-Indo-European *ǵʰelh₃- ("to shine, to gleam; to be yellow or green").[91][92]
|
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The symbol Au is from the Latin: aurum, the Latin word for "gold".[93] The Proto-Indo-European ancestor of aurum was *h₂é-h₂us-o-, meaning "glow". This word is derived from the same root (Proto-Indo-European *h₂u̯es- "to dawn") as *h₂éu̯sōs, the ancestor of the Latin word Aurora, "dawn".[94] This etymological relationship is presumably behind the frequent claim in scientific publications that aurum meant "shining dawn".[95]
|
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Outside chemistry, gold is mentioned in a variety of expressions, most often associated with intrinsic worth.[42] Great human achievements are frequently rewarded with gold, in the form of gold medals, gold trophies and other decorations. Winners of athletic events and other graded competitions are usually awarded a gold medal. Many awards such as the Nobel Prize are made from gold as well. Other award statues and prizes are depicted in gold or are gold plated (such as the Academy Awards, the Golden Globe Awards, the Emmy Awards, the Palme d'Or, and the British Academy Film Awards).
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Aristotle in his ethics used gold symbolism when referring to what is now known as the golden mean. Similarly, gold is associated with perfect or divine principles, such as in the case of the golden ratio and the golden rule.
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Gold is further associated with the wisdom of aging and fruition. The fiftieth wedding anniversary is golden. A person's most valued or most successful latter years are sometimes considered "golden years". The height of a civilization is referred to as a golden age.
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In some forms of Christianity and Judaism, gold has been associated both with holiness and evil. In the Book of Exodus, the Golden Calf is a symbol of idolatry, while in the Book of Genesis, Abraham was said to be rich in gold and silver, and Moses was instructed to cover the Mercy Seat of the Ark of the Covenant with pure gold. In Byzantine iconography the halos of Christ, Mary and the Christian saints are often golden.
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In Islam,[96] gold (along with silk)[97][98] is often cited as being forbidden for men to wear.[99] Abu Bakr al-Jazaeri, quoting a hadith, said that "[t]he wearing of silk and gold are forbidden on the males of my nation, and they are lawful to their women".[100] This, however, has not been enforced consistently throughout history, e.g. in the Ottoman Empire.[101] Further, small gold accents on clothing, such as in embroidery, may be permitted.[102]
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According to Christopher Columbus, those who had something of gold were in possession of something of great value on Earth and a substance to even help souls to paradise.[103]
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Wedding rings are typically made of gold. It is long lasting and unaffected by the passage of time and may aid in the ring symbolism of eternal vows before God and the perfection the marriage signifies. In Orthodox Christian wedding ceremonies, the wedded couple is adorned with a golden crown (though some opt for wreaths, instead) during the ceremony, an amalgamation of symbolic rites.
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The World Gold Council states that as of the end of 2017, "there were 187,200 tonnes of stocks in existence above ground". This can be represented by a cube with an edge length of about 21 meters.[104] At $1,349 per troy ounce, 187,200 metric tonnes of gold would have a value of $8.9 trillion. According to the United States Geological Survey in 2016, about 5,726,000,000 troy ounces (178,100 t) of gold has been produced since the beginning of civilization, of which 85% remains in use.[105]
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In 2017, the world's largest gold producer by far was China with 440 tonnes. The second-largest producer, Australia, mined 300 tonnes in the same year, followed by Russia with 255 tonnes.[10]
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Since the 1880s, South Africa has been the source of a large proportion of the world's gold supply, and about 22% of the gold presently accounted is from South Africa. Production in 1970 accounted for 79% of the world supply, about 1,480 tonnes. In 2007 China (with 276 tonnes) overtook South Africa as the world's largest gold producer, the first time since 1905 that South Africa has not been the largest.[106]
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As of 2017[update], China was the world's leading gold-mining country, followed in order by Australia, Russia, the United States, Canada, and Peru. South Africa, which had dominated world gold production for most of the 20th century, had declined to sixth place.[10] Other major producers are the Ghana, Burkina Faso, Mali, Indonesia and Uzbekistan.
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In South America, the controversial project Pascua Lama aims at exploitation of rich fields in the high mountains of Atacama Desert, at the border between Chile and Argentina.
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Today about one-quarter of the world gold output is estimated to originate from artisanal or small scale mining.[107]
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The city of Johannesburg located in South Africa was founded as a result of the Witwatersrand Gold Rush which resulted in the discovery of some of the largest natural gold deposits in recorded history. The gold fields are confined to the northern and north-western edges of the Witwatersrand basin, which is a 5–7 km thick layer of archean rocks located, in most places, deep under the Free State, Gauteng and surrounding provinces.[108] These Witwatersrand rocks are exposed at the surface on the Witwatersrand, in and around Johannesburg, but also in isolated patches to the south-east and south-west of Johannesburg, as well as in an arc around the Vredefort Dome which lies close to the center of the Witwatersrand basin.[62][108] From these surface exposures the basin dips extensively, requiring some of the mining to occur at depths of nearly 4000 m, making them, especially the Savuka and TauTona mines to the south-west of Johannesburg, the deepest mines on earth. The gold is found only in six areas where archean rivers from the north and north-west formed extensive pebbly Braided river deltas before draining into the "Witwatersrand sea" where the rest of the Witwatersrand sediments were deposited.[108]
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The Second Boer War of 1899–1901 between the British Empire and the Afrikaner Boers was at least partly over the rights of miners and possession of the gold wealth in South Africa.
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|
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During the 19th century, gold rushes occurred whenever large gold deposits were discovered. The first documented discovery of gold in the United States was at the Reed Gold Mine near Georgeville, North Carolina in 1803.[109] The first major gold strike in the United States occurred in a small north Georgia town called Dahlonega.[110] Further gold rushes occurred in California, Colorado, the Black Hills, Otago in New Zealand, a number of locations across Australia, Witwatersrand in South Africa, and the Klondike in Canada.
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Gold extraction is most economical in large, easily mined deposits. Ore grades as little as 0.5 parts per million (ppm) can be economical. Typical ore grades in open-pit mines are 1–5 ppm; ore grades in underground or hard rock mines are usually at least 3 ppm. Because ore grades of 30 ppm are usually needed before gold is visible to the naked eye, in most gold mines the gold is invisible.
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|
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The average gold mining and extraction costs were about $317 per troy ounce in 2007, but these can vary widely depending on mining type and ore quality; global mine production amounted to 2,471.1 tonnes.[114]
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After initial production, gold is often subsequently refined industrially by the Wohlwill process which is based on electrolysis or by the Miller process, that is chlorination in the melt. The Wohlwill process results in higher purity, but is more complex and is only applied in small-scale installations.[115][116] Other methods of assaying and purifying smaller amounts of gold include parting and inquartation as well as cupellation, or refining methods based on the dissolution of gold in aqua regia.[117]
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The consumption of gold produced in the world is about 50% in jewelry, 40% in investments, and 10% in industry.[9][118]
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According to World Gold Council, China is the world's largest single consumer of gold in 2013 and toppled India for the first time with Chinese consumption increasing by 32 percent in a year, while that of India only rose by 13 percent and world consumption rose by 21 percent. Unlike India where gold is mainly used for jewelry, China uses gold for manufacturing and retail.[119]
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|
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Gold production is associated with contribution to hazardous pollution.[120][121]
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Low-grade gold ore may contain less than one ppm gold metal; such ore is ground and mixed with sodium cyanide to dissolve the gold. Cyanide is a highly poisonous chemical, which can kill living creatures when exposed in minute quantities. Many cyanide spills[122] from gold mines have occurred in both developed and developing countries which killed aquatic life in long stretches of affected rivers. Environmentalists consider these events major environmental disasters.[123][124] Thirty tons of used ore is dumped as waste for producing one troy ounce of gold.[125] Gold ore dumps are the source of many heavy elements such as cadmium, lead, zinc, copper, arsenic, selenium and mercury. When sulfide-bearing minerals in these ore dumps are exposed to air and water, the sulfide transforms into sulfuric acid which in turn dissolves these heavy metals facilitating their passage into surface water and ground water. This process is called acid mine drainage. These gold ore dumps are long term, highly hazardous wastes second only to nuclear waste dumps.[125]
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It was once common to use mercury to recover gold from ore, but today the use of mercury is largely limited to small-scale individual miners.[126] Minute quantities of mercury compounds can reach water bodies, causing heavy metal contamination. Mercury can then enter into the human food chain in the form of methylmercury. Mercury poisoning in humans causes incurable brain function damage and severe retardation.
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Gold extraction is also a highly energy intensive industry, extracting ore from deep mines and grinding the large quantity of ore for further chemical extraction requires nearly 25 kWh of electricity per gram of gold produced.[127]
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|
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Gold has been widely used throughout the world as money,[128] for efficient indirect exchange (versus barter), and to store wealth in hoards. For exchange purposes, mints produce standardized gold bullion coins, bars and other units of fixed weight and purity.
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|
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The first known coins containing gold were struck in Lydia, Asia Minor, around 600 BC.[84] The talent coin of gold in use during the periods of Grecian history both before and during the time of the life of Homer weighed between 8.42 and 8.75 grams.[129] From an earlier preference in using silver, European economies re-established the minting of gold as coinage during the thirteenth and fourteenth centuries.[130]
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Bills (that mature into gold coin) and gold certificates (convertible into gold coin at the issuing bank) added to the circulating stock of gold standard money in most 19th century industrial economies.
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In preparation for World War I the warring nations moved to fractional gold standards, inflating their currencies to finance the war effort.
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Post-war, the victorious countries, most notably Britain, gradually restored gold-convertibility, but international flows of gold via bills of exchange remained embargoed; international shipments were made exclusively for bilateral trades or to pay war reparations.
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After World War II gold was replaced by a system of nominally convertible currencies related by fixed exchange rates following the Bretton Woods system. Gold standards and the direct convertibility of currencies to gold have been abandoned by world governments, led in 1971 by the United States' refusal to redeem its dollars in gold. Fiat currency now fills most monetary roles. Switzerland was the last country to tie its currency to gold; it backed 40% of its value until the Swiss joined the International Monetary Fund in 1999.[131]
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Central banks continue to keep a portion of their liquid reserves as gold in some form, and metals exchanges such as the London Bullion Market Association still clear transactions denominated in gold, including future delivery contracts.
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Today, gold mining output is declining.[132]
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With the sharp growth of economies in the 20th century, and increasing foreign exchange, the world's gold reserves and their trading market have become a small fraction of all markets and fixed exchange rates of currencies to gold have been replaced by floating prices for gold and gold future contract.
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Though the gold stock grows by only 1 or 2% per year, very little metal is irretrievably consumed. Inventory above ground would satisfy many decades of industrial and even artisan uses at current prices.
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The gold proportion (fineness) of alloys is measured by karat (k). Pure gold (commercially termed fine gold) is designated as 24 karat, abbreviated 24k. English gold coins intended for circulation from 1526 into the 1930s were typically a standard 22k alloy called crown gold,[133] for hardness (American gold coins for circulation after 1837 contain an alloy of 0.900 fine gold, or 21.6 kt).[134]
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Although the prices of some platinum group metals can be much higher, gold has long been considered the most desirable of precious metals, and its value has been used as the standard for many currencies. Gold has been used as a symbol for purity, value, royalty, and particularly roles that combine these properties. Gold as a sign of wealth and prestige was ridiculed by Thomas More in his treatise Utopia. On that imaginary island, gold is so abundant that it is used to make chains for slaves, tableware, and lavatory seats. When ambassadors from other countries arrive, dressed in ostentatious gold jewels and badges, the Utopians mistake them for menial servants, paying homage instead to the most modestly dressed of their party.
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The ISO 4217 currency code of gold is XAU.[135] Many holders of gold store it in form of bullion coins or bars as a hedge against inflation or other economic disruptions, though its efficacy as such has been questioned; historically, it has not proven itself reliable as a hedging instrument.[136] Modern bullion coins for investment or collector purposes do not require good mechanical wear properties; they are typically fine gold at 24k, although the American Gold Eagle and the British gold sovereign continue to be minted in 22k (0.92) metal in historical tradition, and the South African Krugerrand, first released in 1967, is also 22k (0.92).[137]
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The special issue Canadian Gold Maple Leaf coin contains the highest purity gold of any bullion coin, at 99.999% or 0.99999, while the popular issue Canadian Gold Maple Leaf coin has a purity of 99.99%. In 2006, the United States Mint began producing the American Buffalo gold bullion coin with a purity of 99.99%. The Australian Gold Kangaroos were first coined in 1986 as the Australian Gold Nugget but changed the reverse design in 1989. Other modern coins include the Austrian Vienna Philharmonic bullion coin and the Chinese Gold Panda.
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As of September 2017[update], gold is valued at around $42 per gram ($1,300 per troy ounce).
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Like other precious metals, gold is measured by troy weight and by grams. The proportion of gold in the alloy is measured by karat (k), with 24 karat (24k) being pure gold, and lower karat numbers proportionally less. The purity of a gold bar or coin can also be expressed as a decimal figure ranging from 0 to 1, known as the millesimal fineness, such as 0.995 being nearly pure.
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The price of gold is determined through trading in the gold and derivatives markets, but a procedure known as the Gold Fixing in London, originating in September 1919, provides a daily benchmark price to the industry. The afternoon fixing was introduced in 1968 to provide a price when US markets are open.[138]
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Historically gold coinage was widely used as currency; when paper money was introduced, it typically was a receipt redeemable for gold coin or bullion. In a monetary system known as the gold standard, a certain weight of gold was given the name of a unit of currency. For a long period, the United States government set the value of the US dollar so that one troy ounce was equal to $20.67 ($0.665 per gram), but in 1934 the dollar was devalued to $35.00 per troy ounce ($0.889/g). By 1961, it was becoming hard to maintain this price, and a pool of US and European banks agreed to manipulate the market to prevent further currency devaluation against increased gold demand.[139]
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On 17 March 1968, economic circumstances caused the collapse of the gold pool, and a two-tiered pricing scheme was established whereby gold was still used to settle international accounts at the old $35.00 per troy ounce ($1.13/g) but the price of gold on the private market was allowed to fluctuate; this two-tiered pricing system was abandoned in 1975 when the price of gold was left to find its free-market level.[citation needed] Central banks still hold historical gold reserves as a store of value although the level has generally been declining.[citation needed] The largest gold depository in the world is that of the U.S. Federal Reserve Bank in New York, which holds about 3%[140] of the gold known to exist and accounted for today, as does the similarly laden U.S. Bullion Depository at Fort Knox.
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In 2005 the World Gold Council estimated total global gold supply to be 3,859 tonnes and demand to be 3,754 tonnes, giving a surplus of 105 tonnes.[141]
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After 15 August 1971 Nixon shock, the price began to greatly increase,[142] and between 1968 and 2000 the price of gold ranged widely, from a high of $850 per troy ounce ($27.33/g) on 21 January 1980, to a low of $252.90 per troy ounce ($8.13/g) on 21 June 1999 (London Gold Fixing).[143] Prices increased rapidly from 2001, but the 1980 high was not exceeded until 3 January 2008, when a new maximum of $865.35 per troy ounce was set.[144] Another record price was set on 17 March 2008, at $1023.50 per troy ounce ($32.91/g).[144]
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In late 2009, gold markets experienced renewed momentum upwards due to increased demand and a weakening US dollar.[citation needed] On 2 December 2009, gold reached a new high closing at $1,217.23.[145] Gold further rallied hitting new highs in May 2010 after the European Union debt crisis prompted further purchase of gold as a safe asset.[146][147] On 1 March 2011, gold hit a new all-time high of $1432.57, based on investor concerns regarding ongoing unrest in North Africa as well as in the Middle East.[148]
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From April 2001 to August 2011, spot gold prices more than quintupled in value against the US dollar, hitting a new all-time high of $1,913.50 on 23 August 2011,[149] prompting speculation that the long secular bear market had ended and a bull market had returned.[150] However, the price then began a slow decline towards $1200 per troy ounce in late 2014 and 2015.
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Because of the softness of pure (24k) gold, it is usually alloyed with base metals for use in jewelry, altering its hardness and ductility, melting point, color and other properties. Alloys with lower karat rating, typically 22k, 18k, 14k or 10k, contain higher percentages of copper or other base metals or silver or palladium in the alloy.[25] Nickel is toxic, and its release from nickel white gold is controlled by legislation in Europe.[25] Palladium-gold alloys are more expensive than those using nickel.[151] High-karat white gold alloys are more resistant to corrosion than are either pure silver or sterling silver. The Japanese craft of Mokume-gane exploits the color contrasts between laminated colored gold alloys to produce decorative wood-grain effects.
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By 2014, the gold jewelry industry was escalating despite a dip in gold prices. Demand in the first quarter of 2014 pushed turnover to $23.7 billion according to a World Gold Council report.
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Gold solder is used for joining the components of gold jewelry by high-temperature hard soldering or brazing. If the work is to be of hallmarking quality, the gold solder alloy must match the fineness (purity) of the work, and alloy formulas are manufactured to color-match yellow and white gold. Gold solder is usually made in at least three melting-point ranges referred to as Easy, Medium and Hard. By using the hard, high-melting point solder first, followed by solders with progressively lower melting points, goldsmiths can assemble complex items with several separate soldered joints. Gold can also be made into thread and used in embroidery.
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Only 10% of the world consumption of new gold produced goes to industry,[9] but by far the most important industrial use for new gold is in fabrication of corrosion-free electrical connectors in computers and other electrical devices. For example, according to the World Gold Council, a typical cell phone may contain 50 mg of gold, worth about 50 cents. But since nearly one billion cell phones are produced each year, a gold value of 50 cents in each phone adds to $500 million in gold from just this application.[152]
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Though gold is attacked by free chlorine, its good conductivity and general resistance to oxidation and corrosion in other environments (including resistance to non-chlorinated acids) has led to its widespread industrial use in the electronic era as a thin-layer coating on electrical connectors, thereby ensuring good connection. For example, gold is used in the connectors of the more expensive electronics cables, such as audio, video and USB cables. The benefit of using gold over other connector metals such as tin in these applications has been debated; gold connectors are often criticized by audio-visual experts as unnecessary for most consumers and seen as simply a marketing ploy. However, the use of gold in other applications in electronic sliding contacts in highly humid or corrosive atmospheres, and in use for contacts with a very high failure cost (certain computers, communications equipment, spacecraft, jet aircraft engines) remains very common.[153]
|
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+
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+
Besides sliding electrical contacts, gold is also used in electrical contacts because of its resistance to corrosion, electrical conductivity, ductility and lack of toxicity.[154] Switch contacts are generally subjected to more intense corrosion stress than are sliding contacts. Fine gold wires are used to connect semiconductor devices to their packages through a process known as wire bonding.
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The concentration of free electrons in gold metal is 5.91×1022 cm−3.[155] Gold is highly conductive to electricity, and has been used for electrical wiring in some high-energy applications (only silver and copper are more conductive per volume, but gold has the advantage of corrosion resistance). For example, gold electrical wires were used during some of the Manhattan Project's atomic experiments, but large high-current silver wires were used in the calutron isotope separator magnets in the project.
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It is estimated that 16% of the world's presently-accounted-for gold and 22% of the world's silver is contained in electronic technology in Japan.[156]
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Metallic and gold compounds have long been used for medicinal purposes. Gold, usually as the metal, is perhaps the most anciently administered medicine (apparently by shamanic practitioners)[157] and known to Dioscorides.[158][159] In medieval times, gold was often seen as beneficial for the health, in the belief that something so rare and beautiful could not be anything but healthy. Even some modern esotericists and forms of alternative medicine assign metallic gold a healing power.
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In the 19th century gold had a reputation as a "nervine", a therapy for nervous disorders. Depression, epilepsy, migraine, and glandular problems such as amenorrhea and impotence were treated, and most notably alcoholism (Keeley, 1897).[160]
|
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The apparent paradox of the actual toxicology of the substance suggests the possibility of serious gaps in the understanding of the action of gold in physiology.[161] Only salts and radioisotopes of gold are of pharmacological value, since elemental (metallic) gold is inert to all chemicals it encounters inside the body (i.e., ingested gold cannot be attacked by stomach acid). Some gold salts do have anti-inflammatory properties and at present two are still used as pharmaceuticals in the treatment of arthritis and other similar conditions in the US (sodium aurothiomalate and auranofin). These drugs have been explored as a means to help to reduce the pain and swelling of rheumatoid arthritis, and also (historically) against tuberculosis and some parasites.[162]
|
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|
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Gold alloys are used in restorative dentistry, especially in tooth restorations, such as crowns and permanent bridges. The gold alloys' slight malleability facilitates the creation of a superior molar mating surface with other teeth and produces results that are generally more satisfactory than those produced by the creation of porcelain crowns. The use of gold crowns in more prominent teeth such as incisors is favored in some cultures and discouraged in others.
|
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Colloidal gold preparations (suspensions of gold nanoparticles) in water are intensely red-colored, and can be made with tightly controlled particle sizes up to a few tens of nanometers across by reduction of gold chloride with citrate or ascorbate ions. Colloidal gold is used in research applications in medicine, biology and materials science. The technique of immunogold labeling exploits the ability of the gold particles to adsorb protein molecules onto their surfaces. Colloidal gold particles coated with specific antibodies can be used as probes for the presence and position of antigens on the surfaces of cells.[163] In ultrathin sections of tissues viewed by electron microscopy, the immunogold labels appear as extremely dense round spots at the position of the antigen.[164]
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Gold, or alloys of gold and palladium, are applied as conductive coating to biological specimens and other non-conducting materials such as plastics and glass to be viewed in a scanning electron microscope. The coating, which is usually applied by sputtering with an argon plasma, has a triple role in this application. Gold's very high electrical conductivity drains electrical charge to earth, and its very high density provides stopping power for electrons in the electron beam, helping to limit the depth to which the electron beam penetrates the specimen. This improves definition of the position and topography of the specimen surface and increases the spatial resolution of the image. Gold also produces a high output of secondary electrons when irradiated by an electron beam, and these low-energy electrons are the most commonly used signal source used in the scanning electron microscope.[165]
|
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|
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The isotope gold-198 (half-life 2.7 days) is used in nuclear medicine, in some cancer treatments and for treating other diseases.[166][167]
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Pure metallic (elemental) gold is non-toxic and non-irritating when ingested[179] and is sometimes used as a food decoration in the form of gold leaf.[180] Metallic gold is also a component of the alcoholic drinks Goldschläger, Gold Strike, and Goldwasser. Metallic gold is approved as a food additive in the EU (E175 in the Codex Alimentarius). Although the gold ion is toxic, the acceptance of metallic gold as a food additive is due to its relative chemical inertness, and resistance to being corroded or transformed into soluble salts (gold compounds) by any known chemical process which would be encountered in the human body.
|
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|
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Soluble compounds (gold salts) such as gold chloride are toxic to the liver and kidneys. Common cyanide salts of gold such as potassium gold cyanide, used in gold electroplating, are toxic by virtue of both their cyanide and gold content. There are rare cases of lethal gold poisoning from potassium gold cyanide.[181][182] Gold toxicity can be ameliorated with chelation therapy with an agent such as dimercaprol.
|
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|
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+
Gold metal was voted Allergen of the Year in 2001 by the American Contact Dermatitis Society; gold contact allergies affect mostly women.[183] Despite this, gold is a relatively non-potent contact allergen, in comparison with metals like nickel.[184]
|
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|
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+
A sample of the fungus Aspergillus niger was found growing from gold mining solution; and was found to contain cyano metal complexes, such as gold, silver, copper, iron and zinc. The fungus also plays a role in the solubilization of heavy metal sulfides.[185]
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en/4314.html.txt
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Hominina
|
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|
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Australopithecina or Hominina is a subtribe in the tribe Hominini. The members of the subtribe are generally Australopithecus (cladistically including the genera Homo, Paranthropus,[2] and Kenyanthropus), and it typically includes the earlier Ardipithecus, Orrorin, Sahelanthropus, and Graecopithecus. All these related species are now sometimes collectively termed australopithecines or homininians.[3][4] They are the extinct, close relatives of humans and, with the extant genus Homo, comprise the human clade. Members of the human clade, i.e. the Hominini after the split from the chimpanzees, are now called Hominina[5] (see Hominidae; terms "hominids" and hominins).
|
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|
7 |
+
While none of the groups normally directly assigned to this group survived, the australopithecines do not appear to be literally extinct (in the sense of having no living descendants) as the genera Kenyanthropus, Paranthropus and Homo probably emerged as sister of a late Australopithecus species such as A. africanus and/or A. sediba.
|
8 |
+
|
9 |
+
The terms australopithecine, et al., come from a former classification as members of a distinct subfamily, the Australopithecinae.[6] Members of Australopithecus are sometimes referred to as the "gracile australopithecines", while Paranthropus are called the "robust australopithecines".[7][8]
|
10 |
+
|
11 |
+
The australopithecines occurred in the Plio-Pleistocene era and were bipedal, and they were dentally similar to humans, but with a brain size not much larger than that of modern apes, with lesser encephalization than in the genus Homo.[9] Humans (genus Homo) may have descended from australopithecine ancestors and the genera Ardipithecus, Orrorin, Sahelanthropus, and Graecopithecus are the possible ancestors of the australopithecines.[8]
|
12 |
+
|
13 |
+
Phylogeny of subtribe Australopithecina according to Briggs & Crowther 2008, p. 124.
|
14 |
+
|
15 |
+
The post-cranial remains of australopithecines show they were adapted to bipedal locomotion, but did not walk identically to humans. They have a high brachial index (forearm/upper arm ratio) when compared to other hominins, and they exhibit greater sexual dimorphism than members of Homo or Pan but less so than Gorilla or Pongo. It is thought that they averaged heights of 1.2–1.5 metres (3.9–4.9 ft) and weighed between 30 and 55 kilograms (66 and 121 lb). The brain size may have been 350 cc to 600 cc. The postcanines (the teeth behind the canines) were relatively large, and had more enamel compared to contemporary apes and humans, whereas the incisors and canines were relatively small, and there was little difference between the males' and females' canines compared to modern apes.[8]
|
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|
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+
Most scientists maintain that one of the australopithecine species evolved into the genus Homo in Africa around two million years ago. However, there is no consensus on which species:
|
18 |
+
|
19 |
+
"Determining which species of australopithecine (if any) is ancestral to the genus Homo is a question that is a top priority for many paleoanthropologists, but one that will likely elude any conclusive answers for years to come. Nearly every possible species has been suggested as a likely candidate, but none are overwhelmingly convincing. Presently, it appears that A. garhi has the potential to occupy this coveted place in paleoanthropology, but the lack of fossil evidence is a serious problem. Another problem presents itself in the fact that it has been very difficult to assess which hominid [now "hominin"] represents the first member of the genus Homo. Without knowing this, it is not possible to determine which species of australopithecine may have been ancestral to Homo."[8]
|
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+
|
21 |
+
Marc Verhaegen has argued that an australopithecine species could have also been ancestral to the genus Pan (i.e. chimpanzees).[10]
|
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+
|
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+
A minority held viewpoint among palaeoanthropologists is that australopithecines moved outside Africa. A notable proponent of this theory is Jens Lorenz Franzen, formerly Head of Paleoanthropology at the Research Institute Senckenberg. Franzen argues that robust australopithecines had reached not only Indonesia, as Meganthropus, but also China:
|
24 |
+
|
25 |
+
"In this way we arrive at the conclusion that the recognition of australopithecines in Asia would not confuse but could help to clarify the early evolution of hominids ["hominins"] on that continent. This concept would explain the scanty remains from Java and China as relic of an Asian offshoot of an early radiation of Australopithecus, which was followed much later by an [African] immigration of Homo erectus, and finally became extinct after a period of coexistence."[11]
|
26 |
+
|
27 |
+
In 1957, an Early Pleistocene Chinese fossil tooth of unknown province was described as resembling P. robustus. Three fossilized molars from Jianshi, China (Longgudong Cave) were later identified as belonging to an Australopithecus species (Gao, 1975). However further examination questioned this interpretation; Zhang (1984) argued the Jianshi teeth and unidentified tooth belong to H. erectus. Liu et al. (2010) also dispute the Jianshi-australopithecine link and argue the Jianshi molars fall within the range of Homo erectus:
|
28 |
+
|
29 |
+
"No marked difference in dental crown shape is shown between the Jianshi hominin and other Chinese Homo erectus, and there is also no evidence in support of the Jianshi hominin's closeness to Australopithecus."
|
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+
|
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+
But, Wolpoff (1999) notes that in China "persistent claims of australopithecine or australopithecine-like remains continue".
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en/4315.html.txt
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Cervus alces Linnaeus, 1758
|
6 |
+
|
7 |
+
The moose (North America) or elk (Eurasia), Alces alces, is a member of the New World deer subfamily and is the largest and heaviest extant species in the deer family. Most adult male moose have distinctive broad, palmate ("open-hand shaped") antlers; most other members of the deer family have antlers with a dendritic ("twig-like") configuration. Moose typically inhabit boreal forests and temperate broadleaf and mixed forests of the Northern Hemisphere in temperate to subarctic climates. Hunting and other human activities have caused a reduction in the size of the moose's range over time. It has been reintroduced to some of its former habitats. Currently, most moose occur in Canada, Alaska, New England (with Maine having the most of the lower 48 states), Fennoscandia, the Baltic states, and Russia. Its diet consists of both terrestrial and aquatic vegetation. The most common moose predators are the gray wolf along with bears and humans. Unlike most other deer species, moose do not form herds and are solitary animals, aside from calves who remain with their mother until the cow begins estrus (typically at 18 months after birth of the calf), at which point the cow chases away young bulls. Although generally slow-moving and sedentary, moose can become aggressive and move quickly if angered or startled. Their mating season in the autumn features energetic fights between males competing for a female.
|
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+
|
9 |
+
Alces alces is called a "moose" in North American English, but an "elk" in British English;[4] its scientific name comes from its name in Latin. The word "elk" in North American English refers to a completely different species of deer, Cervus canadensis, also called the wapiti. A mature male moose is called a bull, a mature female a cow, and an immature moose of either sex a calf.
|
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+
|
11 |
+
The word "elk" originated in Proto-Germanic, from which Old English evolved and has cognates in other Indo-European languages, e.g. elg in Danish/Norwegian; älg in Swedish; alnis in Latvian; Elch in German; and łoś in Polish (Latin alcē or alcēs and Ancient Greek ἄλκη álkē are probably Germanic loanwords).[5] In the continental-European languages, these forms of the word "elk" almost always refer to Alces alces.
|
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+
|
13 |
+
The word "moose" had first entered English by 1606[6] and is borrowed from the Algonquian languages (compare the Narragansett moos and Eastern Abenaki mos; according to early sources, these were likely derived from moosu, meaning "he strips off"),[7] and possibly involved forms from multiple languages mutually reinforcing one another. The Proto-Algonquian form was *mo·swa.[8]
|
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+
|
15 |
+
The moose became extinct in Britain during the Bronze Age, long before the European arrival in the Americas. The youngest bones were found in Scotland and are roughly 3,900 years old.[9] The word "elk" remained in usage because of its existence in continental Europe; however, without any living animals around to serve as a reference, the meaning became rather vague to most speakers of English, who used "elk" to refer to "large deer" in general.[10] Dictionaries of the 18th century simply described "elk" as a deer that was "as large as a horse".[11]
|
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+
|
17 |
+
Confusingly, the word "elk" is used in North America to refer to a different animal, Cervus canadensis, which is also called by the Algonquian indigenous name, "wapiti". The British began colonizing America in the 17th century, and found two common species of deer for which they had no names. The wapiti appeared very similar to the red deer of Europe (which itself was almost extinct in Southern Britain) although it was much larger and was not red.[10] The moose was a rather strange-looking deer to the colonists, and they often adopted local names for both. In the early days of American colonization, the wapiti was often called a gray moose and the moose was often called a black moose, but early accounts of the animals varied wildly, adding to the confusion.[12]
|
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+
|
19 |
+
The wapiti is superficially very similar to the red deer of central and western Europe, although it is distinctly different behaviorally and genetically. Early European explorers in North America, particularly in Virginia where there were no moose, called the wapiti "elk" because of its size and resemblance to familiar-looking deer like the red deer.[13] The moose resembled the "German elk" (the moose of continental Europe), which was less familiar to the British colonists. For a long time neither species had an official name, but were called a variety of things. Eventually, in North America the wapiti became known as an elk while the moose retained its Anglicized Native-American name.[13] In 1736, Samuel Dale wrote to the Royal Society of Great Britain:
|
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+
|
21 |
+
The common light-grey moose, called by the Indians, Wampoose, and the large or black-moose, which is the beast whose horns I herewith present. As to the grey moose, I take it to be no larger than what Mr. John Clayton, in his account of the Virginia Quadrupeds, calls the Elke ... was in all respects like those of our red-deer or stags, only larger ... The black moose is (by all that have hitherto writ of it) accounted a very large creature. ... The stag, buck, or male of this kind has a palmed horn, not like that of our common or fallow-deer, but the palm is much longer, and more like that of the German elke.[14]
|
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+
|
23 |
+
Moose require habitat with adequate edible plants (e.g., pond grasses, young trees and shrubs), cover from predators, and protection from extremely hot or cold weather. Moose travel among different habitats with the seasons to address these requirements.[15] Moose are cold-adapted mammals with thickened skin, dense, heat-retaining coat, and a low surface:volume ratio, which provides excellent cold tolerance but poor heat tolerance. Moose survive hot weather by accessing shade or cooling wind, or by immersion in cool water. In hot weather, moose are often found wading or swimming in lakes or ponds. When heat-stressed, moose may fail to adequately forage in summer and may not gain adequate body fat to survive the winter. Also, moose cows may not calve without adequate summer weight gain. Moose require access to both young forest for browsing and mature forest for shelter and cover. Forest disturbed by fire and logging promotes the growth of fodder for moose. Moose also require access to mineral licks, safe places for calving and aquatic feeding sites.[15]
|
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+
|
25 |
+
Moose avoid areas with little or no snow as this increases the risk of predation by wolves and avoid areas with deep snow, as this impairs mobility. Thus, moose select habitat on the basis of trade-offs between risk of predation, food availability, and snow depth.[16] With reintroduction of bison into boreal forest, there was some concern that bison would compete with moose for winter habitat, and thereby worsen the population decline of moose. However, this does not appear to be a problem. Moose prefer sub-alpine shrublands in early winter, while bison prefer wet sedge valley meadowlands in early-winter. In late-winter, moose prefer river valleys with deciduous forest cover or alpine terrain above the tree line, while bison preferred wet sedge meadowlands or sunny southern grassy slopes.[17]
|
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+
|
27 |
+
After expanding for most of the 20th century, the moose population of North America has been in steep decline since the 1990s. Populations expanded greatly with improved habitat and protection, but now the moose population is declining rapidly.[18] This decline has been attributed to opening of roads and landscapes into the northern range of moose, allowing deer to become populous in areas where they were not previously common. This encroachment by deer on moose habitat brought moose into contact with previously unfamiliar pathogens, including brainworm and liver fluke, and these parasites are believed to have contributed to the population decline of moose.[19]
|
28 |
+
In North America, the moose range includes almost all of Canada (excluding the arctic and Vancouver Island), most of Alaska, northern New England and upstate New York, the upper Rocky Mountains, northern Minnesota, northern Wisconsin,
|
29 |
+
Michigan's Upper Peninsula, and Isle Royale in Lake Superior. This massive range, containing diverse habitats, contains four of the six North American subspecies. In the West, moose populations extend well north into Canada (British Columbia and Alberta), and more isolated groups have been verified as far south as the mountains of Utah and Colorado and as far west as the Lake Wenatchee area of the Washington Cascades.[20][21] The range includes Wyoming, Montana, Idaho, and smaller areas of Washington and Oregon.[22] Moose have extended their range southwards in the western Rocky Mountains, with initial sightings in Yellowstone National Park in 1868, and then to the northern slope of the Uinta Mountains in Utah in the first half of the twentieth century.[23] This is the southernmost naturally established moose population in the United States.[23] In 1978, a few breeding pairs were reintroduced in western Colorado, and the state's moose population is now more than 1,000.
|
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+
|
31 |
+
In northeastern North America, the Eastern moose's history is very well documented: moose meat was often a staple in the diet of Native Americans going back centuries, with a tribe that occupied present day coastal Rhode Island giving the animal its distinctive name, adopted into American English. The Native Americans often used moose hides for leather and its meat as an ingredient in pemmican, a type of dried jerky used as a source of sustenance in winter or on long journeys.[24] Eastern tribes also valued moose leather as a source for moccasins and other items.
|
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+
|
33 |
+
The historical range of the subspecies extended from well into Quebec, the Maritimes, and Eastern Ontario south to include all of New England finally ending in the very northeastern tip of Pennsylvania in the west, cutting off somewhere near the mouth of the Hudson River in the south . The moose has been extinct in much of the eastern U.S. for as long as 150 years, due to colonial era overhunting and destruction of its habitat: Dutch, French, and British colonial sources all attest to its presence in the mid 17th century from Maine south to areas within a hundred miles of present-day Manhattan. However, by the 1870s, only a handful of moose existed in this entire region in very remote pockets of forest; less than 20% of suitable habitat remained.[25]
|
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+
|
35 |
+
Since the 1980s, however, moose populations have rebounded, thanks to regrowth of plentiful food sources,[25] abandonment of farmland, better land management, clean-up of pollution, and natural dispersal from the Canadian Maritimes and Quebec. South of the Canada–US border, Maine has most of the population with a 2012 headcount of about 76,000 moose.[26] Dispersals from Maine over the years have resulted in healthy, growing populations each in Vermont and New Hampshire, notably near bodies of water and as high up as 3,000 ft (910 m) above sea level in the mountains. In Massachusetts, moose had gone extinct by 1870, but re-colonized the state in the 1960s, with the population expanding from Vermont and New Hampshire; by 2010, the population was estimated at 850–950.[27] Moose reestablished populations in eastern New York and Connecticut and appeared headed south towards the Catskill Mountains, a former habitat.[28][29][30][31][32][33][34]
|
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+
|
37 |
+
In the Midwest U.S., moose are primarily limited to the upper Great Lakes region, but strays, primarily immature males, have been found as far south as eastern Iowa.[35] For unknown reasons, the moose population is declining rapidly in the Midwest.[18]
|
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+
|
39 |
+
Moose were successfully introduced on Newfoundland in 1878 and 1904,[36] where they are now the dominant ungulate, and somewhat less successfully on Anticosti Island in the Gulf of Saint Lawrence.
|
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+
|
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Since the 1990s, moose populations have declined dramatically in much of temperate North America, although they remain stable in arctic and subarctic regions.[37] The exact causes of specific die-offs are not determined, but most documented mortality events were due to wolf predation, bacterial infection due to injuries sustained from predators, and parasites from whitetail deer to which moose have not developed a natural defense, such as liver flukes, brain worms and winter tick infestations.[18][38] Predation of moose calves by brown bear is also significant.[39] One of the leading hypotheses among biologists for generalized, nonhunting declines in moose populations at the southern extent of their range is increasing heat stress brought on by the rapid seasonal temperature upswings as a result of human-induced climate change.[40] Biologists studying moose populations typically use warm-season, heat-stress thresholds of between 14 and 24 °C (57 and 75 °F).[41] However, the minor average temperature increase of 0.83–1.11 °C (1.5–2 °F), over the last 100 years, has resulted in milder winters that induce favorable conditions for ticks, parasites and other invasive species to flourish within the southern range of moose habitat in North America.[40] This leading hypothesis is supported by mathematical models that explore moose-population responses to future climate-change projections.[42]
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The moose population in New Hampshire fell from 7,500 in the early 2000s to a current estimate of 4,000 and in Vermont the numbers were down to 2,200 from a high of 5,000 animals in 2005. Much of the decline has been attributed to the winter tick with about 70% of the moose calf deaths across Maine and New Hampshire due to the parasite. Moose with heavy tick infections will rub their fur down to the skin raw trying to get the ticks off, making them look white when their outer coat rubs off. Locals call them ghost moose.[43][44][45][46] Loss of the insulating winter coat through attempts to rid the moose of winter tick increases the risk of hypothermia in winter.[47]
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In Europe, moose are currently found in large numbers throughout Norway, Sweden, Finland, Latvia, Estonia, Poland, with more modest numbers in the southern Czech Republic, Belarus and northern Ukraine. They are also widespread through Russia on up through the borders with Finland south towards the border with Estonia, Belarus and Ukraine and stretching far away eastwards to the Yenisei River in Siberia. The European moose was native to most temperate areas with suitable habitat on the continent and even Scotland from the end of the last Ice Age, as Europe had a mix of temperate boreal and deciduous forest. Up through Classical times, the species was certainly thriving in both Gaul and Magna Germania, as it appears in military and hunting accounts of the age. However, as the Roman era faded into medieval times, the beast slowly disappeared: soon after the reign of Charlemagne, the moose disappeared from France, where its range extended from Normandy in the north to the Pyrenees in the south. Farther east, it survived in Alsace and the Netherlands until the 9th century as the marshlands in the latter were drained and the forests were cleared away for feudal lands in the former. It was gone from Switzerland by the year 1000, from the western Czech Republic by 1300, from Mecklenburg in Germany by c. 1600, and from Hungary and the Caucasus since the 18th and 19th century, respectively.
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By the early 20th century, the very last strongholds of the European moose appeared to be in Fennoscandian areas and patchy tracts of Russia, with a few migrants found in what is now Estonia and Lithuania. The USSR and Poland managed to restore portions of the range within its borders (such as the 1951 reintroduction into Kampinos National Park and the later 1958 reintroduction in Belarus), but political complications limited the ability to reintroduce it to other portions of its range. Attempts in 1930 and again in 1967 in marshland north of Berlin were unsuccessful. At present in Poland, populations are recorded in the Biebrza river valley, Kampinos, and in Białowieża Forest. It has migrated into other parts of Eastern Europe and has been spotted in eastern and southern Germany.[48][49] Unsuccessful thus far in recolonizing these areas via natural dispersal from source populations in Poland, Belarus, Ukraine, Czech Republic and Slovakia, it appears to be having more success migrating south into the Caucasus. It is listed under Appendix III of the Bern Convention.[50][51]
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In 2008, two moose were reintroduced into the Scottish Highlands[52][53] in Alladale Wilderness Reserve.[54]
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The East Asian moose populations confine themselves mostly to the territory of Russia, with much smaller populations in Mongolia and Northeastern China. Moose populations are relatively stable in Siberia and increasing on the Kamchatka Peninsula. In Mongolia and China, where poaching took a great toll on moose, forcing them to near extinction, they are protected, but enforcement of the policy is weak and demand for traditional medicines derived from deer parts is high.[citation needed] In 1978, the Regional Hunting Department transported 45 young moose to the center of Kamchatka. These moose were brought from Chukotka, home to the largest moose on the planet. Kamchatka now regularly is responsible for the largest trophy moose shot around the world each season. As it is a fertile environment for moose, with a milder climate, less snow, and an abundance of food, moose quickly bred and settled along the valley of the Kamchatka River and many surrounding regions. The population in the past 20 years has risen to over 2,900 animals.
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The size of the moose varies. Following Bergmann's rule, population in the south (A. a. cameloides) usually grow smaller, while moose in the north and northeast (A. a. buturlini) can match the imposing sizes of the Alaskan moose (A. a. gigas) and are prized by trophy hunters.[citation needed]
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In 1900, an attempt to introduce moose into the Hokitika area failed; then in 1910 ten moose (four bulls and six cows) were introduced into Fiordland. This area is considered a less than suitable habitat, and subsequent low numbers of sightings and kills have led to some presumption of this population's failure.[55] The last proven sighting of a moose in New Zealand was in 1952.[56] However, a moose antler was found in 1972, and DNA tests showed that hair collected in 2002 was from a moose. There has been extensive searching, and while automated cameras failed to capture photographs, evidence was seen of bedding spots, browsing, and antler marks.[57]
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North America:
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Europe and Asia:
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Bull moose have antlers like other members of the deer family. Cows select mates based on antler size. Bull moose use dominant displays of antlers to discourage competition and will spar or fight rivals.[84] The size and growth rate of antlers is determined by diet and age; symmetry reflects health.[84]
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The male's antlers grow as cylindrical beams projecting on each side of the head at right angles to the midline of the skull, and then fork. The lower prong of this fork may be either simple, or divided into two or three tines, with some flattening. Most moose have antlers that are broad and palmate (flat) with tines (points) along the outer edge.[84]. Within the ecologic range of the moose in Europe, those in northerly locales display the palmate pattern of antlers, while the antlers of European moose over the southerly portion of its range are typically of the cervina dendritic pattern and comparatively small, perhaps due to evolutionary pressures of hunting by humans, who prize the large palmate antlers. European moose with antlers intermediate between the palmate and the dendritic form are found in the middle of the north-south range.[85] Moose with antlers have more acute hearing than those without antlers; a study of trophy antlers using a microphone found that the palmate antler acts as a parabolic reflector, amplifying sound at the moose's ear.[86]
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The antlers of mature Alaskan adult bull moose (5 to 12 years old) have a normal maximum spread greater than 200 centimeters (79 in). By the age of 13, moose antlers decline in size and symmetry. The widest spread recorded was 210 centimeters (83 in) across. (An Alaskan moose also holds the record for the heaviest weight at 36 kilograms (79 lb).)[84]
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Antler beam diameter, not the number of tines, indicates age.[84] In North America, moose (A. a. americanus) antlers are usually larger than those of Eurasian moose and have two lobes on each side, like a butterfly. Eurasian moose antlers resemble a seashell, with a single lobe on each side.[84] In the North Siberian moose (A. a. bedfordiae), the posterior division of the main fork divides into three tines, with no distinct flattening. In the common moose (A. a. alces) this branch usually expands into a broad palmation, with one large tine at the base and a number of smaller snags on the free border. There is, however, a Scandinavian breed of the common moose in which the antlers are simpler and recall those of the East Siberian animals. The palmation appears to be more marked in North American moose than in the typical Scandinavian moose.
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After the mating season males drop their antlers to conserve energy for the winter. A new set of antlers will then regrow in the spring. Antlers take three to five months to fully develop, making them one of the fastest growing animal organs. Antler growth is "nourished by an extensive system of blood vessels in the skin covering, which contains numerous hair follicles that give it a 'velvet' texture."[84] This requires intense grazing on a highly-nutritious diet. By September the velvet is removed by rubbing and thrashing which changes the colour of the antlers. Immature bulls may not shed their antlers for the winter, but retain them until the following spring. Birds, carnivores and rodents eat dropped antlers as they are full of protein and moose themselves will eat antler velvet for the nutrients.[84]
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If a bull moose is castrated, either by accidental or chemical means, he will quickly shed his current set of antlers and then immediately begin to grow a new set of misshapen and deformed antlers that he will wear the rest of his life without ever shedding again. The distinctive-looking appendages (often referred to as "devil's antlers") are the source of several myths and legends among many groups of Inuit as well as several other tribes of indigenous peoples of North America.[87]
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In extremely rare circumstances, a cow moose may grow antlers. This is usually attributed to a hormone imbalance.[88]
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The moose proboscis is distinctive among the living cervids due to its large size; it also features nares that can be sealed shut when the moose is browsing aquatic vegetation. The moose proboscis likely evolved as an adaptation to aquatic browsing, with loss of the rhinarium, and development of a superior olfactory column separate from an inferior respiratory column.[89] This separation contributes to the moose's keen sense of smell, which they employ to detect water sources, to find food under snow, and to detect mates or predators.[90][89]
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As with all members of the order Artiodactyla (even-toed ungulates), moose feet have two large keratinized hooves corresponding to the third and fourth toe, with two small posterolateral dewclaws (vestigial digits), corresponding to the second and fifth toe. The hoof of the fourth digit is broader than that of the third digit, while the inner hoof of the third digit is longer than that of the fourth digit. This foot configuration may favor striding on soft ground.[91] The moose hoof splays under load, increasing surface area, which limits sinking of the moose foot into soft ground or snow, and which increases efficiency when swimming. The body weight per footprint surface area of the moose foot is intermediate between that of the pronghorn foot, (which have stiff feet lacking dewclaws—optimized for high-speed running) and the caribou foot (which are more rounded with large dewclaws, optimized for walking in deep snow). The moose's body weight per surface area of footprint is about twice that of the caribou's.[92]
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[93][citation needed]
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On firm ground, a bull moose leaves a visible impression of the dewclaws in its footprint, while a cow moose or calf does not leave a dewclaw impression. On soft ground or mud, bull, cow, and calf footprints may all show dewclaw impressions.[citation needed]
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Their fur consist of two layers; top layer of long guard hairs and a soft wooly undercoat. The guard hairs are hollow and filled with air for better insulation, which also helps them stay afloat when swimming.[94]
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Both male and female moose have a dewlap or bell,[95] which is a fold of skin under the chin. Its exact use is unknown, but theories state that it might be used in mating, as a visual and olfactory signal, or as a dominance signal by males, as are the antlers.[96]
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The tail is short (6 cm to 8 cm in length) and vestigial in appearance; unlike other ungulates the moose tail is too short to swish away insects.[97]
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On average, an adult moose stands 1.4–2.1 m (4.6–6.9 ft) high at the shoulder, which is more than a foot higher than the next largest deer on average, the wapiti.[98] Males (or "bulls") normally weigh from 380 to 700 kg (838 to 1,543 lb) and females (or "cows") typically weigh 200 to 490 kg (441 to 1,080 lb), depending on racial or clinal as well as individual age or nutritional variations.[99][100] The head-and-body length is 2.4–3.1 m (7.9–10.2 ft), with the vestigial tail adding only a further 5–12 cm (2.0–4.7 in).[101] The largest of all the races is the Alaskan subspecies (A. a. gigas), which can stand over 2.1 m (6.9 ft) at the shoulder, has a span across the antlers of 1.8 m (5.9 ft) and averages 634.5 kg (1,399 lb) in males and 478 kg (1,054 lb) in females.[102] Typically, however, the antlers of a mature bull are between 1.2 m (3.9 ft) and 1.5 m (4.9 ft). The largest confirmed size for this species was a bull shot at the Yukon River in September 1897 that weighed 820 kg (1,808 lb) and measured 2.33 m (7.6 ft) high at the shoulder.[103] There have been reported cases of even larger moose, including a bull killed in 2004 that weighed 1,043 kg (2,299 lb)[104], and a bull that reportedly scaled 1,180 kg (2,601 lb), but none are authenticated and some may not be considered reliable.[103] Behind only the two species of bison, the moose is the second largest of extant terrestrial wildlife after the bisons in North America, Siberia,[105] and Europe.
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The moose is a herbivore and is capable of consuming many types of plant or fruit. The average adult moose needs to consume 9,770 kcal (40.9 MJ) per day to maintain its body weight.[106] Much of a moose's energy is derived from terrestrial vegetation, mainly consisting of forbs and other non-grasses, and fresh shoots from trees such as willow and birch. These plants are rather low in sodium, and moose generally need to consume a good quantity of aquatic plants. While much lower in energy, aquatic plants provide the moose with its sodium requirements, and as much as half of their diet usually consists of aquatic plant life.[107] In winter, moose are often drawn to roadways, to lick salt that is used as a snow and ice melter.[108] A typical moose, weighing 360 kg (794 lb), can eat up to 32 kg (71 lb) of food per day.[107]
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Moose lack upper front teeth, but have eight sharp incisors on the lower jaw. They also have a tough tongue, lips and gums, which aid in the eating of woody vegetation. Moose have six pairs of large, flat molars and, ahead of those, six pairs of premolars, to grind up their food. A moose's upper lip is very sensitive, to help distinguish between fresh shoots and harder twigs, and is prehensile, for grasping their food. In the summer, moose may use this prehensile lip for grabbing branches and pulling, stripping the entire branch of leaves in a single mouthful, or for pulling forbs, like dandelions, or aquatic plants up by the base, roots and all.[109][110] A moose's diet often depends on its location, but they seem to prefer the new growths from deciduous trees with a high sugar content, such as white birch, trembling aspen and striped maple, among many others.[111] To reach high branches, a moose may bend small saplings down, using its prehensile lip, mouth or body. For larger trees a moose may stand erect and walk upright on its hind legs, allowing it to reach branches up to 4.26 meters (14.0 ft) or higher above the ground.[112][113]
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Moose also eat many aquatic plants, including lilies and pondweed.[114] Moose are excellent swimmers and are known to wade into water to eat aquatic plants. This trait serves a second purpose in cooling down the moose on summer days and ridding itself of black flies. Moose are thus attracted to marshes and river banks during warmer months as both provide suitable vegetation to eat and water to wet themselves in. Moose have been known to dive underwater to reach plants on lake bottoms, and the complex snout may assist the moose in this type of feeding. Moose are the only deer that are capable of feeding underwater.[115] As an adaptation for feeding on plants underwater, the nose is equipped with fatty pads and muscles that close the nostrils when exposed to water pressure, preventing water from entering the nose.[116] Other species can pluck plants from the water too, but these need to raise their heads in order to swallow.
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Moose are not grazing animals but browsers (concentrate selectors). Like giraffes, moose carefully select foods with less fiber and more concentrations of nutrients. Thus, the moose's digestive system has evolved to accommodate this relatively low-fiber diet. Unlike most hooved, domesticated animals (ruminants), moose cannot digest hay, and feeding it to a moose can be fatal.[117][118] The moose's varied and complex diet is typically expensive for humans to provide, and free-range moose require a lot of forested acreage for sustainable survival, which is one of the main reasons moose have never been widely domesticated.[citation needed]
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A full-grown moose has few enemies except Siberian tigers (Panthera tigris altaica) which regularly prey on adult moose,[119][120][121] but a pack of gray wolves (Canis lupus) can still pose a threat, especially to females with calves.[122] Brown bears (Ursus arctos)[102] are also known to prey on moose of various sizes and are the only predator besides the wolf to attack moose both in Eurasia and North America. However, brown bears are more likely to take over a wolf kill or to take young moose than to hunt adult moose on their own.[123][124][125] American black bears (Ursus americanus) and cougars (Puma concolor) can be significant predators of moose calves in May and June and can, in rare instances, prey on adults (mainly cows rather than the larger bulls).[126][127] Wolverine (Gulo gulo) are most likely to eat moose as carrion but have killed moose, including adults, when the large ungulates are weakened by harsh winter conditions.[128][129] Killer whales (Orcinus orca) are the moose's only known marine predator as they have been known to prey on moose swimming between islands out of North America's Northwest Coast,[130] however, there is at least one recorded instance of a moose preyed upon by a Greenland shark.[131]
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In some areas, moose are the primary source of food for wolves. Moose usually flee upon detecting wolves. Wolves usually follow moose at a distance of 100 to 400 meters (330 to 1,310 ft), occasionally at a distance of 2 to 3 kilometers (1.2 to 1.9 mi). Attacks from wolves against young moose may last seconds, though sometimes they can be drawn out for days with adults. Sometimes, wolves will chase moose into shallow streams or onto frozen rivers, where their mobility is greatly impeded. Moose will sometimes stand their ground and defend themselves by charging at the wolves or lashing out at them with their powerful hooves. Wolves typically kill moose by tearing at their haunches and perineum, causing massive blood loss. Occasionally, a wolf may immobilise a moose by biting its sensitive nose, the pain of which can paralyze a moose.[132] Wolf packs primarily target calves and elderly animals, but can and will take healthy, adult moose. Moose between the ages of two and eight are seldom killed by wolves.[133] Though moose are usually hunted by packs, there are cases in which single wolves have successfully killed healthy, fully-grown moose.[134][135]
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Research into moose predation suggests that their response to perceived threats is learned rather than instinctual. In practical terms this means moose are more vulnerable in areas where wolf or bear populations were decimated in the past but are now rebounding. These same studies suggest, however, that moose learn quickly and adapt, fleeing an area if they hear or smell wolves, bears, or scavenger birds such as ravens.[136]
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Moose are also subject to various diseases and forms of parasitism. In northern Europe, the moose botfly is a parasite whose range seems to be spreading.[137]
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Moose are mostly diurnal. They are generally solitary with the strongest bonds between mother and calf. Although moose rarely gather in groups, there may be several in close proximity during the mating season.
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Rutting and mating occurs in September and October. During the rut, mature bulls will cease feeding completely for a period of approximately two weeks; this fasting behavior has been attributed to neurophysiological changes related to redeployment of olfaction for detection of moose urine and moose cows.[138] The males are polygamous and will seek several females to breed with. During this time both sexes will call to each other. Males produce heavy grunting sounds that can be heard from up to 500 meters away, while females produce wail-like sounds.[139] Males will fight for access to females. Initially, the males assess which of them is dominant and one bull may retreat, however, the interaction can escalate to a fight using their antlers.
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Female moose have an eight-month gestation period, usually bearing one calf, or twins if food is plentiful,[140] in May or June.[141] Twinning can run as high as 30% to 40% with good nutrition[142] Newborn moose have fur with a reddish hue in contrast to the brown appearance of an adult. The young will stay with the mother until just before the next young are born. The life span of an average moose is about 15–25 years. Moose populations are stable at 25 calves for every 100 cows at 1 year of age. With availability of adequate nutrition, mild weather, and low predation, moose have a huge potential for population expansion.[142]
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(newborn)Calves nursing in spring.
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(3 months)Calves stay near their mothers at all times.
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(9 months)This calf is almost ready to leave its mother.
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(10–11 months)This yearling was probably recently chased away by its pregnant mother.
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Moose are not usually aggressive towards humans, but can be provoked or frightened to behave with aggression. In terms of raw numbers, they attack more people than bears and wolves combined, but usually with only minor consequences. In the Americas, moose injure more people than any other wild mammal, and worldwide, only hippopotamuses injure more.[143] When harassed or startled by people or in the presence of a dog, moose may charge. Also, as with bears or any wild animal, moose that have become used to being fed by people may act aggressively when denied food. During the fall mating season, bulls may be aggressive toward humans because of the high hormone levels they experience. Cows with young calves are very protective and will attack humans who come too close, especially if they come between mother and calf. Unlike other dangerous animals, moose are not territorial, and do not view humans as food, and will therefore usually not pursue humans if they simply run away.[144]
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Like any wild animal, moose are unpredictable. They are most likely to attack if annoyed or harassed, or if approached too closely. A moose that has been harassed may vent its anger on anyone in the vicinity, and they often do not make distinctions between their tormentors and innocent passers-by.[citation needed] Moose are very limber animals with highly flexible joints and sharp, pointed hooves, and are capable of kicking with both front and back legs. Unlike other large, hooved mammals, such as horses, moose can kick in all directions including sideways. Therefore, there is no safe side from which to approach. However, moose often give warning signs prior to attacking, displaying their aggression by means of body language. Maintained eye contact is usually the first sign of aggression, while laid-back ears or a lowered head is a definite sign of agitation. If the hairs on the back of the moose's neck and shoulders (hackles) stand up, a charge is usually imminent. The Anchorage Visitor Centers warn tourists that "...a moose with its hackles raised is a thing to fear."[145][146][147][148]
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Studies suggest that the calls made by female moose during the rut not only call the males but can actually induce a bull to invade another bull's harem and fight for control of it. This in turn means that the cow moose has at least a small degree of control over which bulls she mates with.[149]
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Moose often show aggression to other animals as well; especially predators. Bears are common predators of moose calves and, rarely, adults. Alaskan moose have been reported to successfully fend off attacks from both black and brown bears. Moose have been known to stomp attacking wolves, which makes them less preferred as prey to the wolves. Moose are fully capable of killing bears and wolves. A moose of either sex that is confronted by danger may let out a loud roar, more resembling that of a predator than a prey animal. European moose are often more aggressive than North American moose, such as the moose in Sweden, which often become very agitated at the sight of a predator. However, like all ungulates known to attack predators, the more aggressive individuals are always darker in color.[115]
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European rock drawings and cave paintings reveal that moose have been hunted since the Stone Age. Excavations in Alby, Sweden, adjacent to the Stora Alvaret have yielded moose antlers in wooden hut remains from 6000 BCE, indicating some of the earliest moose hunting in northern Europe. In northern Scandinavia one can still find remains of trapping pits used for hunting moose. These pits, which can be up to 4 m × 7 m (13 ft 1 in × 23 ft 0 in) in area and 2 m (6 ft 7 in) deep, would have been camouflaged with branches and leaves. They would have had steep sides lined with planks, making it impossible for the moose to escape once it fell in. The pits are normally found in large groups, crossing the moose's regular paths and stretching over several km. Remains of wooden fences designed to guide the animals toward the pits have been found in bogs and peat. In Norway, an early example of these trapping devices has been dated to around 3700 BC. Trapping elk in pits is an extremely effective hunting method. As early as the 16th century the Norwegian government tried to restrict their use; nevertheless, the method was in use until the 19th century.
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The earliest recorded description of the moose is in Julius Caesar's Commentarii de Bello Gallico, where it is described thus:
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There are also [animals], which are called moose. The shape of these, and the varied color of their skins, is much like roes, but in size they surpass them a little and are destitute of horns, and have legs without joints and ligatures; nor do they lie down for the purpose of rest, nor, if they have been thrown down by any accident, can they raise or lift themselves up. Trees serve as beds to them; they lean themselves against them, and thus reclining only slightly, they take their rest; when the huntsmen have discovered from the footsteps of these animals whither they are accustomed to betake themselves, they either undermine all the trees at the roots, or cut into them so far that the upper part of the trees may appear to be left standing. When they have leant upon them, according to their habit, they knock down by their weight the unsupported trees, and fall down themselves along with them.[150]
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In book 8, chapter 16 of Pliny the Elder's Natural History from 77 CE, the elk and an animal called achlis, which is presumably the same animal, are described thus:
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... there is, also, the moose, which strongly resembles our steers, except that it is distinguished by the length of the ears and of the neck. There is also the achlis, which is produced in the land of Scandinavia; it has never been seen in this city, although we have had descriptions of it from many persons; it is not unlike the moose, but has no joints in the hind leg. Hence, it never lies down, but reclines against a tree while it sleeps; it can only be taken by previously cutting into the tree, and thus laying a trap for it, as otherwise, it would escape through its swiftness. Its upper lip is so extremely large, for which reason it is obliged to go backwards when grazing; otherwise, by moving onwards, the lip would get doubled up.[151]
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Moose are hunted as a game species in many of the countries where they are found. Moose meat tastes, wrote Henry David Thoreau in "The Maine Woods", "like tender beef, with perhaps more flavour; sometimes like veal". While the flesh has protein levels similar to those of other comparable red meats (e.g. beef, deer and wapiti), it has a low fat content, and the fat that is present consists of a higher proportion of polyunsaturated fats than saturated fats.[152]
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Dr. Valerius Geist, who emigrated to Canada from the Soviet Union, wrote in his 1999 book Moose: Behaviour, Ecology, Conservation:
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In Sweden, no fall menu is without a mouthwatering moose dish. The Swedes fence their highways to reduce moose fatalities and design moose-proof cars. Sweden is less than half as large as the Canadian province of British Columbia, but the annual take of moose in Sweden—upward of 150,000—is twice that of the total moose harvest in North America.
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Boosting moose populations in Alaska for hunting purposes is one of the reasons given for allowing aerial or airborne methods to remove wolves in designated areas, e.g., Craig Medred: "A kill of 124 wolves would thus translate to [the survival of] 1488 moose or 2976 caribou or some combination thereof".[153] Some scientists believe that this artificial inflation of game populations is actually detrimental to both caribou and moose populations as well as the ecosystem as a whole. This is because studies have shown[citation needed] that when these game populations are artificially boosted, it leads to both habitat destruction and a crash in these populations.[154]
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Cadmium levels are high in Finnish elk liver and kidneys, with the result that consumption of these organs from elk more than one year old is prohibited in Finland.[155] As a result of a study reported in 1988, the Ontario Ministry of Natural Resources recommended against the consumption of moose and deer kidneys and livers. Levels of cadmium were found to be considerably higher than in Scandinavia.[156] The New Brunswick Department of Natural Resources advises hunters not to consume cervid offal.[157]
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Cadmium intake has been found to be elevated amongst all consumers of elk meat, though the elk meat was found to contribute only slightly to the daily cadmium intake. However the consumption of moose liver or kidneys significantly increased cadmium intake, with the study revealing that heavy consumers of moose organs have a relatively narrow safety margin below the levels which would probably cause adverse health effects.[158]
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The center of mass of a moose is above the hood of most passenger cars. In a collision, the impact crushes the front roof beams and individuals in the front seats.[159] Collisions of this type are frequently fatal; seat belts and airbags offer little protection.[160] In collisions with higher vehicles (such as trucks), most of the deformation is to the front of the vehicle and the passenger compartment is largely spared. Moose collisions have prompted the development of a vehicle test referred to as the "moose test" (Swedish: Älgtest, German: Elchtest). A Massachusetts study found that moose–vehicular collisions had a very high human fatality rate and that such collisions caused the death of 3% of the Massachusetts moose population annually.[161]
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Moose warning signs are used on roads in regions where there is a danger of collision with the animal. The triangular warning signs common in Sweden, Norway, and Finland have become coveted souvenirs among tourists traveling in these countries, causing road authorities so much expense that the moose signs have been replaced with imageless generic warning signs in some regions.[162]
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In Ontario, Canada, an estimated 265 moose die each year as a result of collision with trains. Moose–train collisions were more frequent in winters with above-average snowfall.[163] In January 2008, the Norwegian newspaper Aftenposten estimated that some 13,000 moose had died in collisions with Norwegian trains since 2000. The state agency in charge of railroad infrastructure (Jernbaneverket) plans to spend 80 million Norwegian kroner to reduce collision rate in the future by fencing the railways, clearing vegetation from near the tracks, and providing alternative snow-free feeding places for the animals elsewhere.[164]
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In the Canadian province of New Brunswick, collisions between automobiles and moose are frequent enough that all new highways have fences to prevent moose from accessing the road, as has long been done in Finland, Norway, and Sweden. A demonstration project, Highway 7 between Fredericton and Saint John, which has one of the highest frequencies of moose collisions in the province, did not have these fences until 2008, although it was and continues to be extremely well signed.[165][166] Newfoundland and Labrador recommended that motorists use caution between dusk and dawn because that is when moose are most active and most difficult to see, increasing the risk of collisions.[167] Local moose sightings are often reported on radio stations so that motorists can take care while driving in particular areas. An electronic "moose detection system" was installed on two sections of the Trans-Canada Highway in Newfoundland in 2011, but the system proved unreliable and was removed in 2015.[168]
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In Sweden, a road will not be fenced unless it experiences at least one moose accident per km per year.[169]
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In eastern Germany, where the scarce population is slowly increasing, there were two road accidents involving moose since 2000.[50]
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Domestication of moose was investigated in the Soviet Union before World War II. Early experiments were inconclusive, but with the creation of a moose farm at Pechora-Ilych Nature Reserve in 1949, a small-scale moose domestication program was started, involving attempts at selective breeding of animals on the basis of their behavioural characteristics. Since 1963, the program has continued at Kostroma Moose Farm, which had a herd of 33 tame moose as of 2003. Although at this stage the farm is not expected to be a profit-making enterprise, it obtains some income from the sale of moose milk and from visiting tourist groups. Its main value, however, is seen in the opportunities it offers for the research in the physiology and behavior of the moose, as well as in the insights it provides into the general principles of animal domestication.
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In Sweden, there was a debate in the late 18th century about the national value of using the moose as a domestic animal. Among other things, the moose was proposed to be used in postal distribution, and there was a suggestion to develop a moose-mounted cavalry. Such proposals remained unimplemented, mainly because the extensive hunting for moose that was deregulated in the 1790s nearly drove it to extinction. While there have been documented cases of individual moose being used for riding and/or pulling carts and sleds, Björklöf concludes no wide-scale usage has occurred outside fairy tales.[170]
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Moose are an old genus. Like its relatives, Odocoileus and Capreolus, the genus Alces gave rise to very few species that endured for long periods of time. This differs from the Megacerines, such as the Irish elk, which evolved many species before going extinct. Some scientists, such as Adrian Lister, grouped all the species into one genus, while others, such as Augusto Azzaroli, used Alces for the living species, placing the fossil species into the genera Cervalces and Libralces.
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The earliest known species is Libralces gallicus (French moose), which lived in the Pliocene epoch, about 2 million years ago. Libralces gallicus came from the warm savannas of Pliocene Europe, with the best-preserved skeletons being found in southern France. L. gallicus was 1.25 times larger than the Alaskan moose in linear dimensions, making it nearly twice as massive. L. gallicus had many striking differences from its modern descendants. It had a longer, narrower snout and a less-developed nasal cavity, more resembling that of a modern deer, lacking any sign of the modern moose-snout. Its face resembled that of the modern wapiti. However, the rest of its skull structure, skeletal structure and teeth bore strong resemblance to those features that are unmistakable in modern moose, indicating a similar diet. Its antlers consisted of a horizontal bar 2.5 m (8 ft 2 3⁄8 in) long, with no tines, ending in small palmations. Its skull and neck structure suggest an animal that fought using high-speed impacts, much like the Dall sheep, rather than locking and twisting antlers the way modern moose combat. Their long legs and bone structure suggest an animal that was adapted to running at high speeds over rough terrain.[171][172]
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Libralces existed until the middle Pleistocene epoch and were followed briefly by a species called Cervalces carnutorum. The main differences between the two consisted of shortening of the horizontal bar in the antlers and broadening of the palmations, indicating a likely change from open plains to more forested environments, and skeletal changes that suggest an adaptation to marshy environments.
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Cervalces carnutorum was soon followed by a much larger species called Cervalces latifrons (broad-fronted stag-moose). The Pleistocene epoch was a time of gigantism, in which most species were much larger than their descendants of today, including exceptionally large lions, hippopotamuses, mammoths, and deer. Many fossils of Cervalces latifrons have been found in Siberia, dating from about 1.2 to 0.5 million years ago. This is most likely the time at which the species migrated from the Eurasian continent to North America. Like its descendants, it inhabited mostly northern latitudes, and was probably well-adapted to the cold. Cervalces latifrons was the largest deer known to have ever existed, standing more than 2.1 m (6 ft 10 5⁄8 in) tall at the shoulders. This is bigger than even the Irish elk (megacerine), which was 1.8 m (5 ft 10 7⁄8 in) tall at the shoulders. Its antlers were smaller than the Irish elk's, but comparable in size to those of Libralces gallicus. However, the antlers had a shorter horizontal bar and larger palmations, more resembling those of a modern moose.[171][172][173]
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Alces alces (the modern moose) appeared during the late Pleistocene epoch. The species arrived in North America at the end of the Pleistocene and coexisted with a late-surviving variety or relative of Cervalces latifrons, which Azzaroli classified as a separate species called Cervalces scotti, or the American stag-moose.[174]
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en/4316.html.txt
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Coordinates: 05h 30m 00s, +00° 00′ 00″
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Orion is a prominent constellation located on the celestial equator and visible throughout the world. It is one of the most conspicuous[1] and recognizable constellations in the night sky.[2] It is named after Orion, a hunter in Greek mythology. Its brightest stars are blue-white Rigel (Beta Orionis) and red Betelgeuse (Alpha Orionis).
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The earliest known depiction linked to the constellation of Orion is a prehistoric (Aurignacian) mammoth ivory carving found in a cave in the Ach valley in West Germany in 1979. Archaeologists estimate that it is 32,000 to 38,000 years old.[3][4][5] The distinctive pattern of Orion is recognized in numerous cultures around the world, and many myths are associated with it. Orion is used as a symbol in the modern world.
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The Babylonian star catalogues of the Late Bronze Age name Orion MULSIPA.ZI.AN.NA,[note 1] "The Heavenly Shepherd" or "True Shepherd of Anu" – Anu being the chief god of the heavenly realms.[6] The Babylonian constellation is sacred to Papshukal and Ninshubur, both minor gods fulfilling the role of 'messenger to the gods'. Papshukal is closely associated with the figure of a walking bird on Babylonian boundary stones, and on the star map the figure of the Rooster is located below and behind the figure of the True Shepherd—both constellations represent the herald of the gods, in his bird and human forms respectively.[7]
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In ancient Egypt, the stars of Orion were regarded as a god, called Sah. Because Orion rises before Sirius, the star whose heliacal rising was the basis for the Solar Egyptian calendar, Sah was closely linked with Sopdet, the goddess who personified Sirius. The god Sopdu is said to be the son of Sah and Sopdet. Sah is syncretized with Osiris, while Sopdet is syncretized with Osiris' mythological wife, Isis. In the Pyramid Texts, from the 24th and 23rd centuries BC, Sah is one of many gods whose form the dead pharaoh is said to take in the afterlife.[8]
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The Armenians identified their legendary patriarch and founder Hayk with Orion. Hayk is also the name of the Orion constellation in the Armenian translation of the Bible.[9]
|
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The Bible mentions Orion three times, naming it "Kesil" (כסיל, literally – fool). Though, this name perhaps is etymologically connected with "Kislev", the name for the ninth month of the Hebrew calendar (i.e. November–December), which, in turn, may derive from the Hebrew root K-S-L as in the words "kesel, kisla" (כֵּסֶל, כִּסְלָה, hope, positiveness), i.e. hope for winter rains.: Job 9:9 ("He is the maker of the Bear and Orion"), Job 38:31 ("Can you loosen Orion's belt?"), and Amos 5:8 ("He who made the Pleiades and Orion").
|
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In ancient Aram, the constellation was known as Nephîlā′, the Nephilim are said to be Orion's descendants.[10]
|
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In Greek mythology, Orion was a gigantic, supernaturally strong hunter,[11] born to Euryale, a Gorgon, and Poseidon (Neptune), god of the sea. One myth recounts Gaia's rage at Orion, who dared to say that he would kill every animal on Earth. The angry goddess tried to dispatch Orion with a scorpion. This is given as the reason that the constellations of Scorpius and Orion are never in the sky at the same time. However, Ophiuchus, the Serpent Bearer, revived Orion with an antidote. This is said to be the reason that the constellation of Ophiuchus stands midway between the Scorpion and the Hunter in the sky.[12]
|
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|
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The constellation is mentioned in Horace's Odes (Ode 3.27.18), Homer's Odyssey (Book 5, line 283) and Iliad, and Virgil's Aeneid (Book 1, line 535)
|
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In medieval Muslim astronomy, Orion was known as al-jabbar, "the giant".[13] Orion's sixth brightest star, Saiph, is named from the Arabic, saif al-jabbar, meaning "sword of the giant".[14]
|
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In China, Orion was one of the 28 lunar mansions Sieu (Xiu) (宿). It is known as Shen (參), literally meaning "three", for the stars of Orion's Belt. (See Chinese constellations)
|
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The Chinese character 參 (pinyin shēn) originally meant the constellation Orion (Chinese: 參宿; pinyin: shēnxiù); its Shang dynasty version, over three millennia old, contains at the top a representation of the three stars of Orion's belt atop a man's head (the bottom portion representing the sound of the word was added later).[15]
|
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|
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The Rig Veda refers to the Orion Constellation as Mriga (The Deer).[16]
|
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It is said that two bright stars in the front and two bright stars in the rear are The hunting dogs, the one comparatively less bright star in the middle and ahead of two front dogs is The hunter and three aligned bright stars are in the middle of all four hunting dogs is The Deer (The Mriga) and three little aligned but less brighter stars is The Baby Deer. The Mriga means Deer, locally known as Harnu in folk parlance. There are many folk songs narrating the Harnu. The Malay called Orion's Belt Bintang Tiga Beradik (the "Three Brother Star").[citation needed]
|
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|
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In India, Nataraja 'the cosmic dancer' (an avatar of Shiva) is seen in the constellation called Orion.[17]
|
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|
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The Jain Symbol carved in Udayagiri and Khandagiri Caves, India in 1st century BCE[18] has striking resemblance with Orion.
|
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Bugis sailors identified the three stars in Orion's Belt as tanra tellué, meaning "sign of three".[19]
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In old Hungarian tradition, "Orion" is known as (magic) Archer (Íjász), or Reaper (Kaszás). In recently rediscovered myths, he is called Nimrod (Hungarian "Nimród"), the greatest hunter, father of the twins "Hunor" and "Magor". The "π" and "o" stars (on upper right) form together the reflex bow or the lifted scythe. In other Hungarian traditions, "Orion's belt" is known as "Judge's stick" (Bírópálca).[20]
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In Scandinavian tradition, "Orion's belt" was known as Frigg's Distaff (friggerock) or Freyja's distaff.[21]
|
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The Finns call Orion's belt and the stars below it Väinämöisen viikate (Väinämöinen's scythe).[22] Another name for the asterism of Alnilam, Alnitak and Mintaka is Väinämöisen vyö (Väinämöinen's Belt) and the stars "hanging" from the belt as Kalevanmiekka (Kaleva's sword).
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In Siberia, the Chukchi people see Orion as a hunter; an arrow he has shot is represented by Aldebaran (Alpha Tauri), with the same figure as other Western depictions.[23]
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The Seri people of northwestern Mexico call the three stars in the belt of Orion Hapj (a name denoting a hunter) which consists of three stars: Hap (mule deer), Haamoja (pronghorn), and Mojet (bighorn sheep). Hap is in the middle and has been shot by the hunter; its blood has dripped onto Tiburón Island.[24]
|
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The same three stars are known in Spain and most of Latin America as "Las tres Marías" (Spanish for "The Three Marys").
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In Puerto Rico, the three stars are known as the "Los Tres Reyes Magos" (Spanish for The three Wise Men).[25]
|
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The Ojibwa (Chippewa) Native Americans call this constellation Kabibona'kan, the Winter Maker, as its presence in the night sky heralds winter.[citation needed]
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To the Lakota Native Americans, Tayamnicankhu (Orion's Belt) is the spine of a bison. The great rectangle of Orion are the bison's ribs; the Pleiades star cluster in nearby Taurus is the bison's head; and Sirius in Canis Major, known as Tayamnisinte, is its tail. Another Lakota myth mentions that the bottom half of Orion, the Constellation of the Hand, represented the arm of a chief that was ripped off by the Thunder People as a punishment from the gods for his selfishness. His daughter offered to marry the person who can retrieve his arm from the sky, so the young warrior Fallen Star (whose father was a star and whose mother was human) returned his arm and married his daughter, symbolizing harmony between the gods and humanity with the help of the younger generation. The index finger is represented by Rigel; the Orion Nebula is the thumb; the Belt of Orion is the wrist; and the star Beta Eridani is the pinky finger.[26]
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The seven primary stars of Orion make up the Polynesian constellation Heiheionakeiki which represents a child's string figure similar to a cat's cradle.
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The imagery of the belt and sword has found its way into popular western culture, for example in the form of the shoulder insignia of the 27th Infantry Division of the United States Army during both World Wars, probably owing to a pun on the name of the division's first commander, Major General John F. O'Ryan.[citation needed]
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The film distribution company Orion Pictures used the constellation as its logo.[27]
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In J. R. R. Tolkien's mythology surrounding Middle-earth, Orion is known as Menelvagor, which is Sindarin for "The Swordsman of the Sky".[28]
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In the movie Blade Runner, the dying replicant Roy Batty introspectively delivers his "Tears in Rain" soliloquy: "I've seen things you people wouldn't believe. Attack ships on fire off the shoulder of Orion. I watched C-beams glitter in the dark near the Tannhäuser Gate. All those moments will be lost in time, like tears in rain. Time to die."
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In the sci-fi television series Haven, the Orionid meteor shower is known as the Hunter Meteor Shower. It comes to Haven once every twenty-seven years when the Barn, a space between two worlds, comes to take the mysterious woman away. When the Barn leaves, the Hunter passes safely overhead and The Troubles end, making the town of Haven a haven for the Troubled again.
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In artistic renderings, the surrounding constellations are sometimes related to Orion: he is depicted standing next to the river Eridanus with his two hunting dogs Canis Major and Canis Minor, fighting Taurus. He is sometimes depicted hunting Lepus the hare. He sometimes is depicted to have a lion's hide in his hand.
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There are alternative ways to visualise Orion. From the Southern Hemisphere, Orion is oriented south-upward, and the belt and sword are sometimes called the saucepan or pot in Australia and New Zealand. Orion's Belt is called Drie Konings (Three Kings) or the Drie Susters (Three Sisters) by Afrikaans speakers in South Africa[29] and are referred to as les Trois Rois (the Three Kings) in Daudet's Lettres de Mon Moulin (1866). The appellation Driekoningen (the Three Kings) is also often found in 17th- and 18th-century Dutch star charts and seaman's guides. The same three stars are known in Spain, Latin America, and the Philippines as "Las Tres Marías" (The Three Marys), and as "Los Tres Reyes Magos" (The three Wise Men) in Puerto Rico.[25]
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Even traditional depictions of Orion have varied greatly. Cicero drew Orion in a similar fashion to the modern depiction. The Hunter held an unidentified animal skin aloft in his right hand; his hand was represented by Omicron2 Orionis and the skin was represented by the 5 stars designated Pi Orionis. Kappa and Beta Orionis represented his left and right knees, while Eta and Lambda Leporis were his left and right feet, respectively. As in the modern depiction, Delta, Epsilon, and Zeta represented his belt. His left shoulder was represented by Alpha Orionis, and Mu Orionis made up his left arm. Lambda Orionis was his head and Gamma, his right shoulder. The depiction of Hyginus was similar to that of Cicero, though the two differed in a few important areas. Cicero's animal skin became Hyginus's shield (Omicron and Pi Orionis), and instead of an arm marked out by Mu Orionis, he holds a club (Chi Orionis). His right leg is represented by Theta Orionis and his left leg is represented by Lambda, Mu, and Epsilon Leporis. Further Western European and Arabic depictions have followed these two models.[23]
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Orion is bordered by Taurus to the northwest, Eridanus to the southwest, Lepus to the south, Monoceros to the east, and Gemini to the northeast. Covering 594 square degrees, Orion ranks twenty-sixth of the 88 constellations in size. The constellation boundaries, as set by Eugène Delporte in 1930, are defined by a polygon of 26 sides. In the equatorial coordinate system, the right ascension coordinates of these borders lie between 04h 43.3m and 06h 25.5m, while the declination coordinates are between 22.87° and −10.97°.[30] The constellation's three-letter abbreviation, as adopted by the International Astronomical Union in 1922, is "Ori".[31]
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Orion is most visible in the evening sky from January to March,[32] winter in the Northern Hemisphere, and summer in the Southern Hemisphere. In the tropics (less than about 8° from the equator), the constellation transits at the zenith.
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In the period May–July (summer in the Northern Hemisphere, winter in the Southern Hemisphere), Orion is in the daytime sky and thus not visible at most latitudes. However, for much of Antarctica in the Southern Hemisphere's winter months, the Sun is below the horizon even at midday. Stars (and thus Orion) are then visible at twilight for a few hours around local noon, low in the North. At the same time of day at the South Pole itself (Amundsen–Scott South Pole Station), Rigel is only 8° above the horizon, and the Belt sweeps just along it. In the Southern Hemisphere's summer months, when Orion is normally visible in the night sky, the constellation is actually not visible in Antarctica because the sun does not set at that time of year south of the Antarctic Circle.[33][34]
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In countries close to the equator (e.g. Kenya, Indonesia, Colombia, Ecuador), Orion appears overhead in December around midnight and in the February evening sky.
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Orion is very useful as an aid to locating other stars. By extending the line of the Belt southeastward, Sirius (α CMa) can be found; northwestward, Aldebaran (α Tau). A line eastward across the two shoulders indicates the direction of Procyon (α CMi). A line from Rigel through Betelgeuse points to Castor and Pollux (α Gem and β Gem). Additionally, Rigel is part of the Winter Circle asterism. Sirius and Procyon, which may be located from Orion by following imaginary lines (see map), also are points in both the Winter Triangle and the Circle.[35]
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Orion's seven brightest stars form a distinctive hourglass-shaped asterism, or pattern, in the night sky. Four stars—Rigel, Betelgeuse, Bellatrix and Saiph—form a large roughly rectangular shape, in the centre of which lie the three stars of Orion's Belt—Alnitak, Alnilam and Mintaka. Descending from the 'belt' is a smaller line of three stars, Orion's Sword (the middle of which is in fact not a star but the Orion Nebula), also known as the hunter's sword.
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Many of the stars are luminous hot blue supergiants, with the stars of the belt and sword forming the Orion OB1 Association. Standing out by its red hue, Betelgeuse may nevertheless be a runaway member of the same group.
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Orion's Belt or The Belt of Orion is an asterism within the constellation. It consists of the three bright stars Zeta (Alnitak), Epsilon (Alnilam), and Delta (Mintaka). Alnitak is around 800 light years away from earth and is 100,000 times more luminous than the Sun; much of its radiation is in the ultraviolet range, which the human eye cannot see.[43] Alnilam is approximately 1340 light years away from Earth, shines with magnitude 1.70, and with ultraviolet light is 375,000 times more luminous than the Sun.[44] Mintaka is 915 light years away and shines with magnitude 2.21. It is 90,000 times more luminous than the Sun and is a double star: the two orbit each other every 5.73 days.[45] In the Northern Hemisphere, Orion's Belt is best visible in the night sky during the month of January around 9:00 pm, when it is approximately around the local meridian.[2]
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Just southwest of Alnitak lies Sigma Orionis, a multiple star system composed of five stars that have a combined apparent magnitude of 3.7 and lying 1150 light years distant. Southwest of Mintaka lies the quadruple star Eta Orionis.
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Orion's Sword contains the Orion Nebula, the Messier 43 nebula, the Running Man Nebula, and the stars Theta Orionis, Iota Orionis, and 42 Orionis.
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Three stars comprise a small triangle that marks the head. The apex is marked by Meissa (Lambda Orionis), a hot blue giant of spectral type O8 III and apparent magnitude 3.54, which lies some 1100 light years distant. Phi-1 and Phi-2 Orionis make up the base. Also nearby is the very young star FU Orionis.
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Stretching north from Betelgeuse are the stars that make up Orion's club. Mu Orionis marks the elbow, Nu and Xi mark the handle of the club, and Chi1 and Chi2 mark the end of the club. Just east of Chi1 is the Mira-type variable red giant U Orionis.
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West from Bellatrix lie six stars all designated Pi Orionis (π1 Ori, π2 Ori, π3 Ori, π4 Ori, π5 Ori and π6 Ori) which make up Orion's shield.
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Around 20 October each year the Orionid meteor shower (Orionids) reaches its peak. Coming from the border with the constellation Gemini as many as 20 meteors per hour can be seen. The shower's parent body is Halley's Comet.[46]
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Hanging from Orion's belt is his sword, consisting of the multiple stars θ1 and θ2 Orionis, called the Trapezium and the Orion Nebula (M42). This is a spectacular object that can be clearly identified with the naked eye as something other than a star. Using binoculars, its clouds of nascent stars, luminous gas, and dust can be observed. The Trapezium cluster has many newborn stars, including several brown dwarfs, all of which are at an approximate distance of 1,500 light-years. Named for the four bright stars that form a trapezoid, it is largely illuminated by the brightest stars, which are only a few hundred thousand years old. Observations by the Chandra X-ray Observatory show both the extreme temperatures of the main stars—up to 60,000 kelvins—and the star forming regions still extant in the surrounding nebula.[47]
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M78 (NGC 2068) is a nebula in Orion. With an overall magnitude of 8.0, it is significantly dimmer than the Great Orion Nebula that lies to its south; however, it is at approximately the same distance, at 1600 light-years from Earth. It can easily be mistaken for a comet in the eyepiece of a telescope. M78 is associated with the variable star V351 Orionis, whose magnitude changes are visible in very short periods of time.[48] Another fairly bright nebula in Orion is NGC 1999, also close to the Great Orion Nebula. It has an integrated magnitude of 10.5 and is 1500 light-years from Earth. The variable star V380 Orionis is embedded in NGC 1999.[49]
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Another famous nebula is IC 434, the Horsehead Nebula, near ζ Orionis. It contains a dark dust cloud whose shape gives the nebula its name.
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NGC 2174 is an emission nebula located 6400 light-years from Earth.
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Besides these nebulae, surveying Orion with a small telescope will reveal a wealth of interesting deep-sky objects, including M43, M78, as well as multiple stars including Iota Orionis and Sigma Orionis. A larger telescope may reveal objects such as Barnard's Loop and the Flame Nebula (NGC 2024), as well as fainter and tighter multiple stars and nebulae.
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All of these nebulae are part of the larger Orion Molecular Cloud Complex, which is located approximately 1,500 light-years away and is hundreds of light-years across. It is one of the most intense regions of stellar formation visible within our galaxy.
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Orion is located on the celestial equator, but it will not always be so located due to the effects of precession of the Earth's axis. Orion lies well south of the ecliptic, and it only happens to lie on the celestial equator because the point on the ecliptic that corresponds to the June solstice is close to the border of Gemini and Taurus, to the north of Orion. Precession will eventually carry Orion further south, and by AD 14000 Orion will be far enough south that it will become invisible from the latitude of Great Britain.[50]
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Further in the future, Orion's stars will gradually move away from the constellation due to proper motion. However, Orion's brightest stars all lie at a large distance from the Earth on an astronomical scale—much farther away than Sirius, for example. Orion will still be recognizable long after most of the other constellations—composed of relatively nearby stars—have distorted into new configurations, with the exception of a few of its stars eventually exploding as supernovae, for example Betelgeuse, which is predicted to explode sometime in the next million years.[51]
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Orlando Jonathan Blanchard Copeland Bloom[1] (born 13 January 1977)[2] is an English actor. He made his breakthrough as the character Legolas in The Lord of the Rings film series, a role he reprised in The Hobbit film series. He gained further acclaim by appearing in epic fantasy, historical, and adventure films, most notably as Will Turner in the Pirates of the Caribbean film series.
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Bloom subsequently established himself as a leading man in Hollywood films, with roles such as Paris in Troy (2004) and Balian de Ibelin in Kingdom of Heaven (2005). He currently stars in the Amazon Prime Video series Carnival Row (2019–present).
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He made his professional stage debut in In Celebration at the Duke of York's Theatre in the West End in 2007 and starred in a Broadway adaption of Romeo and Juliet in 2013. In 2009, Bloom was named a UNICEF Goodwill Ambassador. In 2015 he received the BAFTA Britannia Humanitarian Award.[3]
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Bloom was born on 13 January 1977 in Canterbury, Kent, and was named after the 16th-century English composer Orlando Gibbons.[4][5] He has an older sister, Samantha Bloom.[6]
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Orlando initially believed that his biological father was his mother's husband (his stepfather), the Jewish South African-born anti-Apartheid novelist Harry Saul Bloom (1913–1981), who died when Bloom was four years old. However, when he was thirteen, Bloom's mother revealed to him that his biological father was actually Colin Stone, his mother's partner and family friend.[7][8][9] Stone, the principal of the Concorde International language school,[10] became Orlando Bloom's legal guardian after Harry Bloom's death.[9]
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Bloom's mother, Sonia Constance Josephine (née Copeland), was born in Kolkata, India, the daughter of Francis John Copeland, a physician and surgeon, and Betty Constance Josephine (née Walker). Through her, Bloom is a cousin of photographer Sebastian Copeland.[11][12] Bloom's mother's family lived in Tasmania (Australia), Japan, and India, and were of English descent, some of them having originally come from Kent.[7]
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Bloom was brought up in the Church of England.[13] He attended St Peter's Methodist Primary School,[14] then the junior school of the King's School before proceeding to St Edmund's School in Canterbury. Bloom was discovered to be dyslexic,[9][15] and was encouraged by his mother to take art and drama classes.[9] After being spurred into action following his school prize submission to panto actor Richard Sieben in 1992, in 1993, he moved to London to follow a two-year A Level course in Drama, Photography and Sculpture at Fine Arts College, Hampstead. He then joined the National Youth Theatre, spending two seasons there and earning a scholarship to train at the British American Drama Academy.[16] Bloom began acting professionally with television roles in episodes of Casualty and Midsomer Murders,[9] and subsequently made his film debut in Wilde (1997), opposite Stephen Fry, before entering the Guildhall School of Music and Drama in London, where he studied acting.[17]
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Bloom's first appearance on the screen was in a small role, as a rent boy, in the 1997 film Wilde. Two days after graduating from Guildhall in 1999,[18] he was cast in his first major role, playing Legolas in The Lord of the Rings film trilogy (2001–2003).[16] He had originally auditioned for the part of Faramir, who does not appear until the second movie, but the director, Peter Jackson, cast him as Legolas instead. While shooting a scene, he broke a rib after falling off a horse, but eventually recovered and continued shooting.[19] At the same time, Bloom also played a brief role in Ridley Scott's war film Black Hawk Down as PFC Todd Blackburn. In 2002, he was chosen as one of the Teen People "25 Hottest Stars Under 25" and was named People's hottest Hollywood bachelor in the magazine's 2004 list.[16] All members of the cast of the Rings films were nominated for Best Ensemble Acting at the Screen Actors Guild Awards for three years in a row, finally winning in 2003 for the third film, The Lord of the Rings: The Return of the King. Bloom has also won other awards, including European Film Awards, Hollywood Festival Award, Empire Awards and Teen Choice Awards, and has been nominated for many others. Most of Bloom's box office successes have been as part of an ensemble cast.[20]
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Bloom next starred opposite Keira Knightley and Johnny Depp in Pirates of the Caribbean: The Curse of the Black Pearl, which was a blockbuster hit during the summer of 2003. After the success of Pirates, Bloom next took to the screen as Paris, the man who effectively started the Trojan War, in the 2004 Spring blockbuster, Troy opposite Brad Pitt, Eric Bana and Peter O'Toole. He subsequently played the lead roles in Ridley Scott's Kingdom of Heaven and Cameron Crowe's Elizabethtown (both 2005). In 2006, Bloom starred in sequel Pirates of the Caribbean: Dead Man's Chest and in the independently made Haven, of which he was also executive producer. In the same year he was one of the guest stars in the sitcom Extras, in which he portrayed an exaggeratedly arrogant, narcissistic version of himself who had a great loathing for Johnny Depp (his co-star in Pirates of the Caribbean); Bloom pushed for Extras to go further by making his part unlikable, and contributed to the gag about him admiring Depp out of sheer jealousy, that Depp was far more talented than he was, not to mention rated higher than him on the 'top hottest' charts.[21] Also in 2006, Bloom was the most searched male on Google News.[9] As of May 2007, Bloom has appeared in four of the top 15 highest-grossing films of all time.[18]
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Bloom then again portrayed Will Turner, in Pirates of the Caribbean: At World's End, released on 24 May 2007. Bloom, who had intended to become a stage actor after graduating from the Guildhall School of Music and Drama, had stated that he would like to leave films for a time and instead appear in stage roles,[16] and is "avidly looking for the right sort of material that [he] can do something with"[22] and go "back to basics."[23] During the summer of 2007 he appeared in a London revival of In Celebration, a play by David Storey.[23][24] His character was one of three brothers returning home for their parents' 40th wedding anniversary.[25] On 24 August 2007 he made his first ever TV commercial appearance on late-night Japanese TV, promoting the Uno brand of cosmetics maker Shiseido. A "one night only", 2-minute version of the Sci-Fi themed commercial kicked off the product's marketing campaign.[26] In 2008 he signed on to play a small role in the British film An Education[27] but dropped out to take the lead in Johnnie To's film Red Circle. Also in 2009, he was one of many stars to appear in New York, I Love You, which contained 12 short films in one.
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Bloom appeared in The Three Musketeers opposite Milla Jovovich, Logan Lerman, Matthew Macfadyen, Ray Stevenson, Luke Evans, Juno Temple and Christoph Waltz, released in 2011. Bloom reprised his role as Legolas in parts two and three of The Hobbit, Peter Jackson's three-part prequel to The Lord of the Rings trilogy.[28] He made his Broadway stage debut as Romeo in Romeo and Juliet in August 2013 at the Richard Rodgers Theatre.[29] The New York Times theatre critic Ben Brantley described Bloom's performance as "a first-rate Broadway debut" in the title role: "For once, we have a Romeo who evolves substantively, from a posturing youth in love with love, to a man who discovers the startling revelation of real love, with a last-act descent into bilious, bitter anger that verges on madness."[30]
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In October 2011, Bloom stated that he would like to return for a fifth Pirates of the Caribbean film. Bloom did reprise his character, Will Turner, in a supporting role in the fifth film, Pirates of the Caribbean: Dead Men Tell No Tales, which was released in May 2017.[31]
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Bloom starred in the Trafalgar Studios production of Killer Joe in 2018. Bloom is currently starring in the Amazon Prime TV show Carnival Row, which was filming its second season in Prague before temporarily postponing production due to COVID-19. Together with a number of his Lord of the Rings co-stars (plus writer Philippa Boyens and director Peter Jackson), on 1 June 2020 he joined Josh Gad's YouTube series Reunited Apart which reunites the cast of popular movies through video-conferencing, and promotes donations to non-profit charities.[32]
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Most recently, Bloom starred in the war film The Outpost (released July 2020) as captain Benjamin Keating.
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Bloom has said that he tries "not to exclude himself from real life as much as possible."[9] During filming in Morocco for Kingdom of Heaven, Bloom rescued and adopted a dog, Sidi (a black Saluki mix with a white mark on his chest).[33][34]
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In 2004, he became a full member of SGI-UK (the UK branch of Soka Gakkai International), a lay Buddhist association affiliated with the teachings of Nichiren.[35][36]
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Bloom has also been a part of Global Green, an environmental company, since the early 2000s.[20] As part of his environmental involvement, he has renovated his London home to use solar panels, incorporate recycled materials, and use energy efficient lightbulbs.[9][20] Bloom has been approached by UNICEF to act as an international ambassador.[20]
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Involved in modelling work, in 2002 he starred opposite English actress Kate Beckinsale in a Gap television advertisement directed by Cameron Crowe.[37]
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Bloom has a tattoo of the Elvish word "nine" on his right wrist, written in fictional Tengwar Elvish script, a reference to his involvement in The Lord of the Rings as one of the nine members of the Fellowship of the Ring. The other actors of "The Fellowship" (Sean Astin, Sean Bean, Billy Boyd, Ian McKellen, Dominic Monaghan, Viggo Mortensen, and Elijah Wood) got the same tattoo with the exception of John Rhys-Davies whose stunt double got the tattoo instead.[38][39][40] Bloom also has a tattoo of a sun on his lower left abdomen, which he got at age 15 before moving to London.
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On 12 February 2009, Bloom actively participated in the 'Australia Unites' fundraiser to raise support for the victims of the Australian bushfires on 7 February 2009.[41]
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On 12 October 2009, Bloom was named a UNICEF Goodwill Ambassador. He has been involved in the organisation since 2007 and has visited schools and villages in Nepal in support of sanitation and education programmes.[42] Bloom also visited the city of Slovyansk in eastern Ukraine in April 2016 to raise awareness of the education crisis facing children during the military conflict with Russia.[43]
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Bloom has sustained several injuries: he broke his left arm and cracked his skull three times, broke his nose while playing rugby union, broke his right leg skiing in Switzerland, broke his left leg in a motorbike crash, and broke his right wrist while snowboarding.[18] He also broke his back when he slipped trying to reach a roof terrace of a friend's house and fell three floors.[44] He also broke some ribs while shooting The Lord of the Rings.[45]
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On 13 July 2009, four hooded teenagers broke into the Hollywood Hills home of Bloom and stole nearly half a million dollars' worth of possessions. The burglars, dubbed the "Bling Ring," targeted the homes of young celebrities. Most of Bloom's stolen items were retrieved.[46]
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Bloom had an "on-and-off relationship"[47] with American actress Kate Bosworth, from 2003 until ending the relationship in September 2006.[48][49]
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In late 2007, Bloom began dating Australian supermodel Miranda Kerr.[50] They announced their engagement in June 2010,[51] and were married the following month.[52] Kerr gave birth to their son on 6 January 2011 in Los Angeles.[53][54] On 25 October 2013, Bloom and Kerr announced that they had separated several months earlier, and intended to end their marriage. They were divorced by the end of 2013.[55][56]
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Bloom began dating American singer and songwriter Katy Perry in January 2016;[57] however, the couple confirmed in February 2017 that they had split. They resumed their relationship in February 2018,[58][59][60] and became engaged on 14 February 2019.[61][62] On 5 March 2020, it was revealed through the music video for Perry's song "Never Worn White" that the couple are expecting their first child together, a girl.[63]
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Bloom was named the twelfth most influential person in the UK in a 2004 poll of cultural experts conducted for the BBC.[64] Bloom appeared in his home town of Canterbury on 13 July 2010, where he was presented with an honorary degree from the University of Kent at Canterbury Cathedral.[65] For contributions to the film industry, Bloom was inducted into the Hollywood Walk of Fame on 2 April 2014.[66][67] His motion pictures walk of fame star is located at 6927 Hollywood Boulevard.[68]
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1996: Adam Sandler vs. Bob Barker – Happy Gilmore
|
58 |
+
1997: Fairuza Balk vs. Robin Tunney – The Craft
|
59 |
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1998: Will Smith vs. Cockroach – Men in Black
|
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1999: Ben Stiller vs. Puffy the Dog – There's Something About Mary
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2000: Keanu Reeves vs. Laurence Fishburne – The Matrix
|
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2001: Zhang Ziyi vs. Entire bar – Crouching Tiger, Hidden Dragon
|
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2002: Jackie Chan and Chris Tucker vs. Hong Kong gang – Rush Hour 2
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2003: Yoda vs. Christopher Lee – Star Wars: Episode II – Attack of the Clones
|
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2004: Uma Thurman vs. Chiaki Kuriyama – Kill Bill: Volume 1
|
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+
2005: Uma Thurman vs. Daryl Hannah – Kill Bill: Volume 2
|
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2006: Angelina Jolie vs. Brad Pitt – Mr. & Mrs. Smith
|
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2007: Gerard Butler vs. Robert Maillet – 300
|
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2008: Sean Faris vs. Cam Gigandet – Never Back Down
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2009: Robert Pattinson vs. Cam Gigandet – Twilight
|
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2010: Beyoncé Knowles vs. Ali Larter – Obsessed
|
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2011: Robert Pattinson vs. Bryce Dallas Howard and Xavier Samuel – The Twilight Saga: Eclipse
|
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2012: Jennifer Lawrence and Josh Hutcherson vs. Alexander Ludwig – The Hunger Games
|
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2013: Robert Downey Jr., Chris Evans, Mark Ruffalo, Chris Hemsworth, Scarlett Johansson and Jeremy Renner vs. Tom Hiddleston – The Avengers
|
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2014: Orlando Bloom and Evangeline Lilly vs. Orcs – The Hobbit: The Desolation of Smaug
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2015: Dylan O'Brien vs. Will Poulter – The Maze Runner
|
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2016: Ryan Reynolds vs. Ed Skrein – Deadpool
|
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2018: Gal Gadot vs. German soldiers – Wonder Woman
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2019: Brie Larson vs. Gemma Chan – Captain Marvel
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The platypus (Ornithorhynchus anatinus), sometimes referred to as the duck-billed platypus, is a semiaquatic egg-laying mammal endemic to eastern Australia, including Tasmania. The platypus is the sole living representative of its family (Ornithorhynchidae) and genus (Ornithorhynchus), though a number of related species appear in the fossil record.
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Together with the four species of echidna, it is one of the five extant species of monotremes, the only mammals that lay eggs instead of giving birth to live young. Like other monotremes it senses prey through electrolocation. It is one of the few species of venomous mammals, as the male platypus has a spur on the hind foot that delivers a venom capable of causing severe pain to humans. The unusual appearance of this egg-laying, duck-billed, beaver-tailed, otter-footed mammal baffled European naturalists when they first encountered it, and the first scientists to examine a preserved platypus body (in 1799), judged it a fake, made of several animals sewn together.
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The unique features of the platypus make it an important subject in the study of evolutionary biology, and a recognisable and iconic symbol of Australia. It has appeared as a mascot at national events and features on the reverse of the Australian twenty-cent coin, and the platypus is the animal emblem of the state of New South Wales. Until the early 20th century humans hunted the platypus for its fur, but it is now protected throughout its range. Although captive-breeding programs have had only limited success, and the platypus is vulnerable to the effects of pollution, it is not under any immediate threat.
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When the platypus was first encountered by Europeans in 1798, a pelt and sketch were sent back to Great Britain by Captain John Hunter, the second Governor of New South Wales.[3] British scientists' initial hunch was that the attributes were a hoax.[4] George Shaw, who produced the first description of the animal in the Naturalist's Miscellany in 1799, stated it was impossible not to entertain doubts as to its genuine nature,[5] and Robert Knox believed it might have been produced by some Asian taxidermist.[4] It was thought that somebody had sewn a duck's beak onto the body of a beaver-like animal. Shaw even took a pair of scissors to the dried skin to check for stitches.[6]
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The common name "platypus" is the latinisation of the Greek word πλατύπους (platupous), "flat-footed",[7] from πλατύς (platus), "broad, wide, flat"[8] and πούς (pous), "foot".[9][10] Shaw assigned the species the Linnaean name Platypus anatinus when he initially described it, but the genus term was quickly discovered to already be in use as the name of the wood-boring ambrosia beetle genus Platypus.[11] It was independently described as Ornithorhynchus paradoxus by Johann Blumenbach in 1800 (from a specimen given to him by Sir Joseph Banks)[12] and following the rules of priority of nomenclature, it was later officially recognised as Ornithorhynchus anatinus.[11]
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The scientific name Ornithorhynchus anatinus is derived from ορνιθόρυγχος (ornithorhynkhos), which literally means "bird snout" in Greek; and anatinus, which means "duck-like" in Latin.
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There is no universally-agreed plural form of "platypus" in the English language. Scientists generally use "platypuses" or simply "platypus". Colloquially, the term "platypi" is also used for the plural, although this is a form of pseudo-Latin;[6] going by the word's Greek roots the plural would be "platypodes". Early British settlers called it by many names, such as "watermole", "duckbill", and "duckmole".[6] The name "platypus" is occasionally prefixed with the adjective "duck-billed" to form "duck-billed platypus".
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In David Collins's account of the new colony 1788–1801, he describes coming across "an amphibious animal, of the mole species". His account includes a drawing of the animal.[13]
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The body and the broad, flat tail of the platypus are covered with dense, brown fur that traps a layer of insulating air to keep the animal warm.[6][11] The fur is waterproof, and the texture is akin to that of a mole.[14] The platypus uses its tail for storage of fat reserves (an adaptation also found in animals such as the Tasmanian devil[15]). The webbing on the feet is more significant on the front feet and is folded back when walking on land. The elongated snout and lower jaw are covered in soft skin, forming the bill. The nostrils are located on the dorsal surface of the snout, while the eyes and ears are located in a groove set just back from it; this groove is closed when swimming.[11] Platypuses have been heard to emit a low growl when disturbed and a range of other vocalisations have been reported in captive specimens.[6]
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Weight varies considerably from 0.7 to 2.4 kg (1 lb 9 oz to 5 lb 5 oz), with males being larger than females. Males average 50 cm (20 in) in total length, while females average 43 cm (17 in),[11] with substantial variation in average size from one region to another. This pattern does not seem to follow any particular climatic rule and may be due to other environmental factors, such as predation and human encroachment.[16]
|
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+
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The platypus has an average body temperature of about 32 °C (90 °F) rather than the 37 °C (99 °F) typical of placental mammals.[17] Research suggests this has been a gradual adaptation to harsh environmental conditions on the part of the small number of surviving monotreme species rather than a historical characteristic of monotremes.[18][19]
|
23 |
+
|
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+
Modern platypus young have three teeth in each of the maxillae (one premolar and two molars) and dentaries (three molars), which they lose before or just after leaving the breeding burrow;[11] adults have heavily keratinised pads in their place.[11] The first upper and third lower cheek teeth of platypus nestlings are small, each having one principal cusp, while the other teeth have two main cusps.[20] The platypus jaw is constructed differently from that of other mammals, and the jaw-opening muscle is different.[11] As in all true mammals, the tiny bones that conduct sound in the middle ear are fully incorporated into the skull, rather than lying in the jaw as in cynodonts and other pre mammalian synapsids. However, the external opening of the ear still lies at the base of the jaw.[11] The platypus has extra bones in the shoulder girdle, including an interclavicle, which is not found in other mammals.[11] As in many other aquatic and semiaquatic vertebrates, the bones show osteosclerosis, increasing their density to provide ballast.[21] It has a reptilian gait, with the legs on the sides of the body, rather than underneath.[11] When on land, it engages in knuckle-walking on its front feet, to protect the webbing between the toes.[22]
|
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|
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While both male and female platypuses are born with ankle spurs, only the male's spurs deliver venom,[23][24][25]
|
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+
composed largely of defensin-like proteins (DLPs), three of which are unique to the platypus.[26] The DLPs are produced by the immune system of the platypus. The function of defensins is to cause lysis in pathogenic bacteria and viruses, but in platypuses they also are formed into venom for defence. Although powerful enough to kill smaller animals such as dogs, the venom is not lethal to humans, but the pain is so excruciating that the victim may be incapacitated.[26][27] Oedema rapidly develops around the wound and gradually spreads throughout the affected limb. Information obtained from case histories and anecdotal evidence indicates the pain develops into a long-lasting hyperalgesia (a heightened sensitivity to pain) that persists for days or even months.[28][29] Venom is produced in the crural glands of the male, which are kidney-shaped alveolar glands connected by a thin-walled duct to a calcaneus spur on each hind limb. The female platypus, in common with echidnas, has rudimentary spur buds that do not develop (dropping off before the end of their first year) and lack functional crural glands.[11]
|
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|
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The venom appears to have a different function from those produced by non-mammalian species; its effects are not life-threatening to humans, but nevertheless powerful enough to seriously impair the victim. Since only males produce venom and production rises during the breeding season, it may be used as an offensive weapon to assert dominance during this period.[26]
|
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|
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Similar spurs are found on many archaic mammal groups, indicating that this is an ancient characteristic for mammals as a whole, and not exclusive to the platypus or other monotremes.[30]
|
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+
|
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Monotremes are the only mammals (apart from at least one species of dolphin)[31] known to have a sense of electroreception: they locate their prey in part by detecting electric fields generated by muscular contractions. The platypus's electroreception is the most sensitive of any monotreme.[32][33]
|
34 |
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+
The electroreceptors are located in rostrocaudal rows in the skin of the bill, while mechanoreceptors (which detect touch) are uniformly distributed across the bill. The electrosensory area of the cerebral cortex is contained within the tactile somatosensory area, and some cortical cells receive input from both electroreceptors and mechanoreceptors, suggesting a close association between the tactile and electric senses. Both electroreceptors and mechanoreceptors in the bill dominate the somatotopic map of the platypus brain, in the same way human hands dominate the Penfield homunculus map.[34][35]
|
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+
|
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+
The platypus can determine the direction of an electric source, perhaps by comparing differences in signal strength across the sheet of electroreceptors. This would explain the characteristic side-to-side motion of the animal's head while hunting. The cortical convergence of electrosensory and tactile inputs suggests a mechanism that determines the distance of prey that, when they move, emit both electrical signals and mechanical pressure pulses. The platypus uses the difference between arrival times of the two signals to sense distance.[33]
|
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+
|
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Feeding by neither sight nor smell,[36] the platypus closes its eyes, ears, and nose each time it dives.[37] Rather, when it digs in the bottom of streams with its bill, its electroreceptors detect tiny electric currents generated by muscular contractions of its prey, so enabling it to distinguish between animate and inanimate objects, which continuously stimulate its mechanoreceptors.[33] Experiments have shown the platypus will even react to an "artificial shrimp" if a small electric current is passed through it.[38]
|
40 |
+
|
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+
Monotreme electrolocation probably evolved in order to allow the animals to forage in murky waters, and may be tied to their tooth loss.[39] The extinct Obdurodon was electroreceptive, but unlike the modern platypus it foraged pelagically (near the ocean surface).[39]
|
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+
In recent studies it has been suggested that the eyes of the platypus are more similar to those of Pacific hagfish or Northern Hemisphere lampreys than to those of most tetrapods. The eyes also contain double cones, which most mammals do not have.[40]
|
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|
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+
Although the platypus's eyes are small and not used under water, several features indicate that vision played an important role in its ancestors. The corneal surface and the adjacent surface of the lens is flat while the posterior surface of the lens is steeply curved, similar to the eyes of other aquatic mammals such as otters and sea-lions. A temporal (ear side) concentration of retinal ganglion cells, important for binocular vision, indicates a role in predation, while the accompanying visual acuity is insufficient for such activities. Furthermore, this limited acuity is matched by a low cortical magnification, a small lateral geniculate nucleus and a large optic tectum, suggesting that the visual midbrain plays a more important role than the visual cortex, as in some rodents. These features suggest that the platypus has adapted to an aquatic and nocturnal lifestyle, developing its electrosensory system at the cost of its visual system; an evolutionary process paralleled by the small number of electroreceptors in the short-beaked echidna, which dwells in dry environments, whilst the long-beaked echidna, which lives in moist environments, is intermediate between the other two monotremes.[34]
|
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|
47 |
+
The platypus is semiaquatic, inhabiting small streams and rivers over an extensive range from the cold highlands of Tasmania and the Australian Alps to the tropical rainforests of coastal Queensland as far north as the base of the Cape York Peninsula.[41] Inland, its distribution is not well known; it is extinct in South Australia (apart from an introduced population on Kangaroo Island)[42] and is no longer found in the main part of the Murray-Darling Basin, possibly due to the declining water quality brought about by extensive land clearing and irrigation schemes.[43] Along the coastal river systems, its distribution is unpredictable; it appears to be absent from some relatively healthy rivers, and yet maintains a presence in others, for example, the lower Maribyrnong, that are quite degraded.[44]
|
48 |
+
|
49 |
+
In captivity, platypuses have survived to 17 years of age, and wild specimens have been recaptured when 11 years old. Mortality rates for adults in the wild appear to be low.[11] Natural predators include snakes, water rats, goannas, hawks, owls, and eagles. Low platypus numbers in northern Australia are possibly due to predation by crocodiles.[45] The introduction of red foxes in 1845 for hunting may have had some impact on its numbers on the mainland.[16] The platypus is generally regarded as nocturnal and crepuscular, but individuals are also active during the day, particularly when the sky is overcast.[46][47] Its habitat bridges rivers and the riparian zone for both a food supply of prey species, and banks where it can dig resting and nesting burrows.[47] It may have a range of up to 7 km (4.3 mi), with a male's home range overlapping those of three or four females.[48]
|
50 |
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|
51 |
+
The platypus is an excellent swimmer and spends much of its time in the water foraging for food. It has a very characteristic swimming style and no external ears.[49] Uniquely among mammals, it propels itself when swimming by an alternate rowing motion of the front feet; although all four feet of the platypus are webbed, the hind feet (which are held against the body) do not assist in propulsion, but are used for steering in combination with the tail.[50] The species is endothermic, maintaining its body temperature at about 32 °C (90 °F), lower than most mammals, even while foraging for hours in water below 5 °C (41 °F).[11]
|
52 |
+
|
53 |
+
Dives normally last around 30 seconds, but can last longer, although few exceed the estimated aerobic limit of 40 seconds. Recovery at the surface between dives commonly takes from 10 to 20 seconds.[51][52]
|
54 |
+
|
55 |
+
When not in the water, the platypus retires to a short, straight resting burrow of oval cross-section, nearly always in the riverbank not far above water level, and often hidden under a protective tangle of roots.[49]
|
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+
|
57 |
+
The average sleep time of a platypus is said to be as long as 14 hours per day, possibly because it eats crustaceans, which provide a high level of calories.[53]
|
58 |
+
|
59 |
+
The platypus is a carnivore: it feeds on annelid worms, insect larvae, freshwater shrimp, and freshwater yabby (crayfish) that it digs out of the riverbed with its snout or catches while swimming. It uses cheek-pouches to carry prey to the surface, where it is eaten.[49] The platypus needs to eat about 20% of its own weight each day, which requires it to spend an average of 12 hours daily looking for food.[51]
|
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|
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+
When the platypus was first encountered by European naturalists, they were divided over whether the female laid eggs. This was not confirmed until 1884, when William Hay Caldwell was sent to Australia, where, after extensive searching assisted by a team of 150 Aborigines, he managed to discover a few eggs.[11][26] Mindful of the high cost per word, Caldwell tersely wired London, "Monotremes oviparous, ovum meroblastic." That is, monotremes lay eggs, and the eggs are similar to those of reptiles in that only part of the egg divides as it develops.
|
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|
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+
The species exhibits a single breeding season; mating occurs between June and October, with some local variation taking place between different populations across its range.[45] Historical observation, mark-and-recapture studies, and preliminary investigations of population genetics indicate the possibility of both resident and transient members of populations, and suggest a polygynous mating system.[54] Females are thought likely to become sexually mature in their second year, with breeding confirmed still to take place in animals over nine years old.[54]
|
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+
|
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Outside the mating season, the platypus lives in a simple ground burrow, the entrance of which is about 30 cm (12 in) above the water level. After mating, the female constructs a deeper, more elaborate burrow up to 20 m (65 ft) long and blocked at intervals with plugs (which may act as a safeguard against rising waters or predators, or as a method of regulating humidity and temperature).[55] The male takes no part in caring for its young, and retreats to his year-long burrow. The female softens the ground in the burrow with dead, folded, wet leaves, and she fills the nest at the end of the tunnel with fallen leaves and reeds for bedding material. This material is dragged to the nest by tucking it underneath her curled tail.[6]
|
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|
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The female platypus has a pair of ovaries, but only the left one is functional.[46] The platypus's genes are a possible evolutionary link between the mammalian XY and bird/reptile ZW sex-determination systems because one of the platypus's five X chromosomes contains the DMRT1 gene, which birds possess on their Z chromosome.[56] It lays one to three (usually two) small, leathery eggs (similar to those of reptiles), about 11 mm (7⁄16 in) in diameter and slightly rounder than bird eggs.[57] The eggs develop in utero for about 28 days, with only about 10 days of external incubation (in contrast to a chicken egg, which spends about one day in tract and 21 days externally).[46] After laying her eggs, the female curls around them. The incubation period is divided into three phases.[58] In the first phase, the embryo has no functional organs and relies on the yolk sac for sustenance. The yolk is absorbed by the developing young.[59] During the second phase, the digits develop, and in the last phase, the egg tooth appears.[58]
|
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|
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Most mammal zygotes go through holoblastic cleavage, meaning that, following fertilisation, the ovum is split due to cell divisions into multiple, divisible daughter cells. This is in comparison to the more ancestral process of meroblastic cleavage, present in monotremes like the platypus and in non-mammals like reptiles and birds. In meroblastic cleavage, the ovum does not split completely. This causes the cells at the edge of the yolk to be cytoplasmically continuous with the egg's cytoplasm. This allows the yolk, which contains the embryo, to exchange waste and nutrients with the cytoplasm.[60]
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|
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The newly hatched young are vulnerable, blind, and hairless, and are fed by the mother's milk. Although possessing mammary glands, the platypus lacks teats. Instead, milk is released through pores in the skin. The milk pools in grooves on her abdomen, allowing the young to lap it up.[6][45] After they hatch, the offspring are suckled for three to four months. During incubation and weaning, the mother initially leaves the burrow only for short periods, to forage. When doing so, she creates a number of thin soil plugs along the length of the burrow, possibly to protect the young from predators; pushing past these on her return forces water from her fur and allows the burrow to remain dry.[61] After about five weeks, the mother begins to spend more time away from her young and, at around four months, the young emerge from the burrow.[45] A platypus is born with teeth, but these drop out at a very early age, leaving the horny plates it uses to grind food.[62]
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Platypus
|
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Echidnas
|
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Marsupials
|
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Eutherians
|
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The platypus and other monotremes were very poorly understood, and some of the 19th century myths that grew up around them—for example, that the monotremes were "inferior" or quasireptilian—still endure.[64] In 1947, William King Gregory theorised that placental mammals and marsupials may have diverged earlier, and a subsequent branching divided the monotremes and marsupials, but later research and fossil discoveries have suggested this is incorrect.[64][65] In fact, modern monotremes are the survivors of an early branching of the mammal tree, and a later branching is thought to have led to the marsupial and placental groups.[64][66] Molecular clock and fossil dating suggest platypuses split from echidnas around 19–48 million years ago.[67]
|
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|
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The oldest discovered fossil of the modern platypus dates back to about 100,000 years ago, during the Quaternary period. The extinct monotremes Teinolophos and Steropodon were once thought to be closely related to the modern platypus,[65] but are now considered more basal taxa.[68] The fossilised Steropodon was discovered in New South Wales and is composed of an opalised lower jawbone with three molar teeth (whereas the adult contemporary platypus is toothless). The molar teeth were initially thought to be tribosphenic, which would have supported a variation of Gregory's theory, but later research has suggested, while they have three cusps, they evolved under a separate process.[69] The fossil is thought to be about 110 million years old, making it the oldest mammal fossil found in Australia. Unlike the modern platypus (and echidnas), Teinolophos lacked a beak.[68]
|
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Monotrematum sudamericanum, another fossil relative of the platypus, has been found in Argentina, indicating monotremes were present in the supercontinent of Gondwana when the continents of South America and Australia were joined via Antarctica (up to about 167 million years ago).[69][70] A fossilised tooth of a giant platypus species, Obdurodon tharalkooschild, was dated 5–15 million years ago. Judging by the tooth, the animal measured 1.3 meters long, making it the largest platypus on record.[71]
|
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Because of the early divergence from the therian mammals and the low numbers of extant monotreme species, the platypus is a frequent subject of research in evolutionary biology. In 2004, researchers at the Australian National University discovered the platypus has ten sex chromosomes, compared with two (XY) in most other mammals. These ten chromosomes form five unique pairs of XY in males and XX in females, i.e. males are X1Y1X2Y2X3Y3X4Y4X5Y5.[72] One of the X chromosomes of the platypus has great homology to the bird Z chromosome.[73] The platypus genome also has both reptilian and mammalian genes associated with egg fertilisation.[36][74] Though the platypus lacks the mammalian sex-determining gene SRY, a study found that the mechanism of sex determination is the AMH gene on the oldest Y chromosome.[75][76] A draft version of the platypus genome sequence was published in Nature on 8 May 2008, revealing both reptilian and mammalian elements, as well as two genes found previously only in birds, amphibians, and fish. More than 80% of the platypus's genes are common to the other mammals whose genomes have been sequenced.[36]
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Except for its loss from the state of South Australia, the platypus occupies the same general distribution as it did prior to European settlement of Australia. However, local changes and fragmentation of distribution due to human modification of its habitat are documented. Its current and historical abundance, however, are less well-known and it has probably declined in numbers, although still being considered as common over most of its current range.[47] The species was extensively hunted for its fur until the early years of the 20th century and, although protected throughout Australia since 1905,[61] until about 1950 it was still at risk of drowning in the nets of inland fisheries.[43] The platypus does not appear to be in immediate danger of extinction, because conservation measures have been successful, but it could be affected by habitat disruption caused by dams, irrigation, pollution, netting, and trapping. Reduction of watercourse flows and water levels through excessive droughts and extraction of water for industrial, agricultural, and domestic supplies are also considered a threat. The IUCN lists the platypus on its Red List as "Near Threatened".[2] More recently, in January 2020, researchers from the University of New South Wales presented evidence that the platypus is at risk of extinction, due to a combination of water resource development, land clearing, climate change and increasingly severe periods of drought.[77]
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Platypuses generally suffer from few diseases in the wild; however, public concern in Tasmania is widespread about the potential impacts of a disease caused by the fungus Mucor amphibiorum. The disease (termed mucormycosis) affects only Tasmanian platypuses, and has not been observed in platypuses in mainland Australia. Affected platypuses can develop skin lesions or ulcers on various parts of their bodies, including their backs, tails, and legs. Mucormycosis can kill platypuses, death arising from secondary infection and by affecting the animals' ability to maintain body temperature and forage efficiently. The Biodiversity Conservation Branch at the Department of Primary Industries and Water are collaborating with NRM north and University of Tasmania researchers to determine the impacts of the disease on Tasmanian platypuses, as well as the mechanism of transmission and current spread of the disease.[78]
|
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Much of the world was introduced to the platypus in 1939 when National Geographic Magazine published an article on the platypus and the efforts to study and raise it in captivity. The latter is a difficult task, and only a few young have been successfully raised since, notably at Healesville Sanctuary in Victoria. The leading figure in these efforts was David Fleay, who established a platypusary—a simulated stream in a tank—at the Healesville Sanctuary, where breeding was successful in 1943.[79] In 1972, he found a dead baby of about 50 days old, which had presumably been born in captivity, at his wildlife park at Burleigh Heads on the Gold Coast, Queensland.[80] Healesville repeated its success in 1998 and again in 2000 with a similar stream tank.[81] Since 2008, platypus has bred regularly at Healesville,[82] including second-generation (captive born themselves breeding in captivity).[83] Taronga Zoo in Sydney bred twins in 2003, and breeding was again successful there in 2006.[81]
|
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|
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+
The platypus is kept, for conservation purposes, in special aquariums at the following Australian wildlife sanctuaries:
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As of 2019, the only platypuses in captivity outside of Australia are in the San Diego Zoo Safari Park in the U.S. state of California.[87][88]
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Three attempts were made to bring the animals to the Bronx Zoo, in 1922, 1947, and 1958; of these, only two of the three animals introduced in 1947 lived longer than eighteen months.[89]
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The platypus has been a subject in the Dreamtime stories of indigenous Australians, who believed the animal was a hybrid of a duck and a water rat.[90]:57–60 According to one story, the major animal groups, the land animals, water animals and birds, all competed for the platypus to join their respective groups, but the platypus ultimately decided to not join any of them, feeling that he did not need to be part of a group to be special.[90]:83–85
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Platypuses have been used several times as mascots: Syd the platypus was one of the three mascots chosen for the Sydney 2000 Olympics along with an echidna and a kookaburra,[91] Expo Oz the platypus was the mascot for World Expo 88, which was held in Brisbane in 1988,[92] and Hexley the platypus is the mascot for the Darwin operating system, the BSD-based core of macOS and other operating systems from Apple Inc.[93]
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The platypus has been featured in songs, such as Green Day's "Platypus (I Hate You)" and Mr. Bungle's "Platypus". It is the subject of a children's poem by Banjo Paterson.
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The platypus has frequently appeared in Australian postage stamps and coins. The earliest appearance is the 9d Australian stamp from 1937. The platypus re-appeared in the 1960–64 Australian Native Animal Series. Souvenir sheet of "from" Laos and Equatorial Guinea has also featured the animal. The platypus has appeared on a 1987 36-cent stamp and an Australian 1996 95-cent stamp. The 2006 Australian Bush Babies stamp series features a A$4.65 stamp of a young platypus. A 5-cent stamp also produced in 2006 features the platypus also. Since the introduction of decimal currency to Australia in 1966, the embossed image of a platypus, designed and sculpted by Stuart Devlin, has appeared on the reverse (tails) side of the 20-cent coin.
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In the animated series Phineas and Ferb, the title characters own a pet platypus, named Perry, who unknown to them, is a secret agent. The choice of a platypus was inspired by media underuse, as well as to exploit the animal's striking appearance.[94] As a character, Perry has been well received by both fans and critics.[95][96]
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9d postage stamp from 1937
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Large carving of a platypus at the Australian Axeman's Hall of Fame
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A platypus fur cape made in 1890. It was donated to the National Gallery of Victoria by Mrs F Smith in 1985
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The orca or killer whale (Orcinus orca) is a toothed whale belonging to the oceanic dolphin family, of which it is the largest member. Killer whales have a diverse diet, although individual populations often specialize in particular types of prey. Some feed exclusively on fish, while others hunt marine mammals such as seals and other species of dolphin. They have been known to attack baleen whale calves, and even adult whales. Killer whales are apex predators, as no animal preys on them. A cosmopolitan species, they can be found in each of the world's oceans in a variety of marine environments, from Arctic and Antarctic regions to tropical seas, absent only from the Baltic and Black seas, and some areas of the Arctic Ocean.
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Killer whales are highly social; some populations are composed of matrilineal family groups (pods) which are the most stable of any animal species. Their sophisticated hunting techniques and vocal behaviours, which are often specific to a particular group and passed across generations, have been described as manifestations of animal culture.
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The International Union for Conservation of Nature assesses the orca's conservation status as data deficient because of the likelihood that two or more killer whale types are separate species. Some local populations are considered threatened or endangered due to prey depletion, habitat loss, pollution (by PCBs), capture for marine mammal parks, and conflicts with human fisheries. In late 2005, the southern resident killer whales, which swim in British Columbia and Washington state waters, were placed on the U.S. Endangered Species list.
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Wild killer whales are not considered a threat to humans and no fatal attack on humans has ever been documented, but there have been cases of captive orcas killing or injuring their handlers at marine theme parks. Killer whales feature strongly in the mythologies of indigenous cultures, and their reputation in different cultures ranges from being the souls of humans to merciless killers.
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Orcinus orca is the only recognized extant species in the genus Orcinus, and one of many animal species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae.[6] Konrad Gessner wrote the first scientific description of a killer whale in his Piscium & aquatilium animantium natura of 1558, part of the larger Historia animalium, based on examination of a dead stranded animal in the Bay of Greifswald that had attracted a great deal of local interest.[7]
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The killer whale is one of 35 species in the oceanic dolphin family, which first appeared about 11 million years ago. The killer whale lineage probably branched off shortly thereafter.[8] Although it has morphological similarities with the false killer whale, the pygmy killer whale and the pilot whales, a study of cytochrome b gene sequences by Richard LeDuc indicated that its closest extant relatives are the snubfin dolphins of the genus Orcaella.[9] However, a more recent (2018) study places the orca as a sister taxon to the Lissodelphininae, a clade that includes Lagenorhynchus and Cephalorhynchus.[10]
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Although the term "orca" is increasingly used, English-speaking scientists most often use the traditional name "killer whale"[citation needed]. The genus name Orcinus means "of the kingdom of the dead",[11] or "belonging to Orcus".[12]
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Ancient Romans originally used orca (pl. orcae) for these animals, possibly borrowing Ancient Greek ὄρυξ (óryx), which referred (among other things) to a whale species. Since the 1960s, "orca" has steadily grown in popularity. The term "orca" is preferred by some as it avoids the negative connotations of "killer",[13] and because, being part of the family Delphinidae, the species is more closely related to other oceanic dolphins than to other whales.[14]
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They are sometimes referred to as "blackfish", a name also used for other whale species. "Grampus" is a former name for the species, but is now seldom used. This meaning of "grampus" should not be confused with the genus Grampus, whose only member is Risso's dolphin.[15]
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The three to five types of killer whales may be distinct enough to be considered different races,[16] subspecies, or possibly even species[17] (see Species problem). The IUCN reported in 2008, "The taxonomy of this genus is clearly in need of review, and it is likely that O. orca will be split into a number of different species or at least subspecies over the next few years."[3] Although large variation in the ecological distinctiveness of different killer whale groups complicate simple differentiation into types,[18] research off the west coast of Canada and the United States in the 1970s and 1980s identified the following three types:
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Transients and residents live in the same areas, but avoid each other.[31][32][33]
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Other populations have not been as well studied, although specialized fish and mammal eating killer whales have been distinguished elsewhere.[34] In addition, separate populations of "generalist" (fish- and mammal-eating) and "specialist" (mammal-eating) killer whales have been identified off northwestern Europe.[35][36] As with residents and transients, the lifestyle of these whales appears to reflect their diet; fish-eating killer whales in Alaska[37] and Norway[38] have resident-like social structures, while mammal-eating killer whales in Argentina and the Crozet Islands behave more like transients.[39]
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Three types have been documented in the Antarctic. Two dwarf species, named Orcinus nanus and Orcinus glacialis, were described during the 1980s by Soviet researchers, but most cetacean researchers are sceptical about their status, and linking these directly to the types described below is difficult.[17]
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Types B and C live close to the ice pack, and diatoms in these waters may be responsible for the yellowish colouring of both types.[17][46] Mitochondrial DNA sequences support the theory that these are recently diverged separate species.[47] More recently, complete mitochondrial sequencing indicates the two Antarctic groups that eat seals and fish should be recognized as distinct species, as should the North Pacific transients, leaving the others as subspecies pending additional data.[48] Advanced methods that sequenced the entire mitochondrial genome revealed systematic differences in DNA between different populations.[49] A 2019 study of Type D orcas also found them to be distinct from other populations and possibly even a unique species.[44]
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Mammal-eating killer whales in different regions were long thought likely to be closely related, but genetic testing has refuted this hypothesis.[50]
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There are seven identified ecotypes inhabiting isolated ecological niches. Of three orca ecotypes in the Antarctic, one preys on minke whales, the second on seals and penguins, and the third on fish. Another ecotype lives in the eastern North Atlantic, while the three Northeast Pacific ecotypes are labelled the transient, resident and offshore populations described above. Research has supported a proposal to reclassify the Antarctic seal- and fish-eating populations and the North Pacific transients as a distinct species, leaving the remaining ecotypes as subspecies. The first split in the orca population, between the North Pacific transients and the rest, occurred an estimated 700,000 years ago. Such a designation would mean that each new species becomes subject to separate conservation assessments.[49]
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A typical killer whale distinctively bears a black back, white chest and sides, and a white patch above and behind the eye. Calves are born with a yellowish or orange tint, which fades to white. It has a heavy and robust body[51] with a large dorsal fin up to 1.8 m (5 ft 11 in) tall.[52] Behind the fin, it has a dark grey "saddle patch" across the back. Antarctic killer whales may have pale grey to nearly white backs. Adult killer whales are very distinctive, seldom confused with any other sea creature.[53] When seen from a distance, juveniles can be confused with other cetacean species, such as the false killer whale or Risso's dolphin.[54]
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The killer whale's teeth are very strong, and its jaws exert a powerful grip; the upper teeth fall into the gaps between the lower teeth when the mouth is closed. The firm middle and back teeth hold prey in place, while the front teeth are inclined slightly forward and outward to protect them from powerful jerking movements.[55]
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Killer whales are the largest extant members of the dolphin family. Males typically range from 6 to 8 metres (20 to 26 ft) long and weigh in excess of 6 tonnes (5.9 long tons; 6.6 short tons). Females are smaller, generally ranging from 5 to 7 m (16 to 23 ft) and weighing about 3 to 4 tonnes (3.0 to 3.9 long tons; 3.3 to 4.4 short tons).[56] Calves at birth weigh about 180 kg (400 lb) and are about 2.4 m (7.9 ft) long.[57][58] The skeleton of the killer whale is of the typical delphinid structure, but more robust.[59] Its integument, unlike that of most other dolphin species, is characterized by a well-developed dermal layer with a dense network of fascicles of collagen fibres.[60]
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Killer whale pectoral fins, analogous to forelimbs, are large and rounded, resembling paddles, with those of males significantly larger than those of females. Dorsal fins also exhibit sexual dimorphism, with those of males about 1.8 m (5.9 ft) high, more than twice the size of the female's, with the male's fin more like a tall, elongated isosceles triangle, whereas the female's is shorter and more curved.[61] Males and females also have different patterns of black and white skin in their genital areas.[62] In the skull, adult males have longer lower jaws than females, as well as larger occipital crests.[60]
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An individual killer whale can often be identified from its dorsal fin and saddle patch. Variations such as nicks, scratches, and tears on the dorsal fin and the pattern of white or grey in the saddle patch are unique. Published directories contain identifying photographs and names for hundreds of North Pacific animals. Photographic identification has enabled the local population of killer whales to be counted each year rather than estimated, and has enabled great insight into life cycles and social structures.[63]
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Occasionally a killer whale is white; they have been spotted in the northern Bering Sea and around St. Lawrence Island, and near the Russian coast.[64][65] In February 2008, a white killer whale was photographed 3.2 km (2.0 mi) off Kanaga Volcano in the Aleutian Islands.[64][65] In 2010, the Far East Russia Orca Project (FEROP), co-founded and co-directed by Alexander M. Burdin and Erich Hoyt, filmed an adult male nicknamed Iceberg.[66][67]
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Killer whales have good eyesight above and below the water, excellent hearing, and a good sense of touch. They have exceptionally sophisticated echolocation abilities, detecting the location and characteristics of prey and other objects in the water by emitting clicks and listening for echoes,[68] as do other members of the dolphin family. The mean body temperature of the orca is 36 to 38 °C (97 to 100 °F).[69][70] Like most marine mammals, orcas have a layer of insulating blubber ranging from 7.6 to 10 cm (3.0 to 3.9 in) thick[69] beneath the skin. The pulse is about 60 heartbeats per minute when the orca is at the surface, dropping to 30 beats/min when submerged.[71]
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Killer whales are found in all oceans and most seas. Due to their enormous range, numbers, and density, relative distribution is difficult to estimate,[72] but they clearly prefer higher latitudes and coastal areas over pelagic environments.[73] Areas which serve as major study sites for the species include the coasts of Iceland, Norway, the Valdes Peninsula of Argentina, the Crozet Islands, New Zealand and parts of the west coast of North America, from California to Alaska.[74]
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Systematic surveys indicate the highest densities of killer whales (>0.40 individuals per 100 km2) in the northeast Atlantic around the Norwegian coast, in the north Pacific along the Aleutian Islands, the Gulf of Alaska and in the Southern Ocean off much of the coast of Antarctica.[72] They are considered "common" (0.20–0.40 individuals per 100 km2) in the eastern Pacific along the coasts of British Columbia, Washington and Oregon, in the North Atlantic Ocean around Iceland and the Faroe Islands. High densities have also been reported but not quantified in the western North Pacific around the Sea of Japan, Sea of Okhotsk, Kuril Islands, Kamchatka and the Commander Islands and in the Southern Hemisphere off southern Brazil and the tip of southern Africa. They are reported as seasonally common in the Canadian Arctic, including Baffin Bay between Greenland and Nunavut, as well as Tasmania and Macquarie Island.[72] Regularly occurring or distinct populations exist off Northwest Europe, California, Patagonia, the Crozet Islands, Marion Island, southern Australia and New Zealand.[36][72][75] The northwest Atlantic population of at least 67 individuals ranges from Labrador and Newfoundland to New England with sightings to Cape Cod and Long Island.[76]
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Information for offshore regions and warmer waters is more scarce, but widespread sightings indicate that the killer whale can survive in most water temperatures. They have been sighted, though more infrequently, in the Mediterranean, the Arabian Sea, the Gulf of Mexico, Banderas Bay on Mexico's west coast and the Caribbean.[72] Over 50 individual whales have been documented in the northern Indian Ocean, including two individuals that were sighted in the Persian Gulf in 2008 and off Sri Lanka in 2015.[77] Those orcas may occasionally enter the Red Sea through the Gulf of Aden.[78] The modern status of the species along coastal mainland China and its vicinity is unknown. Recorded sightings have been made from almost the entire shoreline.[79] A wide-ranging population is likely to exist in the central Pacific, with some sightings off Hawaii.[80][81] Distinct populations may also exist off the west coast of tropical Africa,[82] and Papua New Guinea.[83] In the Mediterranean, killer whales are considered "visitors", likely from the North Atlantic, and sightings become less frequent further east. However, a small year-round population is known to exist in the Strait of Gibraltar, mostly on the Atlantic side.[84][85] Killer whales also appear to regularly occur off the Galápagos Islands.[86]
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In the Antarctic, killer whales range up to the edge of the pack ice and are believed to venture into the denser pack ice, finding open leads much like beluga whales in the Arctic. However, killer whales are merely seasonal visitors to Arctic waters, and do not approach the pack ice in the summer. With the rapid Arctic sea ice decline in the Hudson Strait, their range now extends deep into the northwest Atlantic.[87] Occasionally, killer whales swim into freshwater rivers. They have been documented 100 mi (160 km) up the Columbia River in the United States.[88][89] They have also been found in the Fraser River in Canada and the Horikawa River in Japan.[88]
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Migration patterns are poorly understood. Each summer, the same individuals appear off the coasts of British Columbia and Washington. Despite decades of research, where these animals go for the rest of the year remains unknown. Transient pods have been sighted from southern Alaska to central California.[90]
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Worldwide population estimates are uncertain, but recent consensus suggests a minimum of 50,000 (2006).[91][3][30] Local estimates include roughly 25,000 in the Antarctic, 8,500 in the tropical Pacific, 2,250–2,700 off the cooler northeast Pacific and 500–1,500 off Norway.[92] Japan's Fisheries Agency estimated in the 2000s that 2,321 killer whales were in the seas around Japan.[93][94]
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Killer whales are apex predators, meaning that they themselves have no natural predators. They are sometimes called the wolves of the sea, because they hunt in groups like wolf packs.[95] Killer whales hunt varied prey including fish, cephalopods, mammals, seabirds, and sea turtles.[96] Different populations or ecotypes may specialize, and some can have a dramatic impact on prey species.[97] However, whales in tropical areas appear to have more generalized diets due to lower food productivity.[81][82] Killer whales spend most of their time at shallow depths,[98] but occasionally dive several hundred meters depending on their prey.[99][100]
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Fish-eating killer whales prey on around 30 species of fish. Some populations in the Norwegian and Greenland sea specialize in herring and follow that fish's autumnal migration to the Norwegian coast. Salmon account for 96% of northeast Pacific residents' diet, including 65% of large, fatty Chinook.[101] Chum salmon are also eaten, but smaller sockeye and pink salmon are not a significant food item.[102] Depletion of specific prey species in an area is, therefore, cause for concern for local populations, despite the high diversity of prey. On average, a killer whale eats 227 kilograms (500 lb) each day.[103] While salmon are usually hunted by an individual whale or a small group, herring are often caught using carousel feeding: the killer whales force the herring into a tight ball by releasing bursts of bubbles or flashing their white undersides. They then slap the ball with their tail flukes, stunning or killing up to 15 fish at a time, then eating them one by one. Carousel feeding has only been documented in the Norwegian killer whale population, as well as some oceanic dolphin species.[104]
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In New Zealand, sharks and rays appear to be important prey, including eagle rays, long-tail and short-tail stingrays, common threshers, smooth hammerheads, blue sharks, basking sharks, and shortfin makos.[105][106] With sharks, orcas may herd them to the surface and strike them with their tail flukes,[105] while bottom-dwelling rays are cornered, pinned to the ground and taken to the surface.[107] In other parts of the world, killer whales have preyed on broadnose sevengill sharks,[108] tiger sharks[109] and even small whale sharks.[110] Killer whales have also been recorded attacking and feeding on great white sharks,[28][111][112][113] and appear to target the liver.[28][112] Competition between killer whales and white sharks is probable in regions where their diets overlap.[114] The arrival of orcas in an area can cause white sharks to flee and forage elsewhere.[115]
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Killer whales are very sophisticated and effective predators of marine mammals. Thirty-two cetacean species have been recorded as prey, from observing orcas' feeding activity, examining the stomach contents of dead orcas, and seeing scars on the bodies of surviving prey animals. Groups even attack larger cetaceans such as minke whales, grey whales,[116] and, rarely, sperm whales or blue whales.[34][117][118][119] It has been hypothesized that predation by orcas on whale calves in high-productivity, high-latitude areas is the reason for great whale migrations during breeding season to low-productivity tropical waters where orcas are scarcer.[120]
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Hunting a large whale usually takes several hours. Killer whales generally attack young or weak animals; however, a group of five or more may attack a healthy adult. When hunting a young whale, a group chases it and its mother to exhaustion. Eventually, they separate the pair and surround the calf, drowning it by keeping it from surfacing. Pods of female sperm whales sometimes protect themselves by forming a protective circle around their calves with their flukes facing outwards, using them to repel the attackers.[121] Rarely, large killer whale pods can overwhelm even adult female sperm whales. Adult bull sperm whales, which are large, powerful and aggressive when threatened, and fully grown adult blue whales, which are possibly too large to overwhelm, are not believed to be prey for killer whales.[122]
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Prior to the advent of industrial whaling, great whales may have been the major food source for killer whales. The introduction of modern whaling techniques may have aided killer whales by the sound of exploding harpoons indicating availability of prey to scavenge, and compressed air inflation of whale carcasses causing them to float, thus exposing them to scavenging. However, the devastation of great whale populations by unfettered whaling has possibly reduced their availability for killer whales, and caused them to expand their consumption of smaller marine mammals, thus contributing to the decline of these as well.[120]
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Other marine mammal prey species include nearly 20 species of seal, sea lion and fur seal. Walruses and sea otters are less frequently taken. Often, to avoid injury, killer whales disable their prey before killing and eating it. This may involve throwing it in the air, slapping it with their tails, ramming it, or breaching and landing on it.[123] In the Aleutian Islands, a decline in sea otter populations in the 1990s was controversially attributed by some scientists to killer whale predation, although with no direct evidence.[124] The decline of sea otters followed a decline in harbour seal and Steller sea lion populations, the killer whale's preferred prey,[a][126] which in turn may be substitutes for their original prey, now decimated by industrial whaling.[127][128][129]
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In steeply banked beaches off Península Valdés, Argentina, and the Crozet Islands, killer whales feed on South American sea lions and southern elephant seals in shallow water, even beaching temporarily to grab prey before wriggling back to the sea. Beaching, usually fatal to cetaceans, is not an instinctive behaviour, and can require years of practice for the young.[130] Killer whales can then release the animal near juvenile whales, allowing the younger whales to practice the difficult capture technique on the now-weakened prey.[123][131] "Wave-hunting" killer whales spy-hop to locate Weddell seals, crabeater seals, leopard seals, and penguins resting on ice floes, and then swim in groups to create waves that wash over the floe. This washes the prey into the water, where other killer whales lie in wait.[49][132][133]
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Killer whales have also been observed preying on terrestrial mammals, such as deer swimming between islands off the northwest coast of North America.[125] Killer whale cannibalism has also been reported based on analysis of stomach contents, but this is likely to be the result of scavenging remains dumped by whalers.[134] One killer whale was also attacked by its companions after being shot.[34] Although resident killer whales have never been observed to eat other marine mammals, they occasionally harass and kill porpoises and seals for no apparent reason.[135]
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Killer whales in many areas may prey on cormorants and gulls.[136] A captive killer whale at Marineland of Canada discovered it could regurgitate fish onto the surface, attracting sea gulls, and then eat the birds. Four others then learned to copy the behaviour.[137]
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Day-to-day killer whale behaviour generally consists of foraging, travelling, resting and socializing. Killer whales frequently engage in surface behaviour such as breaching (jumping completely out of the water) and tail-slapping. These activities may have a variety of purposes, such as courtship, communication, dislodging parasites, or play. Spyhopping is a behaviour in which a whale holds its head above water to view its surroundings.[138] Resident killer whales swim alongside porpoises and other dolphins.[139]
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Killer whales are notable for their complex societies. Only elephants and higher primates live in comparably complex social structures.[140] Due to orcas' complex social bonds, many marine experts have concerns about how humane it is to keep them in captivity.[141]
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Resident killer whales in the eastern North Pacific live in particularly complex and stable social groups. Unlike any other known mammal social structure, resident whales live with their mothers for their entire lives. These family groups are based on matrilines consisting of the eldest female (matriarch) and her sons and daughters, and the descendants of her daughters, etc. The average size of a matriline is 5.5 animals.[142] Because females can reach age 90, as many as four generations travel together. These matrilineal groups are highly stable. Individuals separate for only a few hours at a time, to mate or forage. With one exception, a killer whale named Luna, no permanent separation of an individual from a resident matriline has been recorded.[142]
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Closely related matrilines form loose aggregations called pods, usually consisting of one to four matrilines. Unlike matrilines, pods may separate for weeks or months at a time.[142] DNA testing indicates resident males nearly always mate with females from other pods.[143] Clans, the next level of resident social structure, are composed of pods with similar dialects, and common but older maternal heritage. Clan ranges overlap, mingling pods from different clans.[142] The final association layer, perhaps more arbitrarily defined than the familial groupings, is called the community, and is defined as a set of clans that regularly commingle. Clans within a community do not share vocal patterns.[b]
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Transient pods are smaller than resident pods, typically consisting of an adult female and one or two of her offspring. Males typically maintain stronger relationships with their mothers than other females. These bonds can extend well into adulthood. Unlike residents, extended or permanent separation of transient offspring from natal matrilines is common, with juveniles and adults of both sexes participating. Some males become "rovers" and do not form long-term associations, occasionally joining groups that contain reproductive females.[144] As in resident clans, transient community members share an acoustic repertoire, although regional differences in vocalizations have been noted.[145]
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Like all cetaceans, killer whales depend heavily on underwater sound for orientation, feeding, and communication. They produce three categories of sounds: clicks, whistles, and pulsed calls. Clicks are believed to be used primarily for navigation and discriminating prey and other objects in the surrounding environment, but are also commonly heard during social interactions.[30]
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Northeast Pacific resident groups tend to be much more vocal than transient groups in the same waters.[146] Residents feed primarily on Chinook and chum, which are insensitive to killer whale calls (inferred from the audiogram of Atlantic salmon). In contrast, the marine mammal prey of transients hear whale calls well. Transients are typically silent.[146] They sometimes use a single click (called a cryptic click) rather than the long train of clicks observed in other populations. Residents are silent only when resting.
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All members of a resident pod use similar calls, known collectively as a dialect. Dialects are composed of specific numbers and types of discrete, repetitive calls. They are complex and stable over time.[147] Call patterns and structure are distinctive within matrilines.[148] Newborns produce calls similar to their mothers, but have a more limited repertoire.[145] Individuals likely learn their dialect through contact with pod members.[149] Family-specific calls have been observed more frequently in the days following a calf's birth, which may help the calf learn them.[150] Dialects are probably an important means of maintaining group identity and cohesiveness. Similarity in dialects likely reflects the degree of relatedness between pods, with variation growing over time.[151] When pods meet, dominant call types decrease and subset call types increase. The use of both call types is called biphonation. The increased subset call types may be the distinguishing factor between pods and inter-pod relations.[148]
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Dialects also distinguish types. Resident dialects contain seven to 17 (mean = 11) distinctive call types. All members of the North American west coast transient community express the same basic dialect, although minor regional variation in call types is evident. Preliminary research indicates offshore killer whales have group-specific dialects unlike those of residents and transients.[151]
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Norwegian and Icelandic herring-eating orcas appear to have different vocalizations for activities like hunting.[152] A population that live in McMurdo Sound, Antarctica have 28 complex burst-pulse and whistle calls.[153]
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Killer whales have the second-heaviest brains among marine mammals[154] (after sperm whales, which have the largest brain of any animal). They can be trained in captivity and are often described as intelligent,[155][156] although defining and measuring "intelligence" is difficult in a species whose environment and behavioural strategies are very different from those of humans.[156]
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Killer whales imitate others, and seem to deliberately teach skills to their kin. Off the Crozet Islands, mothers push their calves onto the beach, waiting to pull the youngster back if needed.[123][131]
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People who have interacted closely with killer whales offer numerous anecdotes demonstrating the whales' curiosity, playfulness, and ability to solve problems. Alaskan killer whales have not only learned how to steal fish from longlines, but have also overcome a variety of techniques designed to stop them, such as the use of unbaited lines as decoys.[157] Once, fishermen placed their boats several miles apart, taking turns retrieving small amounts of their catch, in the hope that the whales would not have enough time to move between boats to steal the catch as it was being retrieved. A researcher described what happened next:
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It worked really well for a while. Then the whales split into two groups. It didn't even take them an hour to figure it out. They were so thrilled when they figured out what was going on, that we were playing games. They were breaching by the boats.
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In other anecdotes, researchers describe incidents in which wild killer whales playfully tease humans by repeatedly moving objects the humans are trying to reach,[158] or suddenly start to toss around a chunk of ice after a human throws a snowball.[159]
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The killer whale's use of dialects and the passing of other learned behaviours from generation to generation have been described as a form of animal culture.[160]
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The complex and stable vocal and behavioural cultures of sympatric groups of killer whales (Orcinus orca) appear to have no parallel outside humans and represent an independent evolution of cultural faculties.[161]
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(Two species or populations are considered sympatric when they live in the same geographic area and thus regularly encounter one another.)
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Female killer whales begin to mature at around the age of 10 and reach peak fertility around 20,[162] experiencing periods of polyestrous cycling separated by non-cycling periods of three to 16 months. Females can often breed until age 40, followed by a rapid decrease in fertility.[162] As such, orcas are among the few animals that undergo menopause and live for decades after they have finished breeding.[163][164] The lifespans of wild females average 50 to 80 years.[165] Some are claimed to have lived substantially longer: Granny (J2) was estimated by some researchers to have been as old as 105 years at the time of her death, though a biopsy sample indicated her age as 65 to 80 years.[166][167][168] Orcas held in captivity tend to live less than those in the wild, although this is subject to scientific debate.[165][169][170]
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To avoid inbreeding, males mate with females from other pods. Gestation varies from 15 to 18 months. [171] Mothers usually calve a single offspring about once every five years. In resident pods, births occur at any time of year, although winter is the most common. Mortality is extremely high during the first seven months of life, when 37–50% of all calves die.[172] Weaning begins at about 12 months of age, and is complete by two years. According to observations in several regions, all male and female pod members participate in the care of the young.[140]
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Males sexually mature at the age of 15, but do not typically reproduce until age 21. Wild males live around 29 years on average, with a maximum of about 60 years.[166] One male, known as Old Tom, was reportedly spotted every winter between the 1840s and 1930 off New South Wales, Australia. This would have made him up to 90 years old. Examination of his teeth indicated he died around age 35,[173] but this method of age determination is now believed to be inaccurate for older animals.[174] One male known to researchers in the Pacific Northwest (identified as J1) was estimated to have been 59 years old when he died in 2010.[175] Killer whales are unique among cetaceans, as their caudal sections elongate with age, making their heads relatively shorter.[60]
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Infanticide, once thought to occur only in captive killer whales, was observed in wild populations by researchers off British Columbia on December 2, 2016. In this incident, an adult male killed the calf of a female within the same pod, with his mother also joining in the assault. It is theorized that the male killed the young calf in order to mate with its mother (something that occurs in other carnivore species), while the male's mother supported the breeding opportunity for her son. The attack ended when the calf's mother struck and injured the attacking male. Such behaviour matches that of many smaller dolphin species, such as the bottlenose dolphin.[176]
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In 2008, the IUCN (International Union for Conservation of Nature) changed its assessment of the killer whale's conservation status from conservation dependent to data deficient, recognizing that one or more killer whale types may actually be separate, endangered species.[3] Depletion of prey species, pollution, large-scale oil spills, and habitat disturbance caused by noise and conflicts with boats are the most significant worldwide threats.[3] In January 2020, the first killer whale in England and Wales since 2001 was found dead with a large fragment of plastic in its stomach.[177]
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Like other animals at the highest trophic levels, the killer whale is particularly at risk of poisoning from bioaccumulation of toxins, including Polychlorinated biphenyls (PCBs).[178] European harbour seals have problems in reproductive and immune functions associated with high levels of PCBs and related contaminants, and a survey off the Washington coast found PCB levels in killer whales were higher than levels that had caused health problems in harbour seals.[178] Blubber samples in the Norwegian Arctic show higher levels of PCBs, pesticides and brominated flame-retardants than in polar bears. When food is scarce, killer whales metabolize blubber for energy, which increases pollutant concentrations in their blood.
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In the Pacific Northwest, wild salmon stocks, a main resident food source, have declined dramatically in recent years.[3] In the Puget Sound region only 75 whales remain with few births over the last few years.[179] On the west coast of Alaska and the Aleutian Islands, seal and sea lion populations have also substantially declined.[180]
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In 2005, the United States government listed the southern resident community as an endangered population under the Endangered Species Act.[30] This community comprises three pods which live mostly in the Georgia and Haro Straits and Puget Sound in British Columbia and Washington. They do not breed outside of their community, which was once estimated at around 200 animals and later shrank to around 90.[181] In October 2008, the annual survey revealed seven were missing and presumed dead, reducing the count to 83.[182] This is potentially the largest decline in the population in the past 10 years. These deaths can be attributed to declines in Chinook salmon.[182]
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Scientist Ken Balcomb has extensively studied killer whales since 1976; he is the research biologist responsible for discovering U.S. Navy sonar may harm killer whales. He studied killer whales from the Center for Whale Research, located in Friday Harbor, Washington.[183] He was also able to study killer whales from "his home porch perched above Puget Sound, where the animals hunt and play in summer months".[183] In May 2003, Balcomb (along with other whale watchers near the Puget Sound coastline) noticed uncharacteristic behaviour displayed by the killer whales. The whales seemed "agitated and were moving haphazardly, attempting to lift their heads free of the water" to escape the sound of the sonars.[183] "Balcomb confirmed at the time that strange underwater pinging noises detected with underwater microphones were sonar. The sound originated from a U.S. Navy frigate 12 miles (19 kilometres) distant, Balcomb said."[183] The impact of sonar waves on killer whales is potentially life-threatening. Three years prior to Balcomb's discovery, research in the Bahamas showed 14 beaked whales washed up on the shore. These whales were beached on the day U.S. Navy destroyers were activated into sonar exercise.[183] Of the 14 whales beached, six of them died. These six dead whales were studied, and CAT scans of two of the whale heads showed hemorrhaging around the brain and the ears, which is consistent with decompression sickness.[183]
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Another conservation concern was made public in September 2008 when the Canadian government decided it was not necessary to enforce further protections (including the Species at Risk Act in place to protect endangered animals along their habitats) for killer whales aside from the laws already in place. In response to this decision, six environmental groups sued the federal government, claiming killer whales were facing many threats on the British Columbia Coast and the federal government did nothing to protect them from these threats.[184] A legal and scientific nonprofit organization, Ecojustice, led the lawsuit and represented the David Suzuki Foundation, Environmental Defence, Greenpeace Canada, International Fund for Animal Welfare, the Raincoast Conservation Foundation, and the Wilderness Committee.[184] Many scientists involved in this lawsuit, including Bill Wareham, a marine scientist with the David Suzuki Foundation, noted increased boat traffic, water toxic wastes, and low salmon population as major threats, putting approximately 87 killer whales[184] on the British Columbia Coast in danger.
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Underwater noise from shipping, drilling, and other human activities is a significant concern in some key killer whale habitats, including Johnstone Strait and Haro Strait.[185] In the mid-1990s, loud underwater noises from salmon farms were used to deter seals. Killer whales also avoided the surrounding waters.[186] High-intensity sonar used by the Navy disturbs killer whales along with other marine mammals.[187] Killer whales are popular with whale watchers, which may stress the whales and alter their behaviour, particularly if boats approach too closely or block their lines of travel.[188]
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The Exxon Valdez oil spill adversely affected killer whales in Prince William Sound and Alaska's Kenai Fjords region. Eleven members (about half) of one resident pod disappeared in the following year. The spill damaged salmon and other prey populations, which in turn damaged local killer whales. By 2009, scientists estimated the AT1 transient population (considered part of a larger population of 346 transients), numbered only seven individuals and had not reproduced since the spill. This population is expected to die out.[189][190]
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A 2018 study published in Science found that global killer whale populations are poised to dramatically decline due to exposure to toxic chemical and PCB pollution.[191]
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The indigenous peoples of the Pacific Northwest Coast feature killer whales throughout their art, history, spirituality and religion. The Haida regarded killer whales as the most powerful animals in the ocean, and their mythology tells of killer whales living in houses and towns under the sea. According to these myths, they took on human form when submerged, and humans who drowned went to live with them.[192] For the Kwakwaka'wakw, the killer whale was regarded as the ruler of the undersea world, with sea lions for slaves and dolphins for warriors.[192] In Nuu-chah-nulth and Kwakwaka'wakw mythology, killer whales may embody the souls of deceased chiefs.[192] The Tlingit of southeastern Alaska regarded the killer whale as custodian of the sea and a benefactor of humans.[193]
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The Maritime Archaic people of Newfoundland also had great respect for killer whales, as evidenced by stone carvings found in a 4,000-year-old burial at the Port au Choix Archaeological Site.[194][195]
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In the tales and beliefs of the Siberian Yupik people, killer whales are said to appear as wolves in winter, and wolves as killer whales in summer.[196][197][198][199] Killer whales are believed to assist their hunters in driving walrus.[200] Reverence is expressed in several forms: the boat represents the animal, and a wooden carving hung from the hunter's belt.[198] Small sacrifices such as tobacco or meat are strewn into the sea for them.[200][199]
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Indigenous Ainu tribe often referred killer whales in their folklore and myth as Repun Kamuy (God of Sea/Offshore) to bring fortunes (whales) to the coasts, and there had been traditional funerals for stranded or deceased orcas akin to funerals for other animals such as brown bears.[201]
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In Western cultures, killer whales were historically feared as dangerous, savage predators.[202] The first written description of a killer whale was given by Pliny the Elder circa AD 70, who wrote, "Orcas (the appearance of which no image can express, other than an enormous mass of savage flesh with teeth) are the enemy of [other kinds of whale]... they charge and pierce them like warships ramming."[203]
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Of the very few confirmed attacks on humans by wild killer whales, none have been fatal.[204] In one instance, killer whales tried to tip ice floes on which a dog team and photographer of the Terra Nova Expedition were standing.[205] The sled dogs' barking is speculated to have sounded enough like seal calls to trigger the killer whale's hunting curiosity. In the 1970s, a surfer in California was bitten, and in 2005, a boy in Alaska who was splashing in a region frequented by harbour seals was bumped by a killer whale that apparently misidentified him as prey.[206] Unlike wild killer whales, captive killer whales have made nearly two dozen attacks on humans since the 1970s, some of which have been fatal.[207][208]
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Competition with fishermen also led to killer whales being regarded as pests. In the waters of the Pacific Northwest and Iceland, the shooting of killer whales was accepted and even encouraged by governments.[202] As an indication of the intensity of shooting that occurred until fairly recently, about 25% of the killer whales captured in Puget Sound for aquarium through 1970 bore bullet scars.[209] The U.S. Navy claimed to have deliberately killed hundreds of killer whales in Icelandic waters in 1956 with machine guns, rockets, and depth charges.[210][211]
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Western attitudes towards killer whales have changed dramatically in recent decades. In the mid-1960s and early 1970s, killer whales came to much greater public and scientific awareness, starting with the first live-capture and display of a killer whale known as Moby Doll, a resident harpooned off Saturna Island in 1964.[202] So little was known at the time, it was nearly two months before the whale's keepers discovered what food (fish) it was willing to eat. To the surprise of those who saw him, Moby Doll was a docile, non-aggressive whale who made no attempts to attack humans.[212]
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Between 1964 and 1976, 50 killer whales from the Pacific Northwest were captured for display in aquaria, and public interest in the animals grew. In the 1970s, research pioneered by Michael Bigg led to the discovery of the species' complex social structure, its use of vocal communication, and its extraordinarily stable mother–offspring bonds. Through photo-identification techniques, individuals were named and tracked over decades.[213]
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Bigg's techniques also revealed the Pacific Northwest population was in the low hundreds rather than the thousands that had been previously assumed.[202] The southern resident community alone had lost 48 of its members to captivity; by 1976, only 80 remained.[214] In the Pacific Northwest, the species that had unthinkingly been targeted became a cultural icon within a few decades.[181]
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The public's growing appreciation also led to growing opposition to whale–keeping in aquarium. Only one whale has been taken in North American waters since 1976. In recent years, the extent of the public's interest in killer whales has manifested itself in several high-profile efforts surrounding individuals. Following the success of the 1993 film Free Willy, the movie's captive star Keiko was returned to the coast of his native Iceland in 2002. The director of the International Marine Mammal Project for the Earth Island Institute, David Phillips, led the efforts to return Keiko to the Iceland waters.[215] Keiko however did not adapt to the harsh climate of the Arctic Ocean, and died a year into his release after contracting pneumonia, at the age of 27.[216] In 2002, the orphan Springer was discovered in Puget Sound, Washington. She became the first whale to be successfully reintegrated into a wild pod after human intervention, crystallizing decades of research into the vocal behaviour and social structure of the region's killer whales.[217] The saving of Springer raised hopes that another young killer whale named Luna, which had become separated from his pod, could be returned to it. However, his case was marked by controversy about whether and how to intervene, and in 2006, Luna was killed by a boat propeller.[218]
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The earlier of known records of commercial hunting of killer whales date to the 18th century in Japan. During the 19th and early 20th centuries, the global whaling industry caught immense numbers of baleen and sperm whales, but largely ignored killer whales because of their limited amounts of recoverable oil, their smaller populations, and the difficulty of taking them.[143] Once the stocks of larger species were depleted, killer whales were targeted by commercial whalers in the mid-20th century. Between 1954 and 1997, Japan took 1,178 killer whales (although the Ministry of the Environment claims that there had been domestic catches of about 1,600 whales between late 1940s to 1960s[219]) and Norway took 987.[220] Extensive hunting of killer whales, including an Antarctic catch of 916 in 1979–80 alone, prompted the International Whaling Commission to recommend a ban on commercial hunting of the species pending further research.[220] Today, no country carries out a substantial hunt, although Indonesia and Greenland permit small subsistence hunts (see Aboriginal whaling). Other than commercial hunts, killer whales were hunted along Japanese coasts out of public concern for potential conflicts with fisheries. Such cases include a semi-resident male-female pair in Akashi Strait and Harimanada being killed in the Seto Inland Sea in 1957,[221][222] the killing of five whales from a pod of 11 members that swam into Tokyo Bay in 1970,[223] and a catch record in southern Taiwan in the 1990s.[79][224]
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Killer whales have helped humans hunting other whales.[225] One well-known example was the killer whales of Eden, Australia, including the male known as Old Tom. Whalers more often considered them a nuisance, however, as orcas would gather to scavenge meat from the whalers' catch.[225] Some populations, such as in Alaska's Prince William Sound, may have been reduced significantly by whalers shooting them in retaliation.[16]
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Whale watching continues to increase in popularity, but may have some problematic impacts on killer whales. Exposure to exhaust gasses from large amounts of vessel traffic are causing concern for the overall health of the 75 remaining southern resident killer whales (SRKWs) left as of early 2019.[226] This population is followed by approximately 20 vessels for 12 hours a day during the months May–September.[227] Researchers discovered that these vessels are in the line of sight for these whales for 98–99.5% of daylight hours.[227] With so many vessels, the air quality around these whales deteriorates and impacts their health. Air pollutants that bind with exhaust fumes are responsible for the activation of the cytochrome P450 1A gene family.[227] Researchers have successfully identified this gene in skin biopsies of live whales and also the lungs of deceased whales. A direct correlation between activation of this gene and the air pollutants can not be made because there are other known factors that will induce the same gene. Vessels can have either wet or dry exhaust systems, with wet exhaust systems leaving more pollutants in the water due to various gas solubility. A modelling study determined that the lowest-observed-adverse-effect-level (LOAEL) of exhaust pollutants was about 12% of the human dose.[227]
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As a response to this, in 2017 boats off the British Columbia coast now have a minimum approach distance of 200 metres compared to the previous 100 metres. This new rule complements Washington State's minimum approach zone of 180 metres that has been in effect since 2011. If a whale approaches a vessel it must be placed in neutral until the whale passes. The World Health Organization has set air quality standards in an effort to control the emissions produced by these vessels.[228]
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The killer whale's intelligence, trainability, striking appearance, playfulness in captivity and sheer size have made it a popular exhibit at aquaria and aquatic theme parks. From 1976 to 1997, 55 whales were taken from the wild in Iceland, 19 from Japan, and three from Argentina. These figures exclude animals that died during capture. Live captures fell dramatically in the 1990s, and by 1999, about 40% of the 48 animals on display in the world were captive-born.[229]
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Organizations such as World Animal Protection and the Whale and Dolphin Conservation campaign against the practice of keeping them in captivity. In captivity, they often develop pathologies, such as the dorsal fin collapse seen in 60–90% of captive males. Captives have vastly reduced life expectancies, on average only living into their 20s.[c] In the wild, females who survive infancy live 46 years on average, and up to 70–80 years in rare cases. Wild males who survive infancy live 31 years on average, and up to 50–60 years.[230] Captivity usually bears little resemblance to wild habitat, and captive whales' social groups are foreign to those found in the wild. Critics claim captive life is stressful due to these factors and the requirement to perform circus tricks that are not part of wild killer whale behaviour, see above.[231] Wild killer whales may travel up to 160 kilometres (100 mi) in a day, and critics say the animals are too big and intelligent to be suitable for captivity.[155] Captives occasionally act aggressively towards themselves, their tankmates, or humans, which critics say is a result of stress.[207] Between 1991 and 2010, the bull orca known as Tilikum was involved in the death of three people, and was featured in the critically acclaimed 2013 film Blackfish.[232] Tilikum lived at SeaWorld from 1992 until his death in 2017.[233][234][235][236][237]
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A 2015 study coauthored by staff at SeaWorld and the Minnesota Zoo indicates that there is no significant difference in survivorship between free-ranging and captive killer whales. The authors speculate about the future utility of studying captive populations for the purposes of understanding orca biology and the implications of such research of captive animals in the overall health of both wild and marine park populations.[238]
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As of March 2016, SeaWorld has announced that they will be ending their orca breeding program and their theatrical shows. They previously announced, in November 2015, that the shows would be coming to an end in San Diego but it is now to happen in both Orlando and San Antonio as well.[239]
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Athens (/ˈæθɪnz/ ATH-inz;[3] Greek: Αθήνα, romanized: Athína [aˈθina] (listen); Ancient Greek: Ἀθῆναι, romanized: Athênai (pl.) [atʰɛ̂ːnai̯]), is the capital and largest city of Greece. Athens dominates the Attica region and is one of the world's oldest cities, with its recorded history spanning over 3,400 years[citation needed] and its earliest human presence started somewhere between the 11th and 7th millennium BC.[4]
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Classical Athens was a powerful city-state that emerged in conjunction with the seagoing development of the port of Piraeus. A center for the arts, learning and philosophy, home of Plato's Academy and Aristotle's Lyceum,[5][6] it is widely referred to as the cradle of Western civilization and the birthplace of democracy,[7][8] largely because of its cultural and political impact on the European continent, and in particular the Romans.[9] In modern times, Athens is a large cosmopolitan metropolis and central to economic, financial, industrial, maritime, political and cultural life in Greece.
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Athens is a global city and one of the biggest economic centers in southeastern Europe. It has a large financial sector, and its port Piraeus is both the largest passenger port in Europe,[10][11] and the second largest in the world.[12]
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The Municipality of Athens (also City of Athens), which actually constitutes a small administrative unit of the entire city, had a population of 664,046 (in 2011)[2] within its official limits, and a land area of 38.96 km2 (15.04 sq mi).[13][14] The Athens Urban Area (Greater Athens and Greater Piraeus) extends beyond its administrative municipal city limits, with a population of 3,090,508 (in 2011)[15] over an area of 412 km2 (159 sq mi).[14] According to Eurostat[16] in 2011, the functional urban area (FUA) of Athens was the 9th most populous FUA in the European Union (the 6th most populous capital city of the EU), with a population of 3.8 million people. Athens is also the southernmost capital on the European mainland and the warmest major city in Europe.
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The heritage of the classical era is still evident in the city, represented by ancient monuments and works of art, the most famous of all being the Parthenon, considered a key landmark of early Western civilization. The city also retains Roman and Byzantine monuments, as well as a smaller number of Ottoman monuments. Athens is home to two UNESCO World Heritage Sites, the Acropolis of Athens and the medieval Daphni Monastery. Landmarks of the modern era, dating back to the establishment of Athens as the capital of the independent Greek state in 1834, include the Hellenic Parliament and the so-called "architectural trilogy of Athens", consisting of the National Library of Greece, the National and Kapodistrian University of Athens and the Academy of Athens. Athens is also home to several museums and cultural institutions, such as the National Archeological Museum, featuring the world's largest collection of ancient Greek antiquities, the Acropolis Museum, the Museum of Cycladic Art, the Benaki Museum and the Byzantine and Christian Museum. Athens was the host city of the first modern-day Olympic Games in 1896, and 108 years later it hosted the 2004 Summer Olympics, making it one of the few cities to have hosted the Olympics more than once.[17]
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In Ancient Greek, the name of the city was Ἀθῆναι (Athênai, pronounced [atʰɛ̂ːnai̯] in Classical Attic) a plural. In earlier Greek, such as Homeric Greek, the name had been current in the singular form though, as Ἀθήνη (Athḗnē).[18] It was possibly rendered in the plural later on, like those of Θῆβαι (Thêbai) and Μυκῆναι (Μukênai). The root of the word is probably not of Greek or Indo-European origin,[19] and is possibly a remnant of the Pre-Greek substrate of Attica.[19] In antiquity, it was debated whether Athens took its name from its patron goddess Athena (Attic Ἀθηνᾶ, Athēnâ, Ionic Ἀθήνη, Athḗnē, and Doric Ἀθάνα, Athā́nā) or Athena took her name from the city.[20] Modern scholars now generally agree that the goddess takes her name from the city,[20] because the ending -ene is common in names of locations, but rare for personal names.[20] During the medieval period, the name of the city was rendered once again in the singular as Ἀθήνα. However, after the establishment of the modern Greek state, and partly due to the conservatism of the written language, Ἀθῆναι [aˈθine] became again the official name of the city and remained so until the abandonment of Katharevousa in the 1970s, when Ἀθήνα, Athína, became the official name.
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According to the ancient Athenian founding myth, Athena, the goddess of wisdom, competed against Poseidon, the God of the Seas, for patronage of the yet-unnamed city;[21] they agreed that whoever gave the Athenians the better gift would become their patron[21] and appointed Cecrops, the king of Athens, as the judge.[21] According to the account given by Pseudo-Apollodorus, Poseidon struck the ground with his trident and a salt water spring welled up.[21] In an alternative version of the myth from Vergil's Georgics, Poseidon instead gave the Athenians the first horse.[21] In both versions, Athena offered the Athenians the first domesticated olive tree.[21][22] Cecrops accepted this gift[21] and declared Athena the patron goddess of Athens.[21][22] Eight different etymologies, now commonly rejected, were proposed during the 17th century.[citation needed] Christian Lobeck proposed as the root of the name the word ἄθος (áthos) or ἄνθος (ánthos) meaning "flower", to denote Athens as the "flowering city". Ludwig von Döderlein proposed the stem of the verb θάω, stem θη- (tháō, thē-, "to suck") to denote Athens as having fertile soil.[23]
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In classical literature, the city was sometimes referred to as the City of the Violet Crown, first documented in Pindar's ἰοστέφανοι Ἀθᾶναι (iostéphanoi Athânai), or as τὸ κλεινὸν ἄστυ (tò kleinòn ásty, "the glorious city"). In medieval texts, variant names include Setines, Satine, and Astines, all derivations involving false splitting of prepositional phrases.[24] Today the caption η πρωτεύουσα (ī protévousa), "the capital", has become somewhat common. King Alphonse X of Castile in his 'General Storia' (Global History) cites etymology of Athens as: 'Athinai', 'the one without death', death here meaning ignorance.
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Athenians were called cicada-wearers (Ancient Greek: Τεττιγοφόροι) because they used to wear pins of golden cicadas. A symbol of being autochthon (earth-born), because the legendary founder of Athens, Erechtheus was an autochthon or of being musicians, because the cicada is a "musician" insect.[25]
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The oldest known human presence in Athens is the Cave of Schist, which has been dated to between the 11th and 7th millennia BC.[4] Athens has been continuously inhabited for at least 5,000 years.[26][27] By 1400 BC the settlement had become an important centre of the Mycenaean civilization and the Acropolis was the site of a major Mycenaean fortress, whose remains can be recognised from sections of the characteristic Cyclopean walls.[28] Unlike other Mycenaean centers, such as Mycenae and Pylos, it is not known whether Athens suffered destruction in about 1200 BC, an event often attributed to a Dorian invasion, and the Athenians always maintained that they were pure Ionians with no Dorian element. However, Athens, like many other Bronze Age settlements, went into economic decline for around 150 years afterwards.
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Iron Age burials, in the Kerameikos and other locations, are often richly provided for and demonstrate that from 900 BC onwards Athens was one of the leading centres of trade and prosperity in the region.[29] The leading position of Athens may well have resulted from its central location in the Greek world, its secure stronghold on the Acropolis and its access to the sea, which gave it a natural advantage over inland rivals such as Thebes and Sparta.
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By the 6th century BC, widespread social unrest led to the reforms of Solon. These would pave the way for the eventual introduction of democracy by Cleisthenes in 508 BC. Athens had by this time become a significant naval power with a large fleet, and helped the rebellion of the Ionian cities against Persian rule. In the ensuing Greco-Persian Wars Athens, together with Sparta, led the coalition of Greek states that would eventually repel the Persians, defeating them decisively at Marathon in 490 BC, and crucially at Salamis in 480 BC. However, this did not prevent Athens from being captured and sacked twice by the Persians within one year, after a heroic but ultimately failed resistance at Thermopylae by Spartans and other Greeks led by King Leonidas,[30] after both Boeotia and Attica fell to the Persians.
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The decades that followed became known as the Golden Age of Athenian democracy, during which time Athens became the leading city of Ancient Greece, with its cultural achievements laying the foundations for Western civilization. The playwrights Aeschylus, Sophocles and Euripides flourished in Athens during this time, as did the historians Herodotus and Thucydides, the physician Hippocrates, and the philosopher Socrates. Guided by Pericles, who promoted the arts and fostered democracy, Athens embarked on an ambitious building program that saw the construction of the Acropolis of Athens (including the Parthenon), as well as empire-building via the Delian League. Originally intended as an association of Greek city-states to continue the fight against the Persians, the league soon turned into a vehicle for Athens's own imperial ambitions. The resulting tensions brought about the Peloponnesian War (431–404 BC), in which Athens was defeated by its rival Sparta.
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By the mid-4th century BC, the northern Greek kingdom of Macedon was becoming dominant in Athenian affairs. In 338 BC the armies of Philip II defeated an alliance of some of the Greek city-states including Athens and Thebes at the Battle of Chaeronea, effectively ending Athenian independence. Later, under Rome, Athens was given the status of a free city because of its widely admired schools. The Roman emperor Hadrian, in the 2nd century CE, ordered the construction of a library, a gymnasium, an aqueduct which is still in use, several temples and sanctuaries, a bridge and financed the completion of the Temple of Olympian Zeus.
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By the end of Late Antiquity, Athens had shrunk due to sacks by the Herulians, Visigoths, and Early Slavs which caused massive destruction in the city. In this era, the first Christian churches were built in Athens, and the Parthenon and other temples were converted into churches. Athens expanded its settlement in the second half of the Middle Byzantine Period, in the 9th to 10th centuries CE, and was relatively prosperous during the Crusades, benefiting from Italian trade. After the Fourth Crusade the Duchy of Athens was established. In 1458 it was conquered by the Ottoman Empire and entered a long period of decline.
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Following the Greek War of Independence and the establishment of the Greek Kingdom, Athens was chosen as the capital of the newly independent Greek state in 1834, largely because of historical and sentimental reasons. At the time, it was reduced to a town of about 4,000 people in a loose swarm of houses along the foot of the Acropolis. The first King of Greece, Otto of Bavaria, commissioned the architects Stamatios Kleanthis and Eduard Schaubert to design a modern city plan fit for the capital of a state.
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The first modern city plan consisted of a triangle defined by the Acropolis, the ancient cemetery of Kerameikos and the new palace of the Bavarian king (now housing the Greek Parliament), so as to highlight the continuity between modern and ancient Athens. Neoclassicism, the international style of this epoch, was the architectural style through which Bavarian, French and Greek architects such as Hansen, Klenze, Boulanger or Kaftantzoglou designed the first important public buildings of the new capital. In 1896, Athens hosted the first modern Olympic Games. During the 1920s a number of Greek refugees, expelled from Asia Minor after the Greco-Turkish War, swelled Athens's population; nevertheless it was most particularly following World War II, and from the 1950s and 1960s, that the population of the city exploded, and Athens experienced a gradual expansion.
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In the 1980s it became evident that smog from factories and an ever-increasing fleet of automobiles, as well as a lack of adequate free space due to congestion, had evolved into the city's most important challenge. A series of anti-pollution measures taken by the city's authorities in the 1990s, combined with a substantial improvement of the city's infrastructure (including the Attiki Odos motorway, the expansion of the Athens Metro, and the new Athens International Airport), considerably alleviated pollution and transformed Athens into a much more functional city. In 2004 Athens hosted the 2004 Summer Olympics.
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Tondo of the Aison Cup, showing the victory of Theseus over the Minotaur in the presence of Athena. Theseus was responsible, according to the myth, for the synoikismos ("dwelling together")—the political unification of Attica under Athens.
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The earliest coinage of Athens, c. 545–525/15 BC
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Coat of Arms of the Duchy of Athens during the rule of the de la Roche family (13th century)
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The Roman Agora and the Gate of Athena in Plaka district.
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The Temple of Olympian Zeus with river Ilisos by Edward Dodwell, 1821
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The Entry of King Otto in Athens, Peter von Hess, 1839.
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The Stadiou Street in Central Athens in 1908.
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Athens sprawls across the central plain of Attica that is often referred to as the Athens Basin or the Attica Basin (Greek: Λεκανοπέδιο Αθηνών/Αττικής). The basin is bounded by four large mountains: Mount Aigaleo to the west, Mount Parnitha to the north, Mount Pentelicus to the northeast and Mount Hymettus to the east.[31] Beyond Mount Aegaleo lies the Thriasian plain, which forms an extension of the central plain to the west. The Saronic Gulf lies to the southwest. Mount Parnitha is the tallest of the four mountains (1,413 m (4,636 ft)),[32] and has been declared a national park.
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Athens is built around a number of hills. Lycabettus is one of the tallest hills of the city proper and provides a view of the entire Attica Basin. The meteorology of Athens is deemed to be one of the most complex in the world because its mountains cause a temperature inversion phenomenon which, along with the Greek Government's difficulties controlling industrial pollution, was responsible for the air pollution problems the city has faced.[27] This issue is not unique to Athens; for instance, Los Angeles and Mexico City also suffer from similar atmospheric inversion problems.[27]
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The Cephissus river, the Ilisos and the Eridanos stream are the historical rivers of Athens.
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By the late 1970s, the pollution of Athens had become so destructive that according to the then Greek Minister of Culture, Constantine Trypanis, "...the carved details on the five the caryatids of the Erechtheum had seriously degenerated, while the face of the horseman on the Parthenon's west side was all but obliterated."[33] A series of measures taken by the authorities of the city throughout the 1990s resulted in the improvement of air quality; the appearance of smog (or nefos as the Athenians used to call it) has become less common.
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Measures taken by the Greek authorities throughout the 1990s have improved the quality of air over the Attica Basin. Nevertheless, air pollution still remains an issue for Athens, particularly during the hottest summer days. In late June 2007,[34] the Attica region experienced a number of brush fires,[34] including a blaze that burned a significant portion of a large forested national park in Mount Parnitha,[35] considered critical to maintaining a better air quality in Athens all year round.[34] Damage to the park has led to worries over a stalling in the improvement of air quality in the city.[34]
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The major waste management efforts undertaken in the last decade (particularly the plant built on the small island of Psytalia) have greatly improved water quality in the Saronic Gulf, and the coastal waters of Athens are now accessible again to swimmers. In January 2007, Athens faced a waste management problem when its landfill near Ano Liosia, an Athenian suburb, reached capacity.[36] The crisis eased by mid-January when authorities began taking the garbage to a temporary landfill.[36]
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Athens ranks in the lowest percentage for the risk on frequency and severity of terrorist attacks according to the EU Global Terrorism Database (EIU 2007-2016 calculations). The city also ranked 35th in Digital Security, 21st on Health Security, 29th on Infrastructure Security and 41st on Personal Security globally in a 2017 The Economist Intelligence Unit report.[37] It also ranks as a very safe city (39th globally out of 162 cities overall) on the ranking of the safest and most dangerous countries.[38] A 2019 crime index from Numbeo places Athens at 130th position, rating safer than Tampa, Florida or Dublin, Ireland.[39] According to a Mercer 2019 Quality of Living Survey, Athens ranks 89th on the Mercer Quality of Living Survey ranking.[40]
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Athens has a hot-summer Mediterranean climate (Köppen climate classification: Csa). The dominant feature of Athens' climate is alternation between prolonged hot and dry summers and mild winters with moderate rainfall.[41] With an average of 455 millimetres (17.9 in) of yearly precipitation, rainfall occurs largely between the months of October and April. July and August are the driest months when thunderstorms occur sparsely. Furthermore, some coastal areas of Athens, known as the Athens Riviera have a hot semi-arid climate (BSh) according to the climate atlas published by the Hellenic National Meteorological Service (H.N.M.S).[42]
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Owing to the rain shadow of the Pindus Mountains, annual precipitation of Athens is lower than most other parts of Greece, especially western Greece. As an example, Ioannina receives around 1,300 mm (51 in) per year, and Agrinio around 800 mm (31 in) per year. Daily average highs for July have been measured around 34 °C or 93 °F in downtown Athens, but some parts of the city may be even hotter for the higher density of buildings, such as the center,[43] in particular, western areas due to a combination of industrialization and a number of natural factors, knowledge of which has existed since the mid-19th century.[44][45][46]
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Athens is affected by the urban heat island effect in some areas which is caused by human activity,[47][48] altering its temperatures compared to the surrounding rural areas,[49][50][51][52] and leaving detrimental effects on energy usage, expenditure for cooling,[53][54] and health.[48] The urban heat island of the city has also been found to be partially responsible for alterations of the climatological temperature time-series of specific Athens meteorological stations, because of its impact on the temperatures and the temperature trends recorded by some meteorological stations.[55][56][57][58][59] On the other hand, specific meteorological stations, such as the National Garden station and Thiseio meteorological station, are less affected or do not experience the urban heat island.[49][60]
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Athens holds the World Meteorological Organization record for the highest temperature ever recorded in Europe, at 48 °C (118.4 °F), which was recorded in the Elefsina and Tatoi suburbs of Athens on 10 July 1977.[61]
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The Municipality of Athens, the City Centre of the Athens Urban Area, is divided into several districts: Omonoia, Syntagma, Exarcheia, Agios Nikolaos, Neapolis, Lykavittos, Lofos Strefi, Lofos Finopoulou, Lofos Filopappou, Pedion Areos, Metaxourgeio, Aghios Kostantinos, Larissa Station, Kerameikos, Psiri, Monastiraki, Gazi, Thission, Kapnikarea, Aghia Irini, Aerides, Anafiotika, Plaka, Acropolis, Pnyka, Makrygianni, Lofos Ardittou, Zappeion, Aghios Spyridon, Pangrati, Kolonaki, Dexameni, Evaggelismos, Gouva, Aghios Ioannis, Neos Kosmos, Koukaki, Kynosargous, Fix, Ano Petralona, Kato Petralona, Rouf, Votanikos, Profitis Daniil, Akadimia Platonos, Kolonos, Kolokynthou, Attikis Square, Lofos Skouze, Sepolia, Kypseli, Aghios Meletios, Nea Kypseli, Gyzi, Polygono, Ampelokipoi, Panormou-Gerokomeio, Pentagono, Ellinorosson, Nea Filothei, Ano Kypseli, Tourkovounia-Lofos Patatsou, Lofos Elikonos, Koliatsou, Thymarakia, Kato Patisia, Treis Gefyres, Aghios Eleftherios, Ano Patisia, Kypriadou, Menidi, Prompona, Aghios Panteleimonas, Pangrati, Goudi, Vyronas and Ilisia.
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Parnitha National Park is punctuated by well-marked paths, gorges, springs, torrents and caves dotting the protected area. Hiking and mountain-biking in all four mountains are popular outdoor activities for residents of the city. The National Garden of Athens was completed in 1840 and is a green refuge of 15.5 hectares in the centre of the Greek capital. It is to be found between the Parliament and Zappeion buildings, the latter of which maintains its own garden of seven hectares.
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Parts of the City Centre have been redeveloped under a masterplan called the Unification of Archeological Sites of Athens, which has also gathered funding from the EU to help enhance the project.[68][69] The landmark Dionysiou Areopagitou Street has been pedestrianised, forming a scenic route. The route starts from the Temple of Olympian Zeus at Vasilissis Olgas Avenue, continues under the southern slopes of the Acropolis near Plaka, and finishes just beyond the Temple of Hephaestus in Thiseio. The route in its entirety provides visitors with views of the Parthenon and the Agora (the meeting point of ancient Athenians), away from the busy City Centre.
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The hills of Athens also provide green space. Lycabettus, Philopappos hill and the area around it, including Pnyx and Ardettos hill, are planted with pines and other trees, with the character of a small forest rather than typical metropolitan parkland. Also to be found is the Pedion tou Areos (Field of Mars) of 27.7 hectares, near the National Archaeological Museum.
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Athens' largest zoo is the Attica Zoological Park, a 20-hectare (49-acre) private zoo located in the suburb of Spata. The zoo is home to around 2000 animals representing 400 species, and is open 365 days a year. Smaller zoos exist within public gardens or parks, such as the zoo within the National Garden of Athens.
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The Athens Metropolitan Area consists of 58[15] densely populated municipalities, sprawling around the Municipality of Athens (the City Centre) in virtually all directions. For the Athenians, all the urban municipalities surrounding the City Centre are called suburbs. According to their geographic location in relation to the City of Athens, the suburbs are divided into four zones; the northern suburbs (including Agios Stefanos, Dionysos, Ekali, Nea Erythraia, Kifissia, Kryoneri, Maroussi, Pefki, Lykovrysi, Metamorfosi, Nea Ionia, Nea Filadelfeia, Irakleio, Vrilissia, Melissia, Penteli, Chalandri, Agia Paraskevi, Gerakas, Pallini, Galatsi, Psychiko and Filothei); the southern suburbs (including Alimos, Nea Smyrni, Moschato, Tavros, Agios Ioannis Rentis, Kallithea, Piraeus, Agios Dimitrios, Palaio Faliro, Elliniko, Glyfada, Lagonisi, Saronida, Argyroupoli, Ilioupoli, Varkiza, Voula, Vari and Vouliagmeni); the eastern suburbs (including Zografou, Dafni, Vyronas, Kaisariani, Cholargos and Papagou); and the western suburbs (including Peristeri, Ilion, Egaleo, Koridallos, Agia Varvara, Keratsini, Perama, Nikaia, Drapetsona, Chaidari, Petroupoli, Agioi Anargyroi, Ano Liosia, Aspropyrgos, Eleusina, Acharnes and Kamatero).
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The Athens city coastline, extending from the major commercial port of Piraeus to the southernmost suburb of Varkiza for some 25 km (20 mi),[70] is also connected to the City Centre by tram.
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In the northern suburb of Maroussi, the upgraded main Olympic Complex (known by its Greek acronym OAKA) dominates the skyline. The area has been redeveloped according to a design by the Spanish architect Santiago Calatrava, with steel arches, landscaped gardens, fountains, futuristic glass, and a landmark new blue glass roof which was added to the main stadium. A second Olympic complex, next to the sea at the beach of Palaio Faliro, also features modern stadia, shops and an elevated esplanade. Work is underway to transform the grounds of the old Athens Airport – named Elliniko – in the southern suburbs, into one of the largest landscaped parks in Europe, to be named the Hellenikon Metropolitan Park.[71]
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Many of the southern suburbs (such as Alimos, Palaio Faliro, Elliniko, Glyfada, Voula, Vouliagmeni and Varkiza) known as the Athens Riviera, host a number of sandy beaches, most of which are operated by the Greek National Tourism Organisation and require an entrance fee. Casinos operate on both Mount Parnitha, some 25 km (16 mi)[72] from downtown Athens (accessible by car or cable car), and the nearby town of Loutraki (accessible by car via the Athens – Corinth National Highway, or the suburban rail service Proastiakos).
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The large City Centre (Greek: Κέντρο της Αθήνας) of the Greek capital falls directly within the Municipality of Athens or Athens Municipality (Greek: Δήμος Αθηναίων)—also City of Athens. Athens Municipality is the largest in population size in Greece. Piraeus also forms a significant city centre on its own,[73] within the Athens Urban Area and it is the second largest in population size within it, with Peristeri, Kallithea and Kypseli following.
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The Athens Urban Area (Greek: Πολεοδομικό Συγκρότημα Αθηνών), also known as Urban Area of the Capital (Greek: Πολεοδομικό Συγκρότημα Πρωτεύουσας), today consists of 40 municipalities, 35 of which make up what is referred to as the Greater Athens municipalities, located within 4 regional units (North Athens, West Athens, Central Athens, South Athens); and a further 5, which make up the Greater Piraeus municipalities, located within the regional unit of Piraeus as mentioned above. The densely built up urban area of the Greek capital sprawls across 412 km2 (159 sq mi)[14] throughout the Attica Basin and has a total population of 3,074,160 (in 2011).
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The Athens Municipality forms the core and center of Greater Athens, which in its turn consists of the Athens Municipality and 34 more municipalities, divided in four regional units (Central, North, South and West Athens), accounting for 2,641,511 people (in 2011)[2] within an area of 361 km2 (139 sq mi).[14] Until 2010, these four regional units made up the abolished Athens Prefecture. The municipality of Piraeus, the historic Athenian port, with its 4 suburban municipalities make up the regional unit of Piraeus, which in turn forms Greater Piraeus.
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Greater Athens and Greater Piraeus with part of East[74] and West Attica[75] regional units combined make up the continuous Athens Urban Area,[75][76][77] also called the "Urban Area of the Capital" or simply "Athens" (the most common use of the term), spanning over 412 km2 (159 sq mi),[78] with a population of 3,090,508 people as of 2011. The Athens Urban Area is considered to form the city of Athens as a whole, despite its administrative divisions, which is the largest in Greece and one of the most populated urban areas in Europe.
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The Athens Metropolitan Area (Greek: Μητροπολιτική Περιοχή της Αθήνας), also commonly known in English as Greater Athens,[79] spans 2,928.717 km2 (1,131 sq mi) within the Attica region and includes a total of 58 municipalities, which are organized in seven regional units (those outlined above, along with East Attica and West Attica), having reached a population of 3,737,550 based on the preliminary results of the 2011 census. Athens and Piraeus municipalities serve as the two metropolitan centres of the Athens Metropolitan Area.[80] There are also some inter-municipal centres serving specific areas. For example, Kifissia and Glyfada serve as inter-municipal centres for northern and southern suburbs respectively.
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The Municipality of Athens has an official population of 664,046 people.[2] The four regional units that make up what is referred to as Greater Athens have a combined population of 2,640,701. They together with the regional unit of Piraeus (Greater Piraeus) make up the dense Athens Urban Area which reaches a total population of 3,090,508 inhabitants (in 2011).[15] As Eurostat the FUA of Athens had in 2013 3,828,434 inhabitants, being apparently decreasing compared with the pre-economic crisis date of 2009 (4,164,175)[16]
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The municipality (Center) of Athens is the most populous in Greece, with a population of 664,046 people (in 2011)[2] and an area of 38.96 km2 (15.04 sq mi),[13] forming the core of the Athens Urban Area within the Attica Basin. The incumbent Mayor of Athens is Kostas Bakoyannis of New Democracy. The municipality is divided into seven municipal districts which are mainly used for administrative purposes.
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As of the 2011 census, the population for each of the seven municipal districts of Athens is as follows:[81]
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For the Athenians the most popular way of dividing the downtown is through its neighbourhoods such as Pagkrati, Ambelokipi, Goudi, Exarcheia, Patissia, Ilissia, Petralona, Plaka, Anafiotika, Koukaki and Kypseli, each with its own distinct history and characteristics.
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The Athens Metropolitan Area, with an area of 2,928.717 km2 (1,131 sq mi) and inhabited by 3,753,783 people in 2011,[2] consists of the Athens Urban Area with the addition of the towns and villages of East and West Attica, which surround the dense urban area of the Greek capital. It actually sprawls over the whole peninsula of Attica, which is the best part of the region of Attica, excluding the islands.
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Mycenean Athens in 1600–1100 BC could have reached the size of Tiryns; that would put the population at the range of 10,000–15,000.[82] During the Greek Dark Ages the population of Athens was around 4,000 people. In 700 BC the population grew to 10,000. In 500 BC the area probably contained 200,000 people. During the classical period the city's population is estimated from 150,000–350,000 and up to 610,000 according to Thucydides. When Demetrius of Phalerum conducted a population census in 317 BC the population was 21,000 free citizens, plus 10,000 resident aliens and 400,000 slaves. This suggests a total population of 431,000.[83][84][85][86][87][88][89][90][91] This figure is highly suspect because of the lopsided number of slaves and does not include free women and children and resident foreigners: an estimated based on Thucydides is: 40,000 male citizens, 100,000 family members, 70,000 metics (resident foreigners) and 150,000-400,000 slaves. However the numbers would include all of Athenian territory which consisted almost entirely of modern Attica region expect the territory of the city-state of Megaris and the island section.
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The ancient site of the main city is centred on the rocky hill of the acropolis. In the whole of Athenian territory they existed many towns. Acharnae, Afidnes, Cytherus, Colonus, Corydallus, Cropia, Decelea, Euonymos, Vravron among others was important towns in Athenian countryside. The new port of Piraeus was a prototype harbour with the infrastructure and housing located in the site between modern passenger section of the port (named Kantharos in ancient times) and the Pasalimani harbour (named Zea in ancient times). The old one Phaliro was in the site of modern Palaio Faliro and gradually declined after the construction of the new prototype port but remained as a minor port and important settlement with historic significance in late Classical times. The rapid expansion of the modern city, which continues to this day, was initiated in the 1950s and 1960s, because of Greece's transition from an agricultural to an industrial nation.[92] The expansion is now particularly toward the East and North East (a tendency greatly related to the new Eleftherios Venizelos International Airport and the Attiki Odos, the freeway that cuts across Attica). By this process Athens has engulfed many former suburbs and villages in Attica, and continues to do so. The table below shows the historical population of Athens in recent times.
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Athens became the capital of Greece in 1834, following Nafplion, which was the provisional capital from 1829. The municipality (City) of Athens is also the capital of the Attica region. The term Athens can refer either to the Municipality of Athens, to Greater Athens, or to the entire Athens Urban Area.
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The Hellenic Parliament
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The Presidential Mansion, formerly the Crown Prince Palace, in Herodou Attikou Street.
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The Maximos Mansion, official office of the Prime Minister of the Hellenic Republic, in Herodou Attikou Street.
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The Athens City Hall in Kotzia Square was designed by Panagiotis Kolkas and completed in 1874.[98]
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The Embassy of France in Vasilissis Sofias Avenue.
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The Italian Embassy in Vasilissis Sofias Avenue.
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Athens is twinned with:[99]
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Athens is the financial capital of Greece. According to data from 2014, Athens as a metropolitan economic area produced 130 billion US-dollars as GDP in PPP, which consists nearly a half of the production for the whole country. In the list with the strongest economic metropoles of the world Athens was ranked that year 102nd, while the GDP per capita for the same year was 32,000 US-dollars.[112]
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Athens is one of the major economic centres in south-eastern Europe and is considered as a regional economic power in Europe generally. The port of Piraeus, where big investments by COSCO have already been delivered during the recent decade, the completion of the new Cargo Centre in Thriasion,[113] the expansion of the Athens Metro and the Athens Tram, as well as the projected metropolitan park in Elliniko and other economic projects are the economic landmarks of the upcoming years.
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Important Greek companies such as Hellenic Aerospace Industry, Mytilineos Holdings, Titan Cement, Hellenic Petroleum, Papadopoulos E.J., Folli Follie, Jumbo S.A., OPAP, and COSMOTE have their headquarters in the metropolitan area of Athens. Multinational companies such as Ericsson, Sony, Siemens, Motorola, Samsung, Microsoft, Novartis, Mondelez, Coca-Cola, etc. have their regional research and development headquarters also there.
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The banking sector is represented by National Bank of Greece, Alpha Bank, Eurobank, and Piraeus Bank, while the Bank of Greece is also situated in the City Centre. The Athens Stock Exchange, the only in Greece, has been severely hit by the Greek government-debt crisis and the decision of the government to proceed into capital controls during summer 2015. As a whole the economy of Athens and Greece has been severely hit with today's data showing a change from long recession to growth of 1.4% in 2017.[114]
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Tourism is also a great contributor for the economy of the city, which is considered as one of the top destinations in Europe for city-break tourism and is also the gateway for excursions to the islands or the mainland. Greece attracted 26.5 million visitors in 2015, 30.1 million visitors in 2017 and over 33 million in 2018, making Greece one of the most visited countries in Europe and the world, and contributing 18% to the nation's Gross Domestic Product. Athens welcomed more than 5 million tourists in 2018 and 1,4 million of them were "city-breakers" (in 2013 the city-breakers were only 220.000).[115]
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Athens is the major transportation hub of Greece. They combines the largest Airport in Greece as also the largest port in Greece which also the largest port in Mediterranean in containers transport and the largest passenger port in Europe. Is also major national hub for Intercity (Ktel) and international buses as well as for domestic and international rail transport. Public transport is serviced by a variety of transportation means, forming the largest mass transit system of Greece. The Athens Mass Transit System consists of a large bus fleet, a trolleybus fleet that mainly serves Athens's city center, the city's Metro, a Proastiakos service[116] and a tram network, connecting the southern suburbs to the City Centre.[117]
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OSY (Greek: ΟΣΥ) (Odikes Sygkoinonies S.A.) which is subsidiary company of OASA (Athens urban transport organisation), is the main operator of buses and trolleybusses in Athens. Its network consists of about 300 bus lines and 22 trolleybus lines which span the Athens Metropolitan Area,[118] with a fleet of 1,839 buses and 366 trolleybuses.[119] Of those 1,839 buses 416 run on compressed natural gas,[119] making up the largest fleet of natural gas-powered buses in Europe[120] and all trolleybusses are equipped to enable them to run on diesel in case of power failure.[121]
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International links are provided by a number of private companies. National and regional bus links are provided by KTEL from two InterCity Bus Terminals, Kifissos Bus Terminal A and Liosion Bus Terminal B, both located in the north-western part of the city. Kifissos provides connections towards Peloponnese, North Greece, West Greece and some Ionian Islands, whereas Liosion is used for most of Central Greece.
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The Athens Metro is operated by STASY S.A (Greek: ΣΤΑΣΥ) (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation) and provides public transport throughout the Athens Urban Area. While its main purpose is transport, it also houses Greek artifacts found during construction of the system.[122] The Athens Metro has an operating staff of 387 and runs three metro lines, namely the line 1 (Green Line), line 2 (Red Line) and line 3 (Blue Line) lines, of which the first was constructed
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in 1869, and the other two largely during the 1990s, with the initial new sections opened in January 2000. Line 1 mostly runs at ground level and the other two (Line 2 & 3) routes run entirely underground. A fleet of 42 trains, using 252 carriages, operates on the network,[123] with a daily occupancy of 1,353,000 passengers.[124]
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The line 1 (Green Line) serves 24 stations, and is the oldest line of the Athens metro network. It runs from Piraeus station to Kifissia station and covers a distance of 25.6-kilometre (15.9 mi). There are transfer connections with the Blue Line 3 at Monastiraki station and with the Red Line 2 at Omonoia and Attiki stations.
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The line 2 (Red Line) runs from Anthoupoli station to Elliniko station and covers a distance of 17.5 km (10.9 mi).[123] The line connects the western suburbs of Athens with the southeast suburbs, passing through the center of Athens. The Red Line has transfer connections with the Green Line 1 at Attiki and Omonoia stations. There are also transfer connections with the Blue Line 3 at Syntagma station and with the tram at Syntagma, Syngrou Fix and Neos Kosmos stations.
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The line 3 (Blue Line) runs from Agia Marina and Egaleo stations, through the central Monastiraki and Syntagma stations to Doukissis Plakentias avenue in the northeastern suburb of Halandri, covering a distance of 16 km (10 mi).[123] It then ascends to ground level and continues to Eleftherios Venizelos International Airport, using the Suburban Railway infrastructure, extending its total length to 39 km (24 mi).[123] The spring 2007 extension from Monastiraki westwards to Egaleo connected some of the main night life hubs of the city, namely those of Gazi (Kerameikos station) with Psirri (Monastiraki station) and the City Centre (Syntagma station). Extensions are under construction to the west southwest suburbs of Athens, as far as the port of Piraeus. The new stations will be Agia Barvara, Koridallos, Nikaia, Maniatika, Piraeus and Dimotiko Theatro. The completed extension will be ready in 2022 (the first three stations will open in summer 2020), connecting the biggest port of Greece, Piraeus Port, with the biggest airport of Greece the Athens International Airport.
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The Athens commuter rail service, referred to as the "Proastiakós", connects Eleftherios Venizelos International Airport to the city of Kiato, 106 km (66 mi)[125] west of Athens, via Larissa station, the city's central rail station and the port of Piraeus. The length of Athens's commuter rail network extends to 120 km (75 mi),[125] and is expected to stretch to 281 km (175 mi) by 2010.[125]
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Athens tram is operated by STASY S.A (Statheres Sygkoinonies S.A) which is a subsidiary company of OASA (Athens urban transport organisation). It has a fleet of 35 Sirio type vehicles[126] which serve 48 stations,[126] employ 345 people with an average daily occupancy of 65,000 passengers.[126] The tram network spans a total length of 27 km (17 mi) and covers ten Athenian suburbs.[126] The network runs from Syntagma Square to the southwestern suburb of Palaio Faliro, where the line splits in two branches; the first runs along the Athens coastline toward the southern suburb of Voula, while the other heads toward Neo Faliro. The network covers the majority of the Athens coastline.[127] Further extensions are planned towards the major commercial port of Piraeus.[126] The expansion to Piraeus will include 12 new stations, increase the overall length of tram route by 5.4 km (3 mi), and increase the overall transportation network.[128]
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Athens is served by the Athens International Airport (ATH), located near the town of Spata, in the eastern Messoghia plain, some 35 km (22 mi) east of center of Athens.[129] The airport, awarded the "European Airport of the Year 2004" Award,[130] is intended as an expandable hub for air travel in southeastern Europe and was constructed in 51 months, costing 2.2 billion euros. It employs a staff of 14,000.[130]
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The airport is served by the Metro, the suburban rail, buses to Piraeus port, Athens' City Centre, Liosion and Kifisos Intercity bus stations and Elliniko metro's line 2 southern terminal, and also taxis. The airport accommodates 65 landings and take-offs per hour,[129] with its 24-passenger boarding bridges,[129] 144 check-in counters and broader 150,000 m2 (1,614,587 sq ft) main terminal;[129] and a commercial area of 7,000 m2 (75,347 sq ft) which includes cafés, duty-free shops,[130] and a small museum.
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In 2018, the airport handled 24,135,736 a huge increase over the last 4 years. In 2014, the airport handled 15,196,369 passengers, an increase of 21.2% over the previous year of 2013.[131] Of those 15,196,369 passengers, 5,267,593 passed through the airport for domestic flights,[132] and 9,970,006 passengers travelled through for international flights.[132] Beyond the dimensions of its passenger capacity, ATH handled 205,294 total flights in 2007, or approximately 562 flights per day.[133]
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Athens is the hub of the country's national railway system (OSE), connecting the capital with major cities across Greece and abroad (Istanbul, Sofia, Belgrade and Bucharest). The Port of Piraeus is the largest port in Greece and one of the largest in Europe. It is the largest container port in East Mediterranean Sea Basin. It is also the busiest passenger port in Europe and one of the largest passenger ports in the world. Connects Athens to the numerous Greek islands of the Aegean Sea, with ferries departing, while also serving the cruise ships that arrive.[134][135][136] Rafina and Lavrio act as alternative ports of Athens, connects the city with numerous Greek islands of the Aegean Sea, Evia and Cesme in Turkey,[137][138] while also serving the cruise ships that arrive.
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Two main motorways of Greece begin in Athens, namely the A1/E75, heading north towards Greece's second largest city, Thessaloniki; and the border crossing of Evzones and the A8/E94 heading west, towards Greece's third largest city, Patras, which incorporated the GR-8A. Before their completion much of the road traffic used the GR-1 and the GR-8.
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Athens' Metropolitan Area is served by the motorway network of the Attiki Odos toll-motorway (code: A6). Its main section extends from the western industrial suburb of Elefsina to Athens International Airport; while two beltways, namely the Aigaleo Beltway (A65) and the Hymettus Beltway (A64) serve parts of western and eastern Athens respectively. The span of the Attiki Odos in all its length is 65 km (40 mi),[139] making it the largest metropolitan motorway network in all of Greece.
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Located on Panepistimiou Street, the old campus of the University of Athens, the National Library, and the Athens Academy form the "Athens Trilogy" built in the mid-19th century. Most of the university's workings have been moved to a much larger, modern campus located in the eastern suburb of Zografou. The second higher education institution in the city is the Athens Polytechnic School, found in Patission Street. This was the location where on 17 November 1973, more than 13 students were killed and hundreds injured inside the university during the Athens Polytechnic uprising,[140] against the military junta that ruled the nation from 21 April 1967 until 23 July 1974.
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Other universities that lie within Athens are the Athens University of Economics and Business, the Panteion University, the Agricultural University of Athens and the University of Piraeus. There are overall ten state-supported Institutions of Higher (or Tertiary) education located in the Athens Urban Area, these are by chronological order: Athens School of Fine Arts (1837), National Technical University of Athens (1837), National and Kapodistrian University of Athens (1837), Agricultural University of Athens (1920), Athens University of Economics and Business (1920), Panteion University of Social and Political Sciences (1927), University of Piraeus (1938), Harokopio University of Athens (1990), School of Pedagogical and Technological Education (2002), University of West Attica (2018). There are also several other private colleges, as they called formally in Greece, as the establishment of private universities is prohibited by the constitution. Many of them are accredited by a foreign state or university such as the American College of Greece and the Athens Campus of the University of Indianapolis.[141]
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The city is a world centre of archaeological research. Along with national institutions, such as the Athens University and the Archaeological Society, there are multiple archaeological Museums including the National Archaeological Museum, the Cycladic Museum, the Epigraphic Museum, the Byzantine & Christian Museum, as well as museums at the ancient Agora, Acropolis, Kerameikos, and the Kerameikos Archaeological Museum. The city is also home to the Demokritos laboratory for Archaeometry, alongside regional and national archaeological authorities that form part of the Greek Department of Culture.
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Athens hosts 17 Foreign Archaeological Institutes which promote and facilitate research by scholars from their home countries. As a result, Athens has more than a dozen archaeological libraries and three specialized archaeological laboratories, and is the venue of several hundred specialized lectures, conferences and seminars, as well as dozens of archaeological exhibitions, each year. At any given time, hundreds of international scholars and researchers in all disciplines of archaeology are to be found in the city.
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Athens incorporates architectural styles ranging from Greco-Roman and Neoclassical to modern times. They are often to be found in the same areas, as Athens is not marked by a uniformity of architectural style. A visitor will quickly notice the absence of tall buildings: Athens has very strict height restriction laws in order to ensure the Acropolis hill is visible throughout the city. There is evidence of continuity in elements of the architectural environment through the city's history, although this is still a contested subject [142].
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For the greatest part of the 19th century Neoclassicism dominated Athens, as well as some deviations from it such as Eclecticism, especially in the early 20th century. Thus, the Old Royal Palace was the first important public building to be built, between 1836 and 1843. Later in the mid and late 19th century, Theophil Freiherr von Hansen and Ernst Ziller took part in the construction of many neoclassical buildings such as the Athens Academy and the Zappeion Hall. Ziller also designed many private mansions in the centre of Athens which gradually became public, usually through donations, such as Schliemann's Iliou Melathron.
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Beginning in the 1920s, modern architecture including Bauhaus and Art Deco began to exert an influence on almost all Greek architects, and buildings both public and private were constructed in accordance with these styles. Localities with a great number of such buildings include Kolonaki, and some areas of the centre of the city; neighbourhoods developed in this period include Kypseli.[143]
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In the 1950s and 1960s during the extension and development of Athens, other modern movements such as the International style played an important role. The centre of Athens was largely rebuilt, leading to the demolition of a number of neoclassical buildings. The architects of this era employed materials such as glass, marble and aluminium, and many blended modern and classical elements.[144] After World War II, internationally known architects to have designed and built in the city included Walter Gropius, with his design for the US Embassy, and, among others, Eero Saarinen, in his postwar design for the east terminal of the Ellinikon Airport.
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All over the city can be found several statues or busts. Apart from the neoclassicals by Leonidas Drosis at the Academy of Athens (Plato, Socrates, Apollo, Athena), other notable include the statue of Theseus by Georgios Fytalis at Thiseion, of philhellenes like Lord Byron, George Canning and William Gladstone, the equestrian statue of Theodoros Kolokotronis by Lazaros Sochos in front of the Old Parliament, statues of Ioannis Kapodistrias, Rigas Feraios and Adamantios Korais at the University, of Evangelos Zappas and Konstantinos Zappas at Zappeion, of Ioannis Varvakis at the National Garden, the "woodbreaker" by Dimitrios Filippotis, the equestrian statue of Alexandros Papagos at Papagou district and various busts of fighters of Greek independence at the Pedion tou Areos. A significant landmark is also the Tomb of the Unknown Soldier in Syntagma.
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Athens' most important museums include:
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Athens has been a destination for travellers since antiquity. Over the past decade, the city's infrastructure and social amenities have improved, in part because of its successful bid to stage the 2004 Olympic Games. The Greek Government, aided by the EU, has funded major infrastructure projects such as the state-of-the-art Eleftherios Venizelos International Airport,[145] the expansion of the Athens Metro system,[68] and the new Attiki Odos Motorway.[68]
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Athens was voted as the third best European city to visit in 2015 by European Best Destination. More than 240,000 people voted.
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Athens is home to 148 theatrical stages, more than any other city in the world, including the ancient Odeon of Herodes Atticus, home to the Athens Festival, which runs from May to October each year.[146][147] In addition to a large number of multiplexes, Athens plays host to open air garden cinemas. The city also supports music venues, including the Athens Concert Hall (Megaro Moussikis), which attracts world class artists.[148] The Athens Planetarium,[149] located in Andrea Syngrou Avenue, in Palaio Faliro[150] is one of the largest and best equipped digital planetaria in the world.[151] The Stavros Niarchos Foundation Cultural Center, inaugurated in 2016, will house the National Library of Greece and the Greek National Opera.[152]
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The most successful songs during the period 1870–1930 were the so-called Athenian serenades (Αθηναϊκές καντάδες), based on the Heptanesean kantádhes (καντάδες 'serenades'; sing.: καντάδα) and the songs performed on stage (επιθεωρησιακά τραγούδια 'theatrical revue songs') in revues, musical comedies, operettas and nocturnes that were dominating Athens' theatre scene.
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Notable composers of operettas or nocturnes were Kostas Giannidis, Dionysios Lavrangas, Nikos Hatziapostolou, while Theophrastos Sakellaridis' The Godson remains probably the most popular operetta. Despite the fact that the Athenian songs were not autonomous artistic creations (in contrast with the serenades) and despite their original connection with mainly dramatic forms of Art, they eventually became hits as independent songs. Notable actors of Greek operettas, who made also a series of melodies and songs popular at that time, include Orestis Makris, Kalouta sisters, Vasilis Avlonitis, Afroditi Laoutari, Eleni Papadaki, Marika Nezer, Marika Krevata and others. After 1930, wavering among American and European musical influences as well as the Greek musical tradition. Greek composers begin to write music using the tunes of the tango, waltz, swing, foxtrot, some times combined with melodies in the style of Athenian serenades' repertory. Nikos Gounaris was probably the most renowned composer and singer of the time.
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In 1923, after the population exchange between Greece and Turkey, many ethnic Greeks from Asia Minor fled to Athens as a result of the Greco-Turkish War. They settled in poor neighborhoods and brought with them Rebetiko music, making it popular also in Greece, which became later the base for the Laïko music. Other forms of song popular today in Greece are elafrolaika, entechno, dimotika, and skyladika.[153] Greece's most notable, and internationally famous, composers of Greek song, mainly of the entechno form, are Manos Hadjidakis and Mikis Theodorakis. Both composers have achieved fame abroad for their composition of film scores.[153]
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Athens has a long tradition in sports and sporting events, serving as home to the most important clubs in Greek sport and housing a large number of sports facilities. The city has also been host to sports events of international importance.
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Athens has hosted the Summer Olympic Games twice, in 1896 and 2004. The 2004 Summer Olympics required the development of the Athens Olympic Stadium, which has since gained a reputation as one of the most beautiful stadiums in the world, and one of its most interesting modern monuments.[154] The biggest stadium in the country, it hosted two finals of the UEFA Champions League, in 1994 and 2007. Athens' other major stadium, located in the Piraeus area, is the Karaiskakis Stadium, a sports and entertainment complex, host of the 1971 UEFA Cup Winners' Cup Final.
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Athens has hosted the EuroLeague final three times, the first in 1985 and second in 1993, both at the Peace and Friendship Stadium, most known as SEF, a large indoor arena,[155] and the third time in 2007 at the Olympic Indoor Hall. Events in other sports such as athletics, volleyball, water polo etc., have been hosted in the capital's venues.
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Athens is home to three European multi-sport clubs: Olympiacos, Panathinaikos, AEK Athens. In football, Olympiacos have dominated the domestic competitions, Panathinaikos made it to the 1971 European Cup Final, while AEK Athens is the other member of the big three. These clubs also have basketball teams; Panathinaikos and Olympiacos are among the top powers in European basketball, having won the Euroleague six times and three respectively, whilst AEK Athens was the first Greek team to win a European trophy in any team sport.
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Other notable clubs within Athens are Athinaikos, Panionios, Atromitos, Apollon, Panellinios, Ethnikos Piraeus, Maroussi BCE and Peristeri B.C.. Athenian clubs have also had domestic and international success in other sports.
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The Athens area encompasses a variety of terrain, notably hills and mountains rising around the city, and the capital is the only major city in Europe to be bisected by a mountain range. Four mountain ranges extend into city boundaries and thousands of miles of trails criss-cross the city and neighbouring areas, providing exercise and wilderness access on foot and bike.
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Beyond Athens and across the prefecture of Attica, outdoor activities include skiing, rock climbing, hang gliding and windsurfing. Numerous outdoor clubs serve these sports, including the Athens Chapter of the Sierra Club, which leads over 4,000 outings annually in the area.
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Beside the above clubs, inside the boundaries of Athens Municipality there are some more clubs with presence in national divisions or notable action for short periods. Some of them are PAO Rouf (Rouf) with earlier presence in Gamma Ethniki, Petralona F.C.(el) (Petralona), football club founded in 1963, with earlier presence in Beta Ethniki, Attikos F.C.(el) (Kolonos), football club founded in 1919 with short presence in Gamma Ethniki, Athinais Kypselis [es] (Kypseli), football club founded in 1938 with short presence in Gamma Ethniki, Gyziakos (Gyzi), basketball club founded in 1937 with short presence in Beta Ethniki basketball and Aetos B.C.(el) (Agios Panteleimonas), basketball club founded in 1992 with earlier presence in A2 Ethniki Basketball. Another important Athenian sport club is the Athens Tennis Club founded in 1895 with important offer for the Greek tennis.[156]
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The revival of the modern Olympic Games was brought forth in 1896, by Frenchman Pierre de Coubertin. Thanks to his efforts, Athens was awarded the first modern Olympic Games. In 1896, the city had a population of 123,000[93] and the event helped boost the city's international profile. Of the venues used for these Olympics, the Kallimarmaro Stadium, and Zappeion were most crucial. The Kallimarmaro is a replica of the ancient Athenian stadiums, and the only major stadium (in its capacity of 60,000) to be made entirely of white marble from Mount Penteli, the same material used for construction of the Parthenon.
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Fencing before the king of Greece at the 1896 Summer Olympics.
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The Panathenaic Stadium of Athens (Kallimarmaron) dates back to the 4th century BC and has hosted the first modern Olympic Games in 1896.
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The 1906 Summer Olympics, or the 1906 Intercalated games, were held in Athens. The intercalated competitions were intermediate games to the internationally organized Olympics, and were meant to be organized in Greece every four years, between the main Olympics. This idea later lost support from the IOC and these games were discontinued.
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Athens was awarded the 2004 Summer Olympics on 5 September 1997 in Lausanne, Switzerland, after having lost a previous bid to host the 1996 Summer Olympics, to Atlanta, United States.[17] It was to be the second time Athens would host the games, following the inaugural event of 1896. After an unsuccessful bid in 1990, the 1997 bid was radically improved, including an appeal to Greece's Olympic history. In the last round of voting, Athens defeated Rome with 66 votes to 41.[17] Prior to this round, the cities of Buenos Aires, Stockholm and Cape Town had been eliminated from competition, having received fewer votes.[17]
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During the first three years of preparations, the International Olympic Committee had expressed concern over the speed of construction progress for some of the new Olympic venues. In 2000 the Organising Committee's president was replaced by Gianna Angelopoulos-Daskalaki, who was the president of the original Bidding Committee in 1997. From that point forward, preparations continued at a highly accelerated, almost frenzied pace.
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Although the heavy cost was criticized, estimated at $1.5 billion, Athens was transformed into a more functional city that enjoys modern technology both in transportation and in modern urban development.[157] Some of the finest sporting venues in the world were created in the city, all of which were fully ready for the games. The games welcomed over 10,000 athletes from all 202 countries.[157]
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The 2004 Games were judged a success, as both security and organization worked well, and only a few visitors reported minor problems mainly concerning accommodation issues. The 2004 Olympic Games were described as Unforgettable, dream Games, by IOC President Jacques Rogge for their return to the birthplace of the Olympics, and for meeting the challenges of holding the Olympic Games.[157] The only observable problem was a somewhat sparse attendance of some early events. Eventually, however, a total of more than 3.5 million tickets were sold, which was higher than any other Olympics with the exception of Sydney (more than 5 million tickets were sold there in 2000).[158]
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In 2008 it was reported that most of the Olympic venues had fallen into disrepair: according to those reports, 21 of the 22 facilities built for the games had either been left abandoned or are in a state of dereliction, with several squatter camps having sprung up around certain facilities, and a number of venues afflicted by vandalism, graffiti or strewn with rubbish.[159][160] These claims, however, are disputed and likely to be inaccurate, as most of the facilities used for the Athens Olympics are either in use or in the process of being converted for post-Olympics use. The Greek Government has created a corporation, Olympic Properties SA, which is overseeing the post-Olympics management, development and conversion of these facilities, some of which will be sold off (or have already been sold off) to the private sector,[161] while other facilities are still in use just as during the Olympics, or have been converted for commercial use or modified for other sports.[162] Concerts and theatrical shows, such as those by the troupe Cirque du Soleil, have recently been held in the complex.[153]
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The orca or killer whale (Orcinus orca) is a toothed whale belonging to the oceanic dolphin family, of which it is the largest member. Killer whales have a diverse diet, although individual populations often specialize in particular types of prey. Some feed exclusively on fish, while others hunt marine mammals such as seals and other species of dolphin. They have been known to attack baleen whale calves, and even adult whales. Killer whales are apex predators, as no animal preys on them. A cosmopolitan species, they can be found in each of the world's oceans in a variety of marine environments, from Arctic and Antarctic regions to tropical seas, absent only from the Baltic and Black seas, and some areas of the Arctic Ocean.
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Killer whales are highly social; some populations are composed of matrilineal family groups (pods) which are the most stable of any animal species. Their sophisticated hunting techniques and vocal behaviours, which are often specific to a particular group and passed across generations, have been described as manifestations of animal culture.
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The International Union for Conservation of Nature assesses the orca's conservation status as data deficient because of the likelihood that two or more killer whale types are separate species. Some local populations are considered threatened or endangered due to prey depletion, habitat loss, pollution (by PCBs), capture for marine mammal parks, and conflicts with human fisheries. In late 2005, the southern resident killer whales, which swim in British Columbia and Washington state waters, were placed on the U.S. Endangered Species list.
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Wild killer whales are not considered a threat to humans and no fatal attack on humans has ever been documented, but there have been cases of captive orcas killing or injuring their handlers at marine theme parks. Killer whales feature strongly in the mythologies of indigenous cultures, and their reputation in different cultures ranges from being the souls of humans to merciless killers.
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Orcinus orca is the only recognized extant species in the genus Orcinus, and one of many animal species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae.[6] Konrad Gessner wrote the first scientific description of a killer whale in his Piscium & aquatilium animantium natura of 1558, part of the larger Historia animalium, based on examination of a dead stranded animal in the Bay of Greifswald that had attracted a great deal of local interest.[7]
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The killer whale is one of 35 species in the oceanic dolphin family, which first appeared about 11 million years ago. The killer whale lineage probably branched off shortly thereafter.[8] Although it has morphological similarities with the false killer whale, the pygmy killer whale and the pilot whales, a study of cytochrome b gene sequences by Richard LeDuc indicated that its closest extant relatives are the snubfin dolphins of the genus Orcaella.[9] However, a more recent (2018) study places the orca as a sister taxon to the Lissodelphininae, a clade that includes Lagenorhynchus and Cephalorhynchus.[10]
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Although the term "orca" is increasingly used, English-speaking scientists most often use the traditional name "killer whale"[citation needed]. The genus name Orcinus means "of the kingdom of the dead",[11] or "belonging to Orcus".[12]
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Ancient Romans originally used orca (pl. orcae) for these animals, possibly borrowing Ancient Greek ὄρυξ (óryx), which referred (among other things) to a whale species. Since the 1960s, "orca" has steadily grown in popularity. The term "orca" is preferred by some as it avoids the negative connotations of "killer",[13] and because, being part of the family Delphinidae, the species is more closely related to other oceanic dolphins than to other whales.[14]
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They are sometimes referred to as "blackfish", a name also used for other whale species. "Grampus" is a former name for the species, but is now seldom used. This meaning of "grampus" should not be confused with the genus Grampus, whose only member is Risso's dolphin.[15]
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The three to five types of killer whales may be distinct enough to be considered different races,[16] subspecies, or possibly even species[17] (see Species problem). The IUCN reported in 2008, "The taxonomy of this genus is clearly in need of review, and it is likely that O. orca will be split into a number of different species or at least subspecies over the next few years."[3] Although large variation in the ecological distinctiveness of different killer whale groups complicate simple differentiation into types,[18] research off the west coast of Canada and the United States in the 1970s and 1980s identified the following three types:
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Transients and residents live in the same areas, but avoid each other.[31][32][33]
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Other populations have not been as well studied, although specialized fish and mammal eating killer whales have been distinguished elsewhere.[34] In addition, separate populations of "generalist" (fish- and mammal-eating) and "specialist" (mammal-eating) killer whales have been identified off northwestern Europe.[35][36] As with residents and transients, the lifestyle of these whales appears to reflect their diet; fish-eating killer whales in Alaska[37] and Norway[38] have resident-like social structures, while mammal-eating killer whales in Argentina and the Crozet Islands behave more like transients.[39]
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Three types have been documented in the Antarctic. Two dwarf species, named Orcinus nanus and Orcinus glacialis, were described during the 1980s by Soviet researchers, but most cetacean researchers are sceptical about their status, and linking these directly to the types described below is difficult.[17]
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Types B and C live close to the ice pack, and diatoms in these waters may be responsible for the yellowish colouring of both types.[17][46] Mitochondrial DNA sequences support the theory that these are recently diverged separate species.[47] More recently, complete mitochondrial sequencing indicates the two Antarctic groups that eat seals and fish should be recognized as distinct species, as should the North Pacific transients, leaving the others as subspecies pending additional data.[48] Advanced methods that sequenced the entire mitochondrial genome revealed systematic differences in DNA between different populations.[49] A 2019 study of Type D orcas also found them to be distinct from other populations and possibly even a unique species.[44]
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Mammal-eating killer whales in different regions were long thought likely to be closely related, but genetic testing has refuted this hypothesis.[50]
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There are seven identified ecotypes inhabiting isolated ecological niches. Of three orca ecotypes in the Antarctic, one preys on minke whales, the second on seals and penguins, and the third on fish. Another ecotype lives in the eastern North Atlantic, while the three Northeast Pacific ecotypes are labelled the transient, resident and offshore populations described above. Research has supported a proposal to reclassify the Antarctic seal- and fish-eating populations and the North Pacific transients as a distinct species, leaving the remaining ecotypes as subspecies. The first split in the orca population, between the North Pacific transients and the rest, occurred an estimated 700,000 years ago. Such a designation would mean that each new species becomes subject to separate conservation assessments.[49]
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A typical killer whale distinctively bears a black back, white chest and sides, and a white patch above and behind the eye. Calves are born with a yellowish or orange tint, which fades to white. It has a heavy and robust body[51] with a large dorsal fin up to 1.8 m (5 ft 11 in) tall.[52] Behind the fin, it has a dark grey "saddle patch" across the back. Antarctic killer whales may have pale grey to nearly white backs. Adult killer whales are very distinctive, seldom confused with any other sea creature.[53] When seen from a distance, juveniles can be confused with other cetacean species, such as the false killer whale or Risso's dolphin.[54]
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The killer whale's teeth are very strong, and its jaws exert a powerful grip; the upper teeth fall into the gaps between the lower teeth when the mouth is closed. The firm middle and back teeth hold prey in place, while the front teeth are inclined slightly forward and outward to protect them from powerful jerking movements.[55]
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Killer whales are the largest extant members of the dolphin family. Males typically range from 6 to 8 metres (20 to 26 ft) long and weigh in excess of 6 tonnes (5.9 long tons; 6.6 short tons). Females are smaller, generally ranging from 5 to 7 m (16 to 23 ft) and weighing about 3 to 4 tonnes (3.0 to 3.9 long tons; 3.3 to 4.4 short tons).[56] Calves at birth weigh about 180 kg (400 lb) and are about 2.4 m (7.9 ft) long.[57][58] The skeleton of the killer whale is of the typical delphinid structure, but more robust.[59] Its integument, unlike that of most other dolphin species, is characterized by a well-developed dermal layer with a dense network of fascicles of collagen fibres.[60]
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Killer whale pectoral fins, analogous to forelimbs, are large and rounded, resembling paddles, with those of males significantly larger than those of females. Dorsal fins also exhibit sexual dimorphism, with those of males about 1.8 m (5.9 ft) high, more than twice the size of the female's, with the male's fin more like a tall, elongated isosceles triangle, whereas the female's is shorter and more curved.[61] Males and females also have different patterns of black and white skin in their genital areas.[62] In the skull, adult males have longer lower jaws than females, as well as larger occipital crests.[60]
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An individual killer whale can often be identified from its dorsal fin and saddle patch. Variations such as nicks, scratches, and tears on the dorsal fin and the pattern of white or grey in the saddle patch are unique. Published directories contain identifying photographs and names for hundreds of North Pacific animals. Photographic identification has enabled the local population of killer whales to be counted each year rather than estimated, and has enabled great insight into life cycles and social structures.[63]
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Occasionally a killer whale is white; they have been spotted in the northern Bering Sea and around St. Lawrence Island, and near the Russian coast.[64][65] In February 2008, a white killer whale was photographed 3.2 km (2.0 mi) off Kanaga Volcano in the Aleutian Islands.[64][65] In 2010, the Far East Russia Orca Project (FEROP), co-founded and co-directed by Alexander M. Burdin and Erich Hoyt, filmed an adult male nicknamed Iceberg.[66][67]
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Killer whales have good eyesight above and below the water, excellent hearing, and a good sense of touch. They have exceptionally sophisticated echolocation abilities, detecting the location and characteristics of prey and other objects in the water by emitting clicks and listening for echoes,[68] as do other members of the dolphin family. The mean body temperature of the orca is 36 to 38 °C (97 to 100 °F).[69][70] Like most marine mammals, orcas have a layer of insulating blubber ranging from 7.6 to 10 cm (3.0 to 3.9 in) thick[69] beneath the skin. The pulse is about 60 heartbeats per minute when the orca is at the surface, dropping to 30 beats/min when submerged.[71]
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Killer whales are found in all oceans and most seas. Due to their enormous range, numbers, and density, relative distribution is difficult to estimate,[72] but they clearly prefer higher latitudes and coastal areas over pelagic environments.[73] Areas which serve as major study sites for the species include the coasts of Iceland, Norway, the Valdes Peninsula of Argentina, the Crozet Islands, New Zealand and parts of the west coast of North America, from California to Alaska.[74]
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Systematic surveys indicate the highest densities of killer whales (>0.40 individuals per 100 km2) in the northeast Atlantic around the Norwegian coast, in the north Pacific along the Aleutian Islands, the Gulf of Alaska and in the Southern Ocean off much of the coast of Antarctica.[72] They are considered "common" (0.20–0.40 individuals per 100 km2) in the eastern Pacific along the coasts of British Columbia, Washington and Oregon, in the North Atlantic Ocean around Iceland and the Faroe Islands. High densities have also been reported but not quantified in the western North Pacific around the Sea of Japan, Sea of Okhotsk, Kuril Islands, Kamchatka and the Commander Islands and in the Southern Hemisphere off southern Brazil and the tip of southern Africa. They are reported as seasonally common in the Canadian Arctic, including Baffin Bay between Greenland and Nunavut, as well as Tasmania and Macquarie Island.[72] Regularly occurring or distinct populations exist off Northwest Europe, California, Patagonia, the Crozet Islands, Marion Island, southern Australia and New Zealand.[36][72][75] The northwest Atlantic population of at least 67 individuals ranges from Labrador and Newfoundland to New England with sightings to Cape Cod and Long Island.[76]
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Information for offshore regions and warmer waters is more scarce, but widespread sightings indicate that the killer whale can survive in most water temperatures. They have been sighted, though more infrequently, in the Mediterranean, the Arabian Sea, the Gulf of Mexico, Banderas Bay on Mexico's west coast and the Caribbean.[72] Over 50 individual whales have been documented in the northern Indian Ocean, including two individuals that were sighted in the Persian Gulf in 2008 and off Sri Lanka in 2015.[77] Those orcas may occasionally enter the Red Sea through the Gulf of Aden.[78] The modern status of the species along coastal mainland China and its vicinity is unknown. Recorded sightings have been made from almost the entire shoreline.[79] A wide-ranging population is likely to exist in the central Pacific, with some sightings off Hawaii.[80][81] Distinct populations may also exist off the west coast of tropical Africa,[82] and Papua New Guinea.[83] In the Mediterranean, killer whales are considered "visitors", likely from the North Atlantic, and sightings become less frequent further east. However, a small year-round population is known to exist in the Strait of Gibraltar, mostly on the Atlantic side.[84][85] Killer whales also appear to regularly occur off the Galápagos Islands.[86]
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In the Antarctic, killer whales range up to the edge of the pack ice and are believed to venture into the denser pack ice, finding open leads much like beluga whales in the Arctic. However, killer whales are merely seasonal visitors to Arctic waters, and do not approach the pack ice in the summer. With the rapid Arctic sea ice decline in the Hudson Strait, their range now extends deep into the northwest Atlantic.[87] Occasionally, killer whales swim into freshwater rivers. They have been documented 100 mi (160 km) up the Columbia River in the United States.[88][89] They have also been found in the Fraser River in Canada and the Horikawa River in Japan.[88]
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Migration patterns are poorly understood. Each summer, the same individuals appear off the coasts of British Columbia and Washington. Despite decades of research, where these animals go for the rest of the year remains unknown. Transient pods have been sighted from southern Alaska to central California.[90]
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Worldwide population estimates are uncertain, but recent consensus suggests a minimum of 50,000 (2006).[91][3][30] Local estimates include roughly 25,000 in the Antarctic, 8,500 in the tropical Pacific, 2,250–2,700 off the cooler northeast Pacific and 500–1,500 off Norway.[92] Japan's Fisheries Agency estimated in the 2000s that 2,321 killer whales were in the seas around Japan.[93][94]
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Killer whales are apex predators, meaning that they themselves have no natural predators. They are sometimes called the wolves of the sea, because they hunt in groups like wolf packs.[95] Killer whales hunt varied prey including fish, cephalopods, mammals, seabirds, and sea turtles.[96] Different populations or ecotypes may specialize, and some can have a dramatic impact on prey species.[97] However, whales in tropical areas appear to have more generalized diets due to lower food productivity.[81][82] Killer whales spend most of their time at shallow depths,[98] but occasionally dive several hundred meters depending on their prey.[99][100]
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Fish-eating killer whales prey on around 30 species of fish. Some populations in the Norwegian and Greenland sea specialize in herring and follow that fish's autumnal migration to the Norwegian coast. Salmon account for 96% of northeast Pacific residents' diet, including 65% of large, fatty Chinook.[101] Chum salmon are also eaten, but smaller sockeye and pink salmon are not a significant food item.[102] Depletion of specific prey species in an area is, therefore, cause for concern for local populations, despite the high diversity of prey. On average, a killer whale eats 227 kilograms (500 lb) each day.[103] While salmon are usually hunted by an individual whale or a small group, herring are often caught using carousel feeding: the killer whales force the herring into a tight ball by releasing bursts of bubbles or flashing their white undersides. They then slap the ball with their tail flukes, stunning or killing up to 15 fish at a time, then eating them one by one. Carousel feeding has only been documented in the Norwegian killer whale population, as well as some oceanic dolphin species.[104]
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In New Zealand, sharks and rays appear to be important prey, including eagle rays, long-tail and short-tail stingrays, common threshers, smooth hammerheads, blue sharks, basking sharks, and shortfin makos.[105][106] With sharks, orcas may herd them to the surface and strike them with their tail flukes,[105] while bottom-dwelling rays are cornered, pinned to the ground and taken to the surface.[107] In other parts of the world, killer whales have preyed on broadnose sevengill sharks,[108] tiger sharks[109] and even small whale sharks.[110] Killer whales have also been recorded attacking and feeding on great white sharks,[28][111][112][113] and appear to target the liver.[28][112] Competition between killer whales and white sharks is probable in regions where their diets overlap.[114] The arrival of orcas in an area can cause white sharks to flee and forage elsewhere.[115]
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Killer whales are very sophisticated and effective predators of marine mammals. Thirty-two cetacean species have been recorded as prey, from observing orcas' feeding activity, examining the stomach contents of dead orcas, and seeing scars on the bodies of surviving prey animals. Groups even attack larger cetaceans such as minke whales, grey whales,[116] and, rarely, sperm whales or blue whales.[34][117][118][119] It has been hypothesized that predation by orcas on whale calves in high-productivity, high-latitude areas is the reason for great whale migrations during breeding season to low-productivity tropical waters where orcas are scarcer.[120]
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Hunting a large whale usually takes several hours. Killer whales generally attack young or weak animals; however, a group of five or more may attack a healthy adult. When hunting a young whale, a group chases it and its mother to exhaustion. Eventually, they separate the pair and surround the calf, drowning it by keeping it from surfacing. Pods of female sperm whales sometimes protect themselves by forming a protective circle around their calves with their flukes facing outwards, using them to repel the attackers.[121] Rarely, large killer whale pods can overwhelm even adult female sperm whales. Adult bull sperm whales, which are large, powerful and aggressive when threatened, and fully grown adult blue whales, which are possibly too large to overwhelm, are not believed to be prey for killer whales.[122]
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Prior to the advent of industrial whaling, great whales may have been the major food source for killer whales. The introduction of modern whaling techniques may have aided killer whales by the sound of exploding harpoons indicating availability of prey to scavenge, and compressed air inflation of whale carcasses causing them to float, thus exposing them to scavenging. However, the devastation of great whale populations by unfettered whaling has possibly reduced their availability for killer whales, and caused them to expand their consumption of smaller marine mammals, thus contributing to the decline of these as well.[120]
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Other marine mammal prey species include nearly 20 species of seal, sea lion and fur seal. Walruses and sea otters are less frequently taken. Often, to avoid injury, killer whales disable their prey before killing and eating it. This may involve throwing it in the air, slapping it with their tails, ramming it, or breaching and landing on it.[123] In the Aleutian Islands, a decline in sea otter populations in the 1990s was controversially attributed by some scientists to killer whale predation, although with no direct evidence.[124] The decline of sea otters followed a decline in harbour seal and Steller sea lion populations, the killer whale's preferred prey,[a][126] which in turn may be substitutes for their original prey, now decimated by industrial whaling.[127][128][129]
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In steeply banked beaches off Península Valdés, Argentina, and the Crozet Islands, killer whales feed on South American sea lions and southern elephant seals in shallow water, even beaching temporarily to grab prey before wriggling back to the sea. Beaching, usually fatal to cetaceans, is not an instinctive behaviour, and can require years of practice for the young.[130] Killer whales can then release the animal near juvenile whales, allowing the younger whales to practice the difficult capture technique on the now-weakened prey.[123][131] "Wave-hunting" killer whales spy-hop to locate Weddell seals, crabeater seals, leopard seals, and penguins resting on ice floes, and then swim in groups to create waves that wash over the floe. This washes the prey into the water, where other killer whales lie in wait.[49][132][133]
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Killer whales have also been observed preying on terrestrial mammals, such as deer swimming between islands off the northwest coast of North America.[125] Killer whale cannibalism has also been reported based on analysis of stomach contents, but this is likely to be the result of scavenging remains dumped by whalers.[134] One killer whale was also attacked by its companions after being shot.[34] Although resident killer whales have never been observed to eat other marine mammals, they occasionally harass and kill porpoises and seals for no apparent reason.[135]
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Killer whales in many areas may prey on cormorants and gulls.[136] A captive killer whale at Marineland of Canada discovered it could regurgitate fish onto the surface, attracting sea gulls, and then eat the birds. Four others then learned to copy the behaviour.[137]
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Day-to-day killer whale behaviour generally consists of foraging, travelling, resting and socializing. Killer whales frequently engage in surface behaviour such as breaching (jumping completely out of the water) and tail-slapping. These activities may have a variety of purposes, such as courtship, communication, dislodging parasites, or play. Spyhopping is a behaviour in which a whale holds its head above water to view its surroundings.[138] Resident killer whales swim alongside porpoises and other dolphins.[139]
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Killer whales are notable for their complex societies. Only elephants and higher primates live in comparably complex social structures.[140] Due to orcas' complex social bonds, many marine experts have concerns about how humane it is to keep them in captivity.[141]
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Resident killer whales in the eastern North Pacific live in particularly complex and stable social groups. Unlike any other known mammal social structure, resident whales live with their mothers for their entire lives. These family groups are based on matrilines consisting of the eldest female (matriarch) and her sons and daughters, and the descendants of her daughters, etc. The average size of a matriline is 5.5 animals.[142] Because females can reach age 90, as many as four generations travel together. These matrilineal groups are highly stable. Individuals separate for only a few hours at a time, to mate or forage. With one exception, a killer whale named Luna, no permanent separation of an individual from a resident matriline has been recorded.[142]
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Closely related matrilines form loose aggregations called pods, usually consisting of one to four matrilines. Unlike matrilines, pods may separate for weeks or months at a time.[142] DNA testing indicates resident males nearly always mate with females from other pods.[143] Clans, the next level of resident social structure, are composed of pods with similar dialects, and common but older maternal heritage. Clan ranges overlap, mingling pods from different clans.[142] The final association layer, perhaps more arbitrarily defined than the familial groupings, is called the community, and is defined as a set of clans that regularly commingle. Clans within a community do not share vocal patterns.[b]
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Transient pods are smaller than resident pods, typically consisting of an adult female and one or two of her offspring. Males typically maintain stronger relationships with their mothers than other females. These bonds can extend well into adulthood. Unlike residents, extended or permanent separation of transient offspring from natal matrilines is common, with juveniles and adults of both sexes participating. Some males become "rovers" and do not form long-term associations, occasionally joining groups that contain reproductive females.[144] As in resident clans, transient community members share an acoustic repertoire, although regional differences in vocalizations have been noted.[145]
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Like all cetaceans, killer whales depend heavily on underwater sound for orientation, feeding, and communication. They produce three categories of sounds: clicks, whistles, and pulsed calls. Clicks are believed to be used primarily for navigation and discriminating prey and other objects in the surrounding environment, but are also commonly heard during social interactions.[30]
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Northeast Pacific resident groups tend to be much more vocal than transient groups in the same waters.[146] Residents feed primarily on Chinook and chum, which are insensitive to killer whale calls (inferred from the audiogram of Atlantic salmon). In contrast, the marine mammal prey of transients hear whale calls well. Transients are typically silent.[146] They sometimes use a single click (called a cryptic click) rather than the long train of clicks observed in other populations. Residents are silent only when resting.
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All members of a resident pod use similar calls, known collectively as a dialect. Dialects are composed of specific numbers and types of discrete, repetitive calls. They are complex and stable over time.[147] Call patterns and structure are distinctive within matrilines.[148] Newborns produce calls similar to their mothers, but have a more limited repertoire.[145] Individuals likely learn their dialect through contact with pod members.[149] Family-specific calls have been observed more frequently in the days following a calf's birth, which may help the calf learn them.[150] Dialects are probably an important means of maintaining group identity and cohesiveness. Similarity in dialects likely reflects the degree of relatedness between pods, with variation growing over time.[151] When pods meet, dominant call types decrease and subset call types increase. The use of both call types is called biphonation. The increased subset call types may be the distinguishing factor between pods and inter-pod relations.[148]
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Dialects also distinguish types. Resident dialects contain seven to 17 (mean = 11) distinctive call types. All members of the North American west coast transient community express the same basic dialect, although minor regional variation in call types is evident. Preliminary research indicates offshore killer whales have group-specific dialects unlike those of residents and transients.[151]
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Norwegian and Icelandic herring-eating orcas appear to have different vocalizations for activities like hunting.[152] A population that live in McMurdo Sound, Antarctica have 28 complex burst-pulse and whistle calls.[153]
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Killer whales have the second-heaviest brains among marine mammals[154] (after sperm whales, which have the largest brain of any animal). They can be trained in captivity and are often described as intelligent,[155][156] although defining and measuring "intelligence" is difficult in a species whose environment and behavioural strategies are very different from those of humans.[156]
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Killer whales imitate others, and seem to deliberately teach skills to their kin. Off the Crozet Islands, mothers push their calves onto the beach, waiting to pull the youngster back if needed.[123][131]
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People who have interacted closely with killer whales offer numerous anecdotes demonstrating the whales' curiosity, playfulness, and ability to solve problems. Alaskan killer whales have not only learned how to steal fish from longlines, but have also overcome a variety of techniques designed to stop them, such as the use of unbaited lines as decoys.[157] Once, fishermen placed their boats several miles apart, taking turns retrieving small amounts of their catch, in the hope that the whales would not have enough time to move between boats to steal the catch as it was being retrieved. A researcher described what happened next:
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It worked really well for a while. Then the whales split into two groups. It didn't even take them an hour to figure it out. They were so thrilled when they figured out what was going on, that we were playing games. They were breaching by the boats.
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In other anecdotes, researchers describe incidents in which wild killer whales playfully tease humans by repeatedly moving objects the humans are trying to reach,[158] or suddenly start to toss around a chunk of ice after a human throws a snowball.[159]
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The killer whale's use of dialects and the passing of other learned behaviours from generation to generation have been described as a form of animal culture.[160]
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The complex and stable vocal and behavioural cultures of sympatric groups of killer whales (Orcinus orca) appear to have no parallel outside humans and represent an independent evolution of cultural faculties.[161]
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(Two species or populations are considered sympatric when they live in the same geographic area and thus regularly encounter one another.)
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Female killer whales begin to mature at around the age of 10 and reach peak fertility around 20,[162] experiencing periods of polyestrous cycling separated by non-cycling periods of three to 16 months. Females can often breed until age 40, followed by a rapid decrease in fertility.[162] As such, orcas are among the few animals that undergo menopause and live for decades after they have finished breeding.[163][164] The lifespans of wild females average 50 to 80 years.[165] Some are claimed to have lived substantially longer: Granny (J2) was estimated by some researchers to have been as old as 105 years at the time of her death, though a biopsy sample indicated her age as 65 to 80 years.[166][167][168] Orcas held in captivity tend to live less than those in the wild, although this is subject to scientific debate.[165][169][170]
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To avoid inbreeding, males mate with females from other pods. Gestation varies from 15 to 18 months. [171] Mothers usually calve a single offspring about once every five years. In resident pods, births occur at any time of year, although winter is the most common. Mortality is extremely high during the first seven months of life, when 37–50% of all calves die.[172] Weaning begins at about 12 months of age, and is complete by two years. According to observations in several regions, all male and female pod members participate in the care of the young.[140]
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Males sexually mature at the age of 15, but do not typically reproduce until age 21. Wild males live around 29 years on average, with a maximum of about 60 years.[166] One male, known as Old Tom, was reportedly spotted every winter between the 1840s and 1930 off New South Wales, Australia. This would have made him up to 90 years old. Examination of his teeth indicated he died around age 35,[173] but this method of age determination is now believed to be inaccurate for older animals.[174] One male known to researchers in the Pacific Northwest (identified as J1) was estimated to have been 59 years old when he died in 2010.[175] Killer whales are unique among cetaceans, as their caudal sections elongate with age, making their heads relatively shorter.[60]
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Infanticide, once thought to occur only in captive killer whales, was observed in wild populations by researchers off British Columbia on December 2, 2016. In this incident, an adult male killed the calf of a female within the same pod, with his mother also joining in the assault. It is theorized that the male killed the young calf in order to mate with its mother (something that occurs in other carnivore species), while the male's mother supported the breeding opportunity for her son. The attack ended when the calf's mother struck and injured the attacking male. Such behaviour matches that of many smaller dolphin species, such as the bottlenose dolphin.[176]
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In 2008, the IUCN (International Union for Conservation of Nature) changed its assessment of the killer whale's conservation status from conservation dependent to data deficient, recognizing that one or more killer whale types may actually be separate, endangered species.[3] Depletion of prey species, pollution, large-scale oil spills, and habitat disturbance caused by noise and conflicts with boats are the most significant worldwide threats.[3] In January 2020, the first killer whale in England and Wales since 2001 was found dead with a large fragment of plastic in its stomach.[177]
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Like other animals at the highest trophic levels, the killer whale is particularly at risk of poisoning from bioaccumulation of toxins, including Polychlorinated biphenyls (PCBs).[178] European harbour seals have problems in reproductive and immune functions associated with high levels of PCBs and related contaminants, and a survey off the Washington coast found PCB levels in killer whales were higher than levels that had caused health problems in harbour seals.[178] Blubber samples in the Norwegian Arctic show higher levels of PCBs, pesticides and brominated flame-retardants than in polar bears. When food is scarce, killer whales metabolize blubber for energy, which increases pollutant concentrations in their blood.
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In the Pacific Northwest, wild salmon stocks, a main resident food source, have declined dramatically in recent years.[3] In the Puget Sound region only 75 whales remain with few births over the last few years.[179] On the west coast of Alaska and the Aleutian Islands, seal and sea lion populations have also substantially declined.[180]
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In 2005, the United States government listed the southern resident community as an endangered population under the Endangered Species Act.[30] This community comprises three pods which live mostly in the Georgia and Haro Straits and Puget Sound in British Columbia and Washington. They do not breed outside of their community, which was once estimated at around 200 animals and later shrank to around 90.[181] In October 2008, the annual survey revealed seven were missing and presumed dead, reducing the count to 83.[182] This is potentially the largest decline in the population in the past 10 years. These deaths can be attributed to declines in Chinook salmon.[182]
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Scientist Ken Balcomb has extensively studied killer whales since 1976; he is the research biologist responsible for discovering U.S. Navy sonar may harm killer whales. He studied killer whales from the Center for Whale Research, located in Friday Harbor, Washington.[183] He was also able to study killer whales from "his home porch perched above Puget Sound, where the animals hunt and play in summer months".[183] In May 2003, Balcomb (along with other whale watchers near the Puget Sound coastline) noticed uncharacteristic behaviour displayed by the killer whales. The whales seemed "agitated and were moving haphazardly, attempting to lift their heads free of the water" to escape the sound of the sonars.[183] "Balcomb confirmed at the time that strange underwater pinging noises detected with underwater microphones were sonar. The sound originated from a U.S. Navy frigate 12 miles (19 kilometres) distant, Balcomb said."[183] The impact of sonar waves on killer whales is potentially life-threatening. Three years prior to Balcomb's discovery, research in the Bahamas showed 14 beaked whales washed up on the shore. These whales were beached on the day U.S. Navy destroyers were activated into sonar exercise.[183] Of the 14 whales beached, six of them died. These six dead whales were studied, and CAT scans of two of the whale heads showed hemorrhaging around the brain and the ears, which is consistent with decompression sickness.[183]
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Another conservation concern was made public in September 2008 when the Canadian government decided it was not necessary to enforce further protections (including the Species at Risk Act in place to protect endangered animals along their habitats) for killer whales aside from the laws already in place. In response to this decision, six environmental groups sued the federal government, claiming killer whales were facing many threats on the British Columbia Coast and the federal government did nothing to protect them from these threats.[184] A legal and scientific nonprofit organization, Ecojustice, led the lawsuit and represented the David Suzuki Foundation, Environmental Defence, Greenpeace Canada, International Fund for Animal Welfare, the Raincoast Conservation Foundation, and the Wilderness Committee.[184] Many scientists involved in this lawsuit, including Bill Wareham, a marine scientist with the David Suzuki Foundation, noted increased boat traffic, water toxic wastes, and low salmon population as major threats, putting approximately 87 killer whales[184] on the British Columbia Coast in danger.
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Underwater noise from shipping, drilling, and other human activities is a significant concern in some key killer whale habitats, including Johnstone Strait and Haro Strait.[185] In the mid-1990s, loud underwater noises from salmon farms were used to deter seals. Killer whales also avoided the surrounding waters.[186] High-intensity sonar used by the Navy disturbs killer whales along with other marine mammals.[187] Killer whales are popular with whale watchers, which may stress the whales and alter their behaviour, particularly if boats approach too closely or block their lines of travel.[188]
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The Exxon Valdez oil spill adversely affected killer whales in Prince William Sound and Alaska's Kenai Fjords region. Eleven members (about half) of one resident pod disappeared in the following year. The spill damaged salmon and other prey populations, which in turn damaged local killer whales. By 2009, scientists estimated the AT1 transient population (considered part of a larger population of 346 transients), numbered only seven individuals and had not reproduced since the spill. This population is expected to die out.[189][190]
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A 2018 study published in Science found that global killer whale populations are poised to dramatically decline due to exposure to toxic chemical and PCB pollution.[191]
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The indigenous peoples of the Pacific Northwest Coast feature killer whales throughout their art, history, spirituality and religion. The Haida regarded killer whales as the most powerful animals in the ocean, and their mythology tells of killer whales living in houses and towns under the sea. According to these myths, they took on human form when submerged, and humans who drowned went to live with them.[192] For the Kwakwaka'wakw, the killer whale was regarded as the ruler of the undersea world, with sea lions for slaves and dolphins for warriors.[192] In Nuu-chah-nulth and Kwakwaka'wakw mythology, killer whales may embody the souls of deceased chiefs.[192] The Tlingit of southeastern Alaska regarded the killer whale as custodian of the sea and a benefactor of humans.[193]
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The Maritime Archaic people of Newfoundland also had great respect for killer whales, as evidenced by stone carvings found in a 4,000-year-old burial at the Port au Choix Archaeological Site.[194][195]
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In the tales and beliefs of the Siberian Yupik people, killer whales are said to appear as wolves in winter, and wolves as killer whales in summer.[196][197][198][199] Killer whales are believed to assist their hunters in driving walrus.[200] Reverence is expressed in several forms: the boat represents the animal, and a wooden carving hung from the hunter's belt.[198] Small sacrifices such as tobacco or meat are strewn into the sea for them.[200][199]
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Indigenous Ainu tribe often referred killer whales in their folklore and myth as Repun Kamuy (God of Sea/Offshore) to bring fortunes (whales) to the coasts, and there had been traditional funerals for stranded or deceased orcas akin to funerals for other animals such as brown bears.[201]
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In Western cultures, killer whales were historically feared as dangerous, savage predators.[202] The first written description of a killer whale was given by Pliny the Elder circa AD 70, who wrote, "Orcas (the appearance of which no image can express, other than an enormous mass of savage flesh with teeth) are the enemy of [other kinds of whale]... they charge and pierce them like warships ramming."[203]
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Of the very few confirmed attacks on humans by wild killer whales, none have been fatal.[204] In one instance, killer whales tried to tip ice floes on which a dog team and photographer of the Terra Nova Expedition were standing.[205] The sled dogs' barking is speculated to have sounded enough like seal calls to trigger the killer whale's hunting curiosity. In the 1970s, a surfer in California was bitten, and in 2005, a boy in Alaska who was splashing in a region frequented by harbour seals was bumped by a killer whale that apparently misidentified him as prey.[206] Unlike wild killer whales, captive killer whales have made nearly two dozen attacks on humans since the 1970s, some of which have been fatal.[207][208]
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Competition with fishermen also led to killer whales being regarded as pests. In the waters of the Pacific Northwest and Iceland, the shooting of killer whales was accepted and even encouraged by governments.[202] As an indication of the intensity of shooting that occurred until fairly recently, about 25% of the killer whales captured in Puget Sound for aquarium through 1970 bore bullet scars.[209] The U.S. Navy claimed to have deliberately killed hundreds of killer whales in Icelandic waters in 1956 with machine guns, rockets, and depth charges.[210][211]
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Western attitudes towards killer whales have changed dramatically in recent decades. In the mid-1960s and early 1970s, killer whales came to much greater public and scientific awareness, starting with the first live-capture and display of a killer whale known as Moby Doll, a resident harpooned off Saturna Island in 1964.[202] So little was known at the time, it was nearly two months before the whale's keepers discovered what food (fish) it was willing to eat. To the surprise of those who saw him, Moby Doll was a docile, non-aggressive whale who made no attempts to attack humans.[212]
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Between 1964 and 1976, 50 killer whales from the Pacific Northwest were captured for display in aquaria, and public interest in the animals grew. In the 1970s, research pioneered by Michael Bigg led to the discovery of the species' complex social structure, its use of vocal communication, and its extraordinarily stable mother–offspring bonds. Through photo-identification techniques, individuals were named and tracked over decades.[213]
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Bigg's techniques also revealed the Pacific Northwest population was in the low hundreds rather than the thousands that had been previously assumed.[202] The southern resident community alone had lost 48 of its members to captivity; by 1976, only 80 remained.[214] In the Pacific Northwest, the species that had unthinkingly been targeted became a cultural icon within a few decades.[181]
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The public's growing appreciation also led to growing opposition to whale–keeping in aquarium. Only one whale has been taken in North American waters since 1976. In recent years, the extent of the public's interest in killer whales has manifested itself in several high-profile efforts surrounding individuals. Following the success of the 1993 film Free Willy, the movie's captive star Keiko was returned to the coast of his native Iceland in 2002. The director of the International Marine Mammal Project for the Earth Island Institute, David Phillips, led the efforts to return Keiko to the Iceland waters.[215] Keiko however did not adapt to the harsh climate of the Arctic Ocean, and died a year into his release after contracting pneumonia, at the age of 27.[216] In 2002, the orphan Springer was discovered in Puget Sound, Washington. She became the first whale to be successfully reintegrated into a wild pod after human intervention, crystallizing decades of research into the vocal behaviour and social structure of the region's killer whales.[217] The saving of Springer raised hopes that another young killer whale named Luna, which had become separated from his pod, could be returned to it. However, his case was marked by controversy about whether and how to intervene, and in 2006, Luna was killed by a boat propeller.[218]
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The earlier of known records of commercial hunting of killer whales date to the 18th century in Japan. During the 19th and early 20th centuries, the global whaling industry caught immense numbers of baleen and sperm whales, but largely ignored killer whales because of their limited amounts of recoverable oil, their smaller populations, and the difficulty of taking them.[143] Once the stocks of larger species were depleted, killer whales were targeted by commercial whalers in the mid-20th century. Between 1954 and 1997, Japan took 1,178 killer whales (although the Ministry of the Environment claims that there had been domestic catches of about 1,600 whales between late 1940s to 1960s[219]) and Norway took 987.[220] Extensive hunting of killer whales, including an Antarctic catch of 916 in 1979–80 alone, prompted the International Whaling Commission to recommend a ban on commercial hunting of the species pending further research.[220] Today, no country carries out a substantial hunt, although Indonesia and Greenland permit small subsistence hunts (see Aboriginal whaling). Other than commercial hunts, killer whales were hunted along Japanese coasts out of public concern for potential conflicts with fisheries. Such cases include a semi-resident male-female pair in Akashi Strait and Harimanada being killed in the Seto Inland Sea in 1957,[221][222] the killing of five whales from a pod of 11 members that swam into Tokyo Bay in 1970,[223] and a catch record in southern Taiwan in the 1990s.[79][224]
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Killer whales have helped humans hunting other whales.[225] One well-known example was the killer whales of Eden, Australia, including the male known as Old Tom. Whalers more often considered them a nuisance, however, as orcas would gather to scavenge meat from the whalers' catch.[225] Some populations, such as in Alaska's Prince William Sound, may have been reduced significantly by whalers shooting them in retaliation.[16]
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Whale watching continues to increase in popularity, but may have some problematic impacts on killer whales. Exposure to exhaust gasses from large amounts of vessel traffic are causing concern for the overall health of the 75 remaining southern resident killer whales (SRKWs) left as of early 2019.[226] This population is followed by approximately 20 vessels for 12 hours a day during the months May–September.[227] Researchers discovered that these vessels are in the line of sight for these whales for 98–99.5% of daylight hours.[227] With so many vessels, the air quality around these whales deteriorates and impacts their health. Air pollutants that bind with exhaust fumes are responsible for the activation of the cytochrome P450 1A gene family.[227] Researchers have successfully identified this gene in skin biopsies of live whales and also the lungs of deceased whales. A direct correlation between activation of this gene and the air pollutants can not be made because there are other known factors that will induce the same gene. Vessels can have either wet or dry exhaust systems, with wet exhaust systems leaving more pollutants in the water due to various gas solubility. A modelling study determined that the lowest-observed-adverse-effect-level (LOAEL) of exhaust pollutants was about 12% of the human dose.[227]
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As a response to this, in 2017 boats off the British Columbia coast now have a minimum approach distance of 200 metres compared to the previous 100 metres. This new rule complements Washington State's minimum approach zone of 180 metres that has been in effect since 2011. If a whale approaches a vessel it must be placed in neutral until the whale passes. The World Health Organization has set air quality standards in an effort to control the emissions produced by these vessels.[228]
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The killer whale's intelligence, trainability, striking appearance, playfulness in captivity and sheer size have made it a popular exhibit at aquaria and aquatic theme parks. From 1976 to 1997, 55 whales were taken from the wild in Iceland, 19 from Japan, and three from Argentina. These figures exclude animals that died during capture. Live captures fell dramatically in the 1990s, and by 1999, about 40% of the 48 animals on display in the world were captive-born.[229]
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Organizations such as World Animal Protection and the Whale and Dolphin Conservation campaign against the practice of keeping them in captivity. In captivity, they often develop pathologies, such as the dorsal fin collapse seen in 60–90% of captive males. Captives have vastly reduced life expectancies, on average only living into their 20s.[c] In the wild, females who survive infancy live 46 years on average, and up to 70–80 years in rare cases. Wild males who survive infancy live 31 years on average, and up to 50–60 years.[230] Captivity usually bears little resemblance to wild habitat, and captive whales' social groups are foreign to those found in the wild. Critics claim captive life is stressful due to these factors and the requirement to perform circus tricks that are not part of wild killer whale behaviour, see above.[231] Wild killer whales may travel up to 160 kilometres (100 mi) in a day, and critics say the animals are too big and intelligent to be suitable for captivity.[155] Captives occasionally act aggressively towards themselves, their tankmates, or humans, which critics say is a result of stress.[207] Between 1991 and 2010, the bull orca known as Tilikum was involved in the death of three people, and was featured in the critically acclaimed 2013 film Blackfish.[232] Tilikum lived at SeaWorld from 1992 until his death in 2017.[233][234][235][236][237]
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A 2015 study coauthored by staff at SeaWorld and the Minnesota Zoo indicates that there is no significant difference in survivorship between free-ranging and captive killer whales. The authors speculate about the future utility of studying captive populations for the purposes of understanding orca biology and the implications of such research of captive animals in the overall health of both wild and marine park populations.[238]
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As of March 2016, SeaWorld has announced that they will be ending their orca breeding program and their theatrical shows. They previously announced, in November 2015, that the shows would be coming to an end in San Diego but it is now to happen in both Orlando and San Antonio as well.[239]
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A triangle is a polygon with three edges and three vertices. It is one of the basic shapes in geometry. A triangle with vertices A, B, and C is denoted
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+
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△
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A
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9 |
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B
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C
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{\displaystyle \triangle ABC}
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+
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.
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In Euclidean geometry any three points, when non-collinear, determine a unique triangle and simultaneously, a unique plane (i.e. a two-dimensional Euclidean space). In other words, there is only one plane that contains that triangle, and every triangle is contained in some plane. If the entire geometry is only the Euclidean plane, there is only one plane and all triangles are contained in it; however, in higher-dimensional Euclidean spaces, this is no longer true. This article is about triangles in Euclidean geometry, and in particular, the Euclidean plane, except where otherwise noted.
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+
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Triangles can be classified according to the lengths of their sides:
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+
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Hatch marks, also called tick marks, are used in diagrams of triangles and other geometric figures to identify sides of equal lengths. A side can be marked with a pattern of "ticks", short line segments in the form of tally marks; two sides have equal lengths if they are both marked with the same pattern. In a triangle, the pattern is usually no more than 3 ticks. An equilateral triangle has the same pattern on all 3 sides, an isosceles triangle has the same pattern on just 2 sides, and a scalene triangle has different patterns on all sides since no sides are equal. Similarly, patterns of 1, 2, or 3 concentric arcs inside the angles are used to indicate equal angles. An equilateral triangle has the same pattern on all 3 angles, an isosceles triangle has the same pattern on just 2 angles, and a scalene triangle has different patterns on all angles since no angles are equal.
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+
|
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Triangles can also be classified according to their internal angles, measured here in degrees.
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+
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A triangle that has two angles with the same measure also has two sides with the same length, and therefore it is an isosceles triangle. It follows that in a triangle where all angles have the same measure, all three sides have the same length, and such a triangle is therefore equilateral.
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+
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Triangles are assumed to be two-dimensional plane figures, unless the context provides otherwise (see Non-planar triangles, below). In rigorous treatments, a triangle is therefore called a 2-simplex (see also Polytope). Elementary facts about triangles were presented by Euclid in books 1–4 of his Elements, around 300 BC.
|
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+
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The sum of the measures of the interior angles of a triangle in Euclidean space is always 180 degrees.[5] This fact is equivalent to Euclid's parallel postulate. This allows determination of the measure of the third angle of any triangle given the measure of two angles. An exterior angle of a triangle is an angle that is a linear pair (and hence supplementary) to an interior angle. The measure of an exterior angle of a triangle is equal to the sum of the measures of the two interior angles that are not adjacent to it; this is the exterior angle theorem. The sum of the measures of the three exterior angles (one for each vertex) of any triangle is 360 degrees.[note 2]
|
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+
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Two triangles are said to be similar if every angle of one triangle has the same measure as the corresponding angle in the other triangle. The corresponding sides of similar triangles have lengths that are in the same proportion, and this property is also sufficient to establish similarity.
|
32 |
+
|
33 |
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Some basic theorems about similar triangles are:
|
34 |
+
|
35 |
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Two triangles that are congruent have exactly the same size and shape:[note 4] all pairs of corresponding interior angles are equal in measure, and all pairs of corresponding sides have the same length. (This is a total of six equalities, but three are often sufficient to prove congruence.)
|
36 |
+
|
37 |
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Some individually necessary and sufficient conditions for a pair of triangles to be congruent are:
|
38 |
+
|
39 |
+
Some individually sufficient conditions are:
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40 |
+
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An important condition is:
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+
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Using right triangles and the concept of similarity, the trigonometric functions sine and cosine can be defined. These are functions of an angle which are investigated in trigonometry.
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+
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A central theorem is the Pythagorean theorem, which states in any right triangle, the square of the length of the hypotenuse equals the sum of the squares of the lengths of the two other sides. If the hypotenuse has length c, and the legs have lengths a and b, then the theorem states that
|
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+
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The converse is true: if the lengths of the sides of a triangle satisfy the above equation, then the triangle has a right angle opposite side c.
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+
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Some other facts about right triangles:
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+
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For all triangles, angles and sides are related by the law of cosines and law of sines (also called the cosine rule and sine rule).
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+
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The triangle inequality states that the sum of the lengths of any two sides of a triangle must be greater than or equal to the length of the third side. That sum can equal the length of the third side only in the case of a degenerate triangle, one with collinear vertices. It is not possible for that sum to be less than the length of the third side. A triangle with three given positive side lengths exists if and only if those side lengths satisfy the triangle inequality.
|
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+
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55 |
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Three given angles form a non-degenerate triangle (and indeed an infinitude of them) if and only if both of these conditions hold: (a) each of the angles is positive, and (b) the angles sum to 180°. If degenerate triangles are permitted, angles of 0° are permitted.
|
56 |
+
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57 |
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Three positive angles α, β, and γ, each of them less than 180°, are the angles of a triangle if and only if any one of the following conditions holds:
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+
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59 |
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the last equality applying only if none of the angles is 90° (so the tangent function's value is always finite).
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+
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There are thousands of different constructions that find a special point associated with (and often inside) a triangle, satisfying some unique property: see the article Encyclopedia of Triangle Centers for a catalogue of them. Often they are constructed by finding three lines associated in a symmetrical way with the three sides (or vertices) and then proving that the three lines meet in a single point: an important tool for proving the existence of these is Ceva's theorem, which gives a criterion for determining when three such lines are concurrent. Similarly, lines associated with a triangle are often constructed by proving that three symmetrically constructed points are collinear: here Menelaus' theorem gives a useful general criterion. In this section just a few of the most commonly encountered constructions are explained.
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+
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A perpendicular bisector of a side of a triangle is a straight line passing through the midpoint of the side and being perpendicular to it, i.e. forming a right angle with it. The three perpendicular bisectors meet in a single point, the triangle's circumcenter, usually denoted by O; this point is the center of the circumcircle, the circle passing through all three vertices. The diameter of this circle, called the circumdiameter, can be found from the law of sines stated above. The circumcircle's radius is called the circumradius.
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+
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Thales' theorem implies that if the circumcenter is located on a side of the triangle, then the opposite angle is a right one. If the circumcenter is located inside the triangle, then the triangle is acute; if the circumcenter is located outside the triangle, then the triangle is obtuse.
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+
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An altitude of a triangle is a straight line through a vertex and perpendicular to (i.e. forming a right angle with) the opposite side. This opposite side is called the base of the altitude, and the point where the altitude intersects the base (or its extension) is called the foot of the altitude. The length of the altitude is the distance between the base and the vertex. The three altitudes intersect in a single point, called the orthocenter of the triangle, usually denoted by H. The orthocenter lies inside the triangle if and only if the triangle is acute.
|
68 |
+
|
69 |
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An angle bisector of a triangle is a straight line through a vertex which cuts the corresponding angle in half. The three angle bisectors intersect in a single point, the incenter, usually denoted by I, the center of the triangle's incircle. The incircle is the circle which lies inside the triangle and touches all three sides. Its radius is called the inradius. There are three other important circles, the excircles; they lie outside the triangle and touch one side as well as the extensions of the other two. The centers of the in- and excircles form an orthocentric system.
|
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+
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A median of a triangle is a straight line through a vertex and the midpoint of the opposite side, and divides the triangle into two equal areas. The three medians intersect in a single point, the triangle's centroid or geometric barycenter, usually denoted by G. The centroid of a rigid triangular object (cut out of a thin sheet of uniform density) is also its center of mass: the object can be balanced on its centroid in a uniform gravitational field. The centroid cuts every median in the ratio 2:1, i.e. the distance between a vertex and the centroid is twice the distance between the centroid and the midpoint of the opposite side.
|
72 |
+
|
73 |
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The midpoints of the three sides and the feet of the three altitudes all lie on a single circle, the triangle's nine-point circle. The remaining three points for which it is named are the midpoints of the portion of altitude between the vertices and the orthocenter. The radius of the nine-point circle is half that of the circumcircle. It touches the incircle (at the Feuerbach point) and the three excircles.
|
74 |
+
|
75 |
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The orthocenter (blue point), center of the nine-point circle (red), centroid (orange), and circumcenter (green) all lie on a single line, known as Euler's line (red line). The center of the nine-point circle lies at the midpoint between the orthocenter and the circumcenter, and the distance between the centroid and the circumcenter is half that between the centroid and the orthocenter.
|
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+
|
77 |
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The center of the incircle is not in general located on Euler's line.
|
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+
|
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If one reflects a median in the angle bisector that passes through the same vertex, one obtains a symmedian. The three symmedians intersect in a single point, the symmedian point of the triangle.
|
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+
|
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There are various standard methods for calculating the length of a side or the measure of an angle. Certain methods are suited to calculating values in a right-angled triangle; more complex methods may be required in other situations.
|
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+
|
83 |
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In right triangles, the trigonometric ratios of sine, cosine and tangent can be used to find unknown angles and the lengths of unknown sides. The sides of the triangle are known as follows:
|
84 |
+
|
85 |
+
The sine of an angle is the ratio of the length of the opposite side to the length of the hypotenuse. In our case
|
86 |
+
|
87 |
+
This ratio does not depend on the particular right triangle chosen, as long as it contains the angle A, since all those triangles are similar.
|
88 |
+
|
89 |
+
The cosine of an angle is the ratio of the length of the adjacent side to the length of the hypotenuse. In our case
|
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+
|
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+
The tangent of an angle is the ratio of the length of the opposite side to the length of the adjacent side. In our case
|
92 |
+
|
93 |
+
The acronym "SOH-CAH-TOA" is a useful mnemonic for these ratios.
|
94 |
+
|
95 |
+
The inverse trigonometric functions can be used to calculate the internal angles for a right angled triangle with the length of any two sides.
|
96 |
+
|
97 |
+
Arcsin can be used to calculate an angle from the length of the opposite side and the length of the hypotenuse.
|
98 |
+
|
99 |
+
Arccos can be used to calculate an angle from the length of the adjacent side and the length of the hypotenuse.
|
100 |
+
|
101 |
+
Arctan can be used to calculate an angle from the length of the opposite side and the length of the adjacent side.
|
102 |
+
|
103 |
+
In introductory geometry and trigonometry courses, the notation sin−1, cos−1, etc., are often used in place of arcsin, arccos, etc. However, the arcsin, arccos, etc., notation is standard in higher mathematics where trigonometric functions are commonly raised to powers, as this avoids confusion between multiplicative inverse and compositional inverse.
|
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+
|
105 |
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The law of sines, or sine rule,[8] states that the ratio of the length of a side to the sine of its corresponding opposite angle is constant, that is
|
106 |
+
|
107 |
+
This ratio is equal to the diameter of the circumscribed circle of the given triangle. Another interpretation of this theorem is that every triangle with angles α, β and γ is similar to a triangle with side lengths equal to sin α, sin β and sin γ. This triangle can be constructed by first constructing a circle of diameter 1, and inscribing in it two of the angles of the triangle. The length of the sides of that triangle will be sin α, sin β and sin γ. The side whose length is sin α is opposite to the angle whose measure is α, etc.
|
108 |
+
|
109 |
+
The law of cosines, or cosine rule, connects the length of an unknown side of a triangle to the length of the other sides and the angle opposite to the unknown side.[8] As per the law:
|
110 |
+
|
111 |
+
For a triangle with length of sides a, b, c and angles of α, β, γ respectively, given two known lengths of a triangle a and b, and the angle between the two known sides γ (or the angle opposite to the unknown side c), to calculate the third side c, the following formula can be used:
|
112 |
+
|
113 |
+
If the lengths of all three sides of any triangle are known the three angles can be calculated:
|
114 |
+
|
115 |
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The law of tangents, or tangent rule, can be used to find a side or an angle when two sides and an angle or two angles and a side are known. It states that:[9]
|
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+
|
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+
"Solution of triangles" is the main trigonometric problem: to find missing characteristics of a triangle (three angles, the lengths of the three sides etc.) when at least three of these characteristics are given. The triangle can be located on a plane or on a sphere. This problem often occurs in various trigonometric applications, such as geodesy, astronomy, construction, navigation etc.
|
118 |
+
|
119 |
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Calculating the area T of a triangle is an elementary problem encountered often in many different situations. The best known and simplest formula is:
|
120 |
+
|
121 |
+
where b is the length of the base of the triangle, and h is the height or altitude of the triangle. The term "base" denotes any side, and "height" denotes the length of a perpendicular from the vertex opposite the base onto the line containing the base. In 499 CE Aryabhata, used this illustrated method in the Aryabhatiya (section 2.6).[10]
|
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+
|
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+
Although simple, this formula is only useful if the height can be readily found, which is not always the case. For example, the surveyor of a triangular field might find it relatively easy to measure the length of each side, but relatively difficult to construct a 'height'. Various methods may be used in practice, depending on what is known about the triangle. The following is a selection of frequently used formulae for the area of a triangle.[11]
|
124 |
+
|
125 |
+
The height of a triangle can be found through the application of trigonometry.
|
126 |
+
|
127 |
+
Knowing SAS: Using the labels in the image on the right, the altitude is h = a sin
|
128 |
+
|
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+
|
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+
|
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+
γ
|
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+
|
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+
|
134 |
+
{\displaystyle \gamma }
|
135 |
+
|
136 |
+
. Substituting this in the formula
|
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+
|
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+
|
139 |
+
|
140 |
+
T
|
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+
=
|
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+
|
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+
|
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+
1
|
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+
2
|
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+
|
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+
|
148 |
+
b
|
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+
h
|
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+
|
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+
|
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+
{\displaystyle T={\frac {1}{2}}bh}
|
153 |
+
|
154 |
+
derived above, the area of the triangle can be expressed as:
|
155 |
+
|
156 |
+
(where α is the interior angle at A, β is the interior angle at B,
|
157 |
+
|
158 |
+
|
159 |
+
|
160 |
+
γ
|
161 |
+
|
162 |
+
|
163 |
+
{\displaystyle \gamma }
|
164 |
+
|
165 |
+
is the interior angle at C and c is the line AB).
|
166 |
+
|
167 |
+
Furthermore, since sin α = sin (π − α) = sin (β +
|
168 |
+
|
169 |
+
|
170 |
+
|
171 |
+
γ
|
172 |
+
|
173 |
+
|
174 |
+
{\displaystyle \gamma }
|
175 |
+
|
176 |
+
), and similarly for the other two angles:
|
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+
|
178 |
+
Knowing AAS:
|
179 |
+
|
180 |
+
and analogously if the known side is a or c.
|
181 |
+
|
182 |
+
Knowing ASA:[12]
|
183 |
+
|
184 |
+
and analogously if the known side is b or c.
|
185 |
+
|
186 |
+
The shape of the triangle is determined by the lengths of the sides. Therefore, the area can also be derived from the lengths of the sides. By Heron's formula:
|
187 |
+
|
188 |
+
where
|
189 |
+
|
190 |
+
|
191 |
+
|
192 |
+
s
|
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+
=
|
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+
|
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+
|
196 |
+
|
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+
|
198 |
+
a
|
199 |
+
+
|
200 |
+
b
|
201 |
+
+
|
202 |
+
c
|
203 |
+
|
204 |
+
2
|
205 |
+
|
206 |
+
|
207 |
+
|
208 |
+
|
209 |
+
|
210 |
+
{\displaystyle s={\tfrac {a+b+c}{2}}}
|
211 |
+
|
212 |
+
is the semiperimeter, or half of the triangle's perimeter.
|
213 |
+
|
214 |
+
Three other equivalent ways of writing Heron's formula are
|
215 |
+
|
216 |
+
The area of a parallelogram embedded in a three-dimensional Euclidean space can be calculated using vectors. Let vectors AB and AC point respectively from A to B and from A to C. The area of parallelogram ABDC is then
|
217 |
+
|
218 |
+
which is the magnitude of the cross product of vectors AB and AC. The area of triangle ABC is half of this,
|
219 |
+
|
220 |
+
The area of triangle ABC can also be expressed in terms of dot products as follows:
|
221 |
+
|
222 |
+
In two-dimensional Euclidean space, expressing vector AB as a free vector in Cartesian space equal to (x1,y1) and AC as (x2,y2), this can be rewritten as:
|
223 |
+
|
224 |
+
If vertex A is located at the origin (0, 0) of a Cartesian coordinate system and the coordinates of the other two vertices are given by B = (xB, yB) and C = (xC, yC), then the area can be computed as 1⁄2 times the absolute value of the determinant
|
225 |
+
|
226 |
+
For three general vertices, the equation is:
|
227 |
+
|
228 |
+
which can be written as
|
229 |
+
|
230 |
+
If the points are labeled sequentially in the counterclockwise direction, the above determinant expressions are positive and the absolute value signs can be omitted.[13] The above formula is known as the shoelace formula or the surveyor's formula.
|
231 |
+
|
232 |
+
If we locate the vertices in the complex plane and denote them in counterclockwise sequence as a = xA + yAi, b = xB + yBi, and c = xC + yCi, and denote their complex conjugates as
|
233 |
+
|
234 |
+
|
235 |
+
|
236 |
+
|
237 |
+
|
238 |
+
|
239 |
+
a
|
240 |
+
¯
|
241 |
+
|
242 |
+
|
243 |
+
|
244 |
+
|
245 |
+
|
246 |
+
{\displaystyle {\bar {a}}}
|
247 |
+
|
248 |
+
,
|
249 |
+
|
250 |
+
|
251 |
+
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
b
|
256 |
+
¯
|
257 |
+
|
258 |
+
|
259 |
+
|
260 |
+
|
261 |
+
|
262 |
+
{\displaystyle {\bar {b}}}
|
263 |
+
|
264 |
+
, and
|
265 |
+
|
266 |
+
|
267 |
+
|
268 |
+
|
269 |
+
|
270 |
+
|
271 |
+
c
|
272 |
+
¯
|
273 |
+
|
274 |
+
|
275 |
+
|
276 |
+
|
277 |
+
|
278 |
+
{\displaystyle {\bar {c}}}
|
279 |
+
|
280 |
+
, then the formula
|
281 |
+
|
282 |
+
is equivalent to the shoelace formula.
|
283 |
+
|
284 |
+
In three dimensions, the area of a general triangle A = (xA, yA, zA), B = (xB, yB, zB) and C = (xC, yC, zC) is the Pythagorean sum of the areas of the respective projections on the three principal planes (i.e. x = 0, y = 0 and z = 0):
|
285 |
+
|
286 |
+
The area within any closed curve, such as a triangle, is given by the line integral around the curve of the algebraic or signed distance of a point on the curve from an arbitrary oriented straight line L. Points to the right of L as oriented are taken to be at negative distance from L, while the weight for the integral is taken to be the component of arc length parallel to L rather than arc length itself.
|
287 |
+
|
288 |
+
This method is well suited to computation of the area of an arbitrary polygon. Taking L to be the x-axis, the line integral between consecutive vertices (xi,yi) and (xi+1,yi+1) is given by the base times the mean height, namely (xi+1 − xi)(yi + yi+1)/2. The sign of the area is an overall indicator of the direction of traversal, with negative area indicating counterclockwise traversal. The area of a triangle then falls out as the case of a polygon with three sides.
|
289 |
+
|
290 |
+
While the line integral method has in common with other coordinate-based methods the arbitrary choice of a coordinate system, unlike the others it makes no arbitrary choice of vertex of the triangle as origin or of side as base. Furthermore, the choice of coordinate system defined by L commits to only two degrees of freedom rather than the usual three, since the weight is a local distance (e.g. xi+1 − xi in the above) whence the method does not require choosing an axis normal to L.
|
291 |
+
|
292 |
+
When working in polar coordinates it is not necessary to convert to Cartesian coordinates to use line integration, since the line integral between consecutive vertices (ri,θi) and (ri+1,θi+1) of a polygon is given directly by riri+1sin(θi+1 − θi)/2. This is valid for all values of θ, with some decrease in numerical accuracy when |θ| is many orders of magnitude greater than π. With this formulation negative area indicates clockwise traversal, which should be kept in mind when mixing polar and cartesian coordinates. Just as the choice of y-axis (x = 0) is immaterial for line integration in cartesian coordinates, so is the choice of zero heading (θ = 0) immaterial here.
|
293 |
+
|
294 |
+
Three formulas have the same structure as Heron's formula but are expressed in terms of different variables. First, denoting the medians from sides a, b, and c respectively as ma, mb, and mc and their semi-sum (ma + mb + mc)/2 as σ, we have[14]
|
295 |
+
|
296 |
+
Next, denoting the altitudes from sides a, b, and c respectively as ha, hb, and hc, and denoting the semi-sum of the reciprocals of the altitudes as
|
297 |
+
|
298 |
+
|
299 |
+
|
300 |
+
H
|
301 |
+
=
|
302 |
+
(
|
303 |
+
|
304 |
+
h
|
305 |
+
|
306 |
+
a
|
307 |
+
|
308 |
+
|
309 |
+
−
|
310 |
+
1
|
311 |
+
|
312 |
+
|
313 |
+
+
|
314 |
+
|
315 |
+
h
|
316 |
+
|
317 |
+
b
|
318 |
+
|
319 |
+
|
320 |
+
−
|
321 |
+
1
|
322 |
+
|
323 |
+
|
324 |
+
+
|
325 |
+
|
326 |
+
h
|
327 |
+
|
328 |
+
c
|
329 |
+
|
330 |
+
|
331 |
+
−
|
332 |
+
1
|
333 |
+
|
334 |
+
|
335 |
+
)
|
336 |
+
|
337 |
+
/
|
338 |
+
|
339 |
+
2
|
340 |
+
|
341 |
+
|
342 |
+
{\displaystyle H=(h_{a}^{-1}+h_{b}^{-1}+h_{c}^{-1})/2}
|
343 |
+
|
344 |
+
we have[15]
|
345 |
+
|
346 |
+
And denoting the semi-sum of the angles' sines as S = [(sin α) + (sin β) + (sin γ)]/2, we have[16]
|
347 |
+
|
348 |
+
where D is the diameter of the circumcircle:
|
349 |
+
|
350 |
+
|
351 |
+
|
352 |
+
D
|
353 |
+
=
|
354 |
+
|
355 |
+
|
356 |
+
|
357 |
+
a
|
358 |
+
|
359 |
+
sin
|
360 |
+
|
361 |
+
α
|
362 |
+
|
363 |
+
|
364 |
+
|
365 |
+
|
366 |
+
=
|
367 |
+
|
368 |
+
|
369 |
+
|
370 |
+
b
|
371 |
+
|
372 |
+
sin
|
373 |
+
|
374 |
+
β
|
375 |
+
|
376 |
+
|
377 |
+
|
378 |
+
|
379 |
+
=
|
380 |
+
|
381 |
+
|
382 |
+
|
383 |
+
c
|
384 |
+
|
385 |
+
sin
|
386 |
+
|
387 |
+
γ
|
388 |
+
|
389 |
+
|
390 |
+
|
391 |
+
|
392 |
+
.
|
393 |
+
|
394 |
+
|
395 |
+
{\displaystyle D={\tfrac {a}{\sin \alpha }}={\tfrac {b}{\sin \beta }}={\tfrac {c}{\sin \gamma }}.}
|
396 |
+
|
397 |
+
See Pick's theorem for a technique for finding the area of any arbitrary lattice polygon (one drawn on a grid with vertically and horizontally adjacent lattice points at equal distances, and with vertices on lattice points).
|
398 |
+
|
399 |
+
The theorem states:
|
400 |
+
|
401 |
+
where
|
402 |
+
|
403 |
+
|
404 |
+
|
405 |
+
I
|
406 |
+
|
407 |
+
|
408 |
+
{\displaystyle I}
|
409 |
+
|
410 |
+
is the number of internal lattice points and B is the number of lattice points lying on the border of the polygon.
|
411 |
+
|
412 |
+
Numerous other area formulas exist, such as
|
413 |
+
|
414 |
+
where r is the inradius, and s is the semiperimeter (in fact, this formula holds for all tangential polygons), and[17]:Lemma 2
|
415 |
+
|
416 |
+
where
|
417 |
+
|
418 |
+
|
419 |
+
|
420 |
+
|
421 |
+
r
|
422 |
+
|
423 |
+
a
|
424 |
+
|
425 |
+
|
426 |
+
,
|
427 |
+
|
428 |
+
|
429 |
+
r
|
430 |
+
|
431 |
+
b
|
432 |
+
|
433 |
+
|
434 |
+
,
|
435 |
+
|
436 |
+
|
437 |
+
r
|
438 |
+
|
439 |
+
c
|
440 |
+
|
441 |
+
|
442 |
+
|
443 |
+
|
444 |
+
{\displaystyle r_{a},\,r_{b},\,r_{c}}
|
445 |
+
|
446 |
+
are the radii of the excircles tangent to sides a, b, c respectively.
|
447 |
+
|
448 |
+
We also have
|
449 |
+
|
450 |
+
and[18]
|
451 |
+
|
452 |
+
for circumdiameter D; and[19]
|
453 |
+
|
454 |
+
for angle α ≠ 90°.
|
455 |
+
|
456 |
+
The area can also be expressed as[20]
|
457 |
+
|
458 |
+
In 1885, Baker[21] gave a collection of over a hundred distinct area formulas for the triangle. These include:
|
459 |
+
|
460 |
+
for circumradius (radius of the circumcircle) R, and
|
461 |
+
|
462 |
+
The area T of any triangle with perimeter p satisfies
|
463 |
+
|
464 |
+
with equality holding if and only if the triangle is equilateral.[22][23]:657
|
465 |
+
|
466 |
+
Other upper bounds on the area T are given by[24]:p.290
|
467 |
+
|
468 |
+
and
|
469 |
+
|
470 |
+
both again holding if and only if the triangle is equilateral.
|
471 |
+
|
472 |
+
There are infinitely many lines that bisect the area of a triangle.[25] Three of them are the medians, which are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides.
|
473 |
+
|
474 |
+
Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter. There can be one, two, or three of these for any given triangle.
|
475 |
+
|
476 |
+
The formulas in this section are true for all Euclidean triangles.
|
477 |
+
|
478 |
+
The medians and the sides are related by[26]:p.70
|
479 |
+
|
480 |
+
and
|
481 |
+
|
482 |
+
and equivalently for mb and mc.
|
483 |
+
|
484 |
+
For angle A opposite side a, the length of the internal angle bisector is given by[27]
|
485 |
+
|
486 |
+
for semiperimeter s, where the bisector length is measured from the vertex to where it meets the opposite side.
|
487 |
+
|
488 |
+
The interior perpendicular bisectors are given by
|
489 |
+
|
490 |
+
where the sides are
|
491 |
+
|
492 |
+
|
493 |
+
|
494 |
+
a
|
495 |
+
≥
|
496 |
+
b
|
497 |
+
≥
|
498 |
+
c
|
499 |
+
|
500 |
+
|
501 |
+
{\displaystyle a\geq b\geq c}
|
502 |
+
|
503 |
+
and the area is
|
504 |
+
|
505 |
+
|
506 |
+
|
507 |
+
T
|
508 |
+
.
|
509 |
+
|
510 |
+
|
511 |
+
{\displaystyle T.}
|
512 |
+
|
513 |
+
[28]:Thm 2
|
514 |
+
|
515 |
+
The altitude from, for example, the side of length a is
|
516 |
+
|
517 |
+
The following formulas involve the circumradius R and the inradius r:
|
518 |
+
|
519 |
+
where ha etc. are the altitudes to the subscripted sides;[26]:p.79
|
520 |
+
|
521 |
+
and
|
522 |
+
|
523 |
+
The product of two sides of a triangle equals the altitude to the third side times the diameter D of the circumcircle:[26]:p.64
|
524 |
+
|
525 |
+
Suppose two adjacent but non-overlapping triangles share the same side of length f and share the same circumcircle, so that the side of length f is a chord of the circumcircle and the triangles have side lengths (a, b, f) and (c, d, f), with the two triangles together forming a cyclic quadrilateral with side lengths in sequence (a, b, c, d). Then[29]:84
|
526 |
+
|
527 |
+
Let G be the centroid of a triangle with vertices A, B, and C, and let P be any interior point. Then the distances between the points are related by[29]:174
|
528 |
+
|
529 |
+
The sum of the squares of the triangle's sides equals three times the sum of the squared distances of the centroid from the vertices:
|
530 |
+
|
531 |
+
Let qa, qb, and qc be the distances from the centroid to the sides of lengths a, b, and c. Then[29]:173
|
532 |
+
|
533 |
+
and
|
534 |
+
|
535 |
+
for area T.
|
536 |
+
|
537 |
+
Carnot's theorem states that the sum of the distances from the circumcenter to the three sides equals the sum of the circumradius and the inradius.[26]:p.83 Here a segment's length is considered to be negative if and only if the segment lies entirely outside the triangle. This method is especially useful for deducing the properties of more abstract forms of triangles, such as the ones induced by Lie algebras, that otherwise have the same properties as usual triangles.
|
538 |
+
|
539 |
+
Euler's theorem states that the distance d between the circumcenter and the incenter is given by[26]:p.85
|
540 |
+
|
541 |
+
or equivalently
|
542 |
+
|
543 |
+
where R is the circumradius and r is the inradius. Thus for all triangles R ≥ 2r, with equality holding for equilateral triangles.
|
544 |
+
|
545 |
+
If we denote that the orthocenter divides one altitude into segments of lengths u and v, another altitude into segment lengths w and x, and the third altitude into segment lengths y and z, then uv = wx = yz.[26]:p.94
|
546 |
+
|
547 |
+
The distance from a side to the circumcenter equals half the distance from the opposite vertex to the orthocenter.[26]:p.99
|
548 |
+
|
549 |
+
The sum of the squares of the distances from the vertices to the orthocenter H plus the sum of the squares of the sides equals twelve times the square of the circumradius:[26]:p.102
|
550 |
+
|
551 |
+
In addition to the law of sines, the law of cosines, the law of tangents, and the trigonometric existence conditions given earlier, for any triangle
|
552 |
+
|
553 |
+
Morley's trisector theorem states that in any triangle, the three points of intersection of the adjacent angle trisectors form an equilateral triangle, called the Morley triangle.
|
554 |
+
|
555 |
+
As discussed above, every triangle has a unique inscribed circle (incircle) that is interior to the triangle and tangent to all three sides.
|
556 |
+
|
557 |
+
Every triangle has a unique Steiner inellipse which is interior to the triangle and tangent at the midpoints of the sides. Marden's theorem shows how to find the foci of this ellipse.[31] This ellipse has the greatest area of any ellipse tangent to all three sides of the triangle.
|
558 |
+
|
559 |
+
The Mandart inellipse of a triangle is the ellipse inscribed within the triangle tangent to its sides at the contact points of its excircles.
|
560 |
+
|
561 |
+
For any ellipse inscribed in a triangle ABC, let the foci be P and Q. Then[32]
|
562 |
+
|
563 |
+
Every convex polygon with area T can be inscribed in a triangle of area at most equal to 2T. Equality holds (exclusively) for a parallelogram.[33]
|
564 |
+
|
565 |
+
The Lemoine hexagon is a cyclic hexagon with vertices given by the six intersections of the sides of a triangle with the three lines that are parallel to the sides and that pass through its symmedian point. In either its simple form or its self-intersecting form, the Lemoine hexagon is interior to the triangle with two vertices on each side of the triangle.
|
566 |
+
|
567 |
+
Every acute triangle has three inscribed squares (squares in its interior such that all four of a square's vertices lie on a side of the triangle, so two of them lie on the same side and hence one side of the square coincides with part of a side of the triangle). In a right triangle two of the squares coincide and have a vertex at the triangle's right angle, so a right triangle has only two distinct inscribed squares. An obtuse triangle has only one inscribed square, with a side coinciding with part of the triangle's longest side. Within a given triangle, a longer common side is associated with a smaller inscribed square. If an inscribed square has side of length qa and the triangle has a side of length a, part of which side coincides with a side of the square, then qa, a, the altitude ha from the side a, and the triangle's area T are related according to[34][35]
|
568 |
+
|
569 |
+
The largest possible ratio of the area of the inscribed square to the area of the triangle is 1/2, which occurs when a2 = 2T, q = a/2, and the altitude of the triangle from the base of length a is equal to a. The smallest possible ratio of the side of one inscribed square to the side of another in the same non-obtuse triangle is
|
570 |
+
|
571 |
+
|
572 |
+
|
573 |
+
2
|
574 |
+
|
575 |
+
|
576 |
+
2
|
577 |
+
|
578 |
+
|
579 |
+
|
580 |
+
/
|
581 |
+
|
582 |
+
3
|
583 |
+
=
|
584 |
+
0.94....
|
585 |
+
|
586 |
+
|
587 |
+
{\displaystyle 2{\sqrt {2}}/3=0.94....}
|
588 |
+
|
589 |
+
[35] Both of these extreme cases occur for the isosceles right triangle.
|
590 |
+
|
591 |
+
From an interior point in a reference triangle, the nearest points on the three sides serve as the vertices of the pedal triangle of that point. If the interior point is the circumcenter of the reference triangle, the vertices of the pedal triangle are the midpoints of the reference triangle's sides, and so the pedal triangle is called the midpoint triangle or medial triangle. The midpoint triangle subdivides the reference triangle into four congruent triangles which are similar to the reference triangle.
|
592 |
+
|
593 |
+
The Gergonne triangle or intouch triangle of a reference triangle has its vertices at the three points of tangency of the reference triangle's sides with its incircle. The extouch triangle of a reference triangle has its vertices at the points of tangency of the reference triangle's excircles with its sides (not extended).
|
594 |
+
|
595 |
+
The tangential triangle of a reference triangle (other than a right triangle) is the triangle whose sides are on the tangent lines to the reference triangle's circumcircle at its vertices.
|
596 |
+
|
597 |
+
As mentioned above, every triangle has a unique circumcircle, a circle passing through all three vertices, whose center is the intersection of the perpendicular bisectors of the triangle's sides.
|
598 |
+
|
599 |
+
Further, every triangle has a unique Steiner circumellipse, which passes through the triangle's vertices and has its center at the triangle's centroid. Of all ellipses going through the triangle's vertices, it has the smallest area.
|
600 |
+
|
601 |
+
The Kiepert hyperbola is the unique conic which passes through the triangle's three vertices, its centroid, and its circumcenter.
|
602 |
+
|
603 |
+
Of all triangles contained in a given convex polygon, there exists a triangle with maximal area whose vertices are all vertices of the given polygon.[36]
|
604 |
+
|
605 |
+
One way to identify locations of points in (or outside) a triangle is to place the triangle in an arbitrary location and orientation in the Cartesian plane, and to use Cartesian coordinates. While convenient for many purposes, this approach has the disadvantage of all points' coordinate values being dependent on the arbitrary placement in the plane.
|
606 |
+
|
607 |
+
Two systems avoid that feature, so that the coordinates of a point are not affected by moving the triangle, rotating it, or reflecting it as in a mirror, any of which give a congruent triangle, or even by rescaling it to give a similar triangle:
|
608 |
+
|
609 |
+
A non-planar triangle is a triangle which is not contained in a (flat) plane. Some examples of non-planar triangles in non-Euclidean geometries are spherical triangles in spherical geometry and hyperbolic triangles in hyperbolic geometry.
|
610 |
+
|
611 |
+
While the measures of the internal angles in planar triangles always sum to 180°, a hyperbolic triangle has measures of angles that sum to less than 180°, and a spherical triangle has measures of angles that sum to more than 180°. A hyperbolic triangle can be obtained by drawing on a negatively curved surface, such as a saddle surface, and a spherical triangle can be obtained by drawing on a positively curved surface such as a sphere. Thus, if one draws a giant triangle on the surface of the Earth, one will find that the sum of the measures of its angles is greater than 180°; in fact it will be between 180° and 540°.[37] In particular it is possible to draw a triangle on a sphere such that the measure of each of its internal angles is equal to 90°, adding up to a total of 270°.
|
612 |
+
|
613 |
+
Specifically, on a sphere the sum of the angles of a triangle is
|
614 |
+
|
615 |
+
where f is the fraction of the sphere's area which is enclosed by the triangle. For example, suppose that we draw a triangle on the Earth's surface with vertices at the North Pole, at a point on the equator at 0° longitude, and a point on the equator at 90° West longitude. The great circle line between the latter two points is the equator, and the great circle line between either of those points and the North Pole is a line of longitude; so there are right angles at the two points on the equator. Moreover, the angle at the North Pole is also 90° because the other two vertices differ by 90° of longitude. So the sum of the angles in this triangle is 90° + 90° + 90° = 270°. The triangle encloses 1/4 of the northern hemisphere (90°/360° as viewed from the North Pole) and therefore 1/8 of the Earth's surface, so in the formula f = 1/8; thus the formula correctly gives the sum of the triangle's angles as 270°.
|
616 |
+
|
617 |
+
From the above angle sum formula we can also see that the Earth's surface is locally flat: If we draw an arbitrarily small triangle in the neighborhood of one point on the Earth's surface, the fraction f of the Earth's surface which is enclosed by the triangle will be arbitrarily close to zero. In this case the angle sum formula simplifies to 180°, which we know is what Euclidean geometry tells us for triangles on a flat surface.
|
618 |
+
|
619 |
+
Rectangles have been the most popular and common geometric form for buildings since the shape is easy to stack and organize; as a standard, it is easy to design furniture and fixtures to fit inside rectangularly shaped buildings. But triangles, while more difficult to use conceptually, provide a great deal of strength. As computer technology helps architects design creative new buildings, triangular shapes are becoming increasingly prevalent as parts of buildings and as the primary shape for some types of skyscrapers as well as building materials. In Tokyo in 1989, architects had wondered whether it was possible to build a 500-story tower to provide affordable office space for this densely packed city, but with the danger to buildings from earthquakes, architects considered that a triangular shape would be necessary if such a building were to be built.[38]
|
620 |
+
|
621 |
+
In New York City, as Broadway crisscrosses major avenues, the resulting blocks are cut like triangles, and buildings have been built on these shapes; one such building is the triangularly shaped Flatiron Building which real estate people admit has a "warren of awkward spaces that do not easily accommodate modern office furniture" but that has not prevented the structure from becoming a landmark icon.[39] Designers have made houses in Norway using triangular themes.[40] Triangle shapes have appeared in churches[41] as well as public buildings including colleges[42] as well as supports for innovative home designs.[43]
|
622 |
+
|
623 |
+
Triangles are sturdy; while a rectangle can collapse into a parallelogram from pressure to one of its points, triangles have a natural strength which supports structures against lateral pressures. A triangle will not change shape unless its sides are bent or extended or broken or if its joints break; in essence, each of the three sides supports the other two. A rectangle, in contrast, is more dependent on the strength of its joints in a structural sense. Some innovative designers have proposed making bricks not out of rectangles, but with triangular shapes which can be combined in three dimensions.[44] It is likely that triangles will be used increasingly in new ways as architecture increases in complexity. It is important to remember that triangles are strong in terms of rigidity, but while packed in a tessellating arrangement triangles are not as strong as hexagons under compression (hence the prevalence of hexagonal forms in nature). Tessellated triangles still maintain superior strength for cantilevering however, and this is the basis for one of the strongest man made structures, the tetrahedral truss.
|
624 |
+
|
625 |
+
|
626 |
+
|
627 |
+
|
628 |
+
|
en/4322.html.txt
ADDED
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1 |
+
An orthography is a set of conventions for writing a language. It includes norms of spelling, hyphenation, capitalization, word breaks, emphasis, and punctuation.
|
2 |
+
|
3 |
+
Most transnational languages in the modern period have a system of writing, and for most such languages a standard orthography has been developed, often based on a standard variety of the language, and thus exhibiting less dialect variation than the spoken language.[1][2] Sometimes there may be variation in a language's orthography, as between American and British spelling in the case of English orthography. In some languages orthography is regulated by language academies, although for many languages (including English) there are no such authorities, and orthography develops in a more natural way. Even in the latter languages, a significant amount of consensus arises naturally, although a maximum of consistency or standardization occurs only when prescriptively imposed according to style guides.
|
4 |
+
|
5 |
+
The English word orthography dates from the 15th century. It comes from the French orthographie, from Latin orthographia, which derives from Greek ὀρθός orthós, "correct", and γράφειν gráphein, "to write".[3]
|
6 |
+
|
7 |
+
Orthography is largely concerned with matters of spelling, and in particular the relationship between phonemes and graphemes in a language.[4][5] Other elements that may be considered part of orthography include hyphenation, capitalization, word breaks, emphasis, and punctuation.[6] Orthography thus describes or defines the set of symbols used in writing a language, and the rules regarding how to use those symbols.
|
8 |
+
|
9 |
+
Most natural languages developed as oral languages, and writing systems have usually been crafted or adapted as ways of representing the spoken language. The rules for doing this tend to become standardized for a given language, leading to the development of an orthography that is generally considered "correct". In linguistics the term orthography is often used to refer to any method of writing a language, without judgment as to right and wrong, with a scientific understanding that orthographic standardization exists on a spectrum of strength of convention. The original sense of the word, though, implies a dichotomy of correct and incorrect, and the word is still most often used to refer specifically to a thoroughly standardized, prescriptively correct, way of writing a language. A distinction may be made here between etic and emic viewpoints: the purely descriptive (etic) approach, which simply considers any system that is actually used—and the emic view, which takes account of language users' perceptions of correctness.
|
10 |
+
|
11 |
+
Orthographic units, such as letters of an alphabet, are technically called graphemes. These are a type of abstraction, analogous to the phonemes of spoken languages; different physical forms of written symbols are considered to represent the same grapheme if the differences between them are not significant for meaning. For example, different forms of the letter "b" are all considered to represent a single grapheme in the orthography of, say, English.
|
12 |
+
|
13 |
+
Graphemes or sequences of them are sometimes placed between angle brackets, as in ⟨b⟩ or ⟨back⟩. This distinguishes them from phonemic transcription, which is placed between slashes (/b/, /bæk/), and from phonetic transcription, which is placed between square brackets ([b], [bæk]).
|
14 |
+
|
15 |
+
The writing systems on which orthographies are based can be divided into a number of types, depending on what type of unit each symbol serves to represent. The principal types are logographic (with symbols representing words or morphemes), syllabic (with symbols representing syllables), and alphabetic (with symbols roughly representing phonemes). Many writing systems combine features of more than one of these types, and a number of detailed classifications have been proposed. Japanese is an example of a writing system that can be written using a combination of logographic kanji characters and syllabic hiragana and katakana characters; as with many non-alphabetic languages, alphabetic romaji characters may also be used as needed.[7]
|
16 |
+
|
17 |
+
Orthographies that use alphabets and syllabaries are based on the principle that the written symbols (graphemes) correspond to units of sound of the spoken language: phonemes in the former case, and syllables in the latter. However, in virtually all cases, this correspondence is not exact. Different languages' orthographies offer different degrees of correspondence between spelling and pronunciation. English orthography, French orthography and Danish orthography, for example, are highly irregular, whereas the orthographies of languages such as Russian, German and Spanish represent pronunciation much more faithfully, although the correspondence between letters and phonemes is still not exact. Finnish, Turkish and Serbo-Croatian orthographies are remarkably consistent: approximation of the principle "one letter per sound".
|
18 |
+
|
19 |
+
An orthography in which the correspondences between spelling and pronunciation are highly complex or inconsistent is called a deep orthography (or less formally, the language is said to have irregular spelling). An orthography with relatively simple and consistent correspondences is called shallow (and the language has regular spelling).
|
20 |
+
|
21 |
+
One of the main reasons for which spelling and pronunciation deviate is that sound changes taking place in the spoken language are not always reflected in the orthography, and hence spellings correspond to historical rather than present-day pronunciation. One consequence of this is that many spellings come to reflect a word's morphophonemic structure rather than its purely phonemic structure (for example, the English regular past tense morpheme is consistently spelled -ed in spite of its different pronunciations in various words). This is discussed further at Phonemic orthography § Morphophonemic features.
|
22 |
+
|
23 |
+
The syllabary systems of Japanese (hiragana and katakana) are examples of almost perfectly shallow orthographies—the kana correspond with almost perfect consistency to the spoken syllables, although with a few exceptions where symbols reflect historical or morphophonemic features: notably the use of ぢ ji and づ zu (rather than じ ji and ず zu, their pronunciation in standard Tokyo dialect) when the character is a voicing of an underlying ち or つ (see rendaku), and the use of は, を, and へ to represent the sounds わ, お, and え, as relics of historical kana usage.
|
24 |
+
|
25 |
+
The Korean hangul system was also originally an extremely shallow orthography, but as a representation of the modern language it frequently also reflects morphophonemic features.
|
26 |
+
|
27 |
+
For full discussion of degrees of correspondence between spelling and pronunciation in alphabetic orthographies, including reasons why such correspondence may break down, see Phonemic orthography.
|
28 |
+
|
29 |
+
An orthography based on the principle that symbols correspond to phonemes may, in some cases, lack characters to represent all the phonemes or all the phonemic distinctions in the language. This is called a defective orthography. An example in English is the lack of any indication of stress. Another is the digraph th, which represents two different phonemes (as in then and thin). A more systematic example is that of abjads like the Arabic and Hebrew alphabets, in which the short vowels are normally left unwritten and must be inferred by the reader.
|
30 |
+
|
31 |
+
When an alphabet is borrowed from its original language for use with a new language—as has been done with the Latin alphabet for many languages, or Japanese Katakana for non-Japanese words—it often proves defective in representing the new language's phonemes. Sometimes this problem is addressed by the use of such devices as digraphs (such as sh and ch in English, where pairs of letters represent single sounds), diacritics (like the caron on the letters š and č, which represent those same sounds in Czech), or the addition of completely new symbols (as some languages have introduced the letter w to the Latin alphabet) or of symbols from another alphabet, such as the rune þ in Icelandic.
|
32 |
+
|
33 |
+
After the classical period, Greek developed a lowercase letter system that introduced diacritic marks to enable foreigners to learn pronunciation and in some cases, grammatical features. However, as pronunciation of letters changed over time, the diacritic marks were reduced to representing the stressed syllable. In Modern Greek typesetting, this system has been simplified to only have a single accent to indicate which syllable is stressed.[8]
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1 |
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Green Book
|
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|
5 |
+
Parasite
|
6 |
+
|
7 |
+
The Academy Awards, popularly known as the Oscars,[1] are awards for artistic and technical merit in the film industry. Given annually by the Academy of Motion Picture Arts and Sciences (AMPAS), the awards are an international recognition of excellence in cinematic achievements, as assessed by the Academy's voting membership. The various category winners are awarded a copy of a golden statuette, officially called the "Academy Award of Merit", although more commonly referred to by its nickname, the "Oscar". The statuette depicts a knight rendered in the Art Deco style.
|
8 |
+
|
9 |
+
The award was originally sculpted by George Stanley from a design sketch by Cedric Gibbons.[2] AMPAS first presented it in 1929 at a private dinner hosted by Douglas Fairbanks in the Hollywood Roosevelt Hotel in what would become known as the 1st Academy Awards.[3][4] The Academy Awards ceremony was first broadcast by radio in 1930 and was televised for the first time in 1953. It is the oldest worldwide entertainment awards ceremony and is now televised live worldwide.[5] It is also the oldest of the four major annual American entertainment awards; its equivalents – the Emmy Awards for television, the Tony Awards for theater, and the Grammy Awards for music – are modeled after the Academy Awards.[6] They are widely cited as the most famous and prestigious competitive awards in the field of entertainment.
|
10 |
+
|
11 |
+
The 92nd Academy Awards ceremony, honoring the best films of 2019, was held on February 9, 2020. As with the previous ceremony, there was no host. The ceremony was broadcast on ABC. It took place at the Dolby Theatre in Los Angeles, California, for the 18th consecutive year. This ceremony marked the first time in Oscar history that more than three films had received ten or more nominations.[7] A total of 3,140 Oscar statuettes have been awarded since its inception in 1929.[8]
|
12 |
+
|
13 |
+
The first Academy Awards presentation was held on May 16, 1929, at a private dinner function at the Hollywood Roosevelt Hotel with an audience of about 270 people.[9]
|
14 |
+
|
15 |
+
The post-awards party was held at the Mayfair Hotel.[5] The cost of guest tickets for that night's ceremony was $5 ($74 in 2020 dollars). Fifteen statuettes were awarded, honoring artists, directors and other participants in the film-making industry of the time, for their works during the 1927–28 period. The ceremony ran for 15 minutes.
|
16 |
+
|
17 |
+
Winners were announced to media three months earlier. That was changed for the second ceremony in 1930. Since then, for the rest of the first decade, the results were given to newspapers for publication at 11:00 pm on the night of the awards.[5] This method was used until 1940, when the Los Angeles Times announced the winners before the ceremony began; as a result, the Academy has, since 1941, used a sealed envelope to reveal the names of the winners.[5]
|
18 |
+
|
19 |
+
The first Best Actor awarded was Emil Jannings, for his performances in The Last Command and The Way of All Flesh. He had to return to Europe before the ceremony, so the Academy agreed to give him the prize earlier; this made him the first Academy Award winner in history. At that time, winners were recognized for the entirety of their work done in a certain category during the qualifying period; for example, Jannings received the award for two movies in which he starred during that period, and Janet Gaynor later won a single Oscar for performances in three films. With the fourth ceremony, however, the system changed, and professionals were honored for a specific performance in a single film. For the first six ceremonies, the eligibility period spanned two calendar years.[5]
|
20 |
+
|
21 |
+
At the 29th ceremony, held in 1957, the Best Foreign Language Film category, now known as Best International Feature Film, was introduced. Until then, foreign-language films had been honored with the Special Achievement Award.
|
22 |
+
|
23 |
+
Perhaps the most widely seen streaker in history was 34-year-old Robert Opel, who streaked across the stage of The Dorothy Chandler Pavilion in Los Angeles flashing a peace sign on national US television at the 46th Academy Awards in 1974. Bemused host David Niven quipped, "Isn't it fascinating to think that probably the only laugh that man will ever get in his life is by stripping off and showing his shortcomings?" Later, evidence arose suggesting that Opel's appearance was facilitated as a publicity stunt by the show's producer Jack Haley Jr. Robert Metzler, the show's business manager, believed that the incident had been planned in some way; during the dress rehearsal Niven had asked Metzler's wife to borrow a pen so he could write down the famous line, which was thus not the ad-lib it appeared to be.[10]
|
24 |
+
|
25 |
+
The 74th Academy Awards, held in 2002, presented the first Academy Award for Best Animated Feature.[11]
|
26 |
+
|
27 |
+
Since 1973, all Academy Awards ceremonies have ended with the Academy Award for Best Picture.
|
28 |
+
|
29 |
+
Traditionally, the previous year's winner for Best Actor and Best Supporting Actor present the awards for Best Actress and Best Supporting Actress, while the previous year's winner for Best Actress and Best Supporting Actress present the awards for Best Actor and Best Supporting Actor.
|
30 |
+
|
31 |
+
Parasite became the first foreign-language film to win Best Picture at the February 9, 2020 award ceremony.[12]
|
32 |
+
|
33 |
+
Tom Hanks announced at the 2020 Oscar Ceremony, the opening of the Academy Museum of Motion Pictures on December 14, 2020.[13] The museum development started in 2017 under Kerry Brougher, but is now led by Bill Kramer.[14] The industry curated exhibits will be geared toward the history of motion picture, the art & science of film making, exhibiting trailblazing directors, actors, film makers, sound editors and more, and will house famous artifacts from acclaimed movies like Dorothy's Ruby Red Slippers.
|
34 |
+
|
35 |
+
Because of COVID-19, Academy president David Rubin and CEO Dawn Hudson announced that for the 2021 Oscar Ceremony, streaming movies not shown in theaters would be eligible, though at some point the requirement that movies be shown in theaters would return.[15]
|
36 |
+
|
37 |
+
The best known award is the Academy Award of Merit, more popularly known as the Oscar statuette.[8] Made of gold-plated bronze on a black metal base, it is 13.5 in (34.3 cm) tall, weighs 8.5 lb (3.856 kg), and depicts a knight rendered in Art Deco style holding a crusader's sword standing on a reel of film with five spokes. The five spokes represent the original branches of the Academy: Actors, Writers, Directors, Producers, and Technicians.[16]
|
38 |
+
|
39 |
+
Sculptor George Stanley (who also did the Muse Fountain at the Hollywood Bowl) sculpted Cedric Gibbons' design. The statuettes presented at the initial ceremonies were gold-plated solid bronze. Within a few years, the bronze was abandoned in favor of Britannia metal, a pewter-like alloy which is then plated in copper, nickel silver, and finally, 24-karat gold.[8] Due to a metal shortage during World War II, Oscars were made of painted plaster for three years. Following the war, the Academy invited recipients to redeem the plaster figures for gold-plated metal ones.[17] The only addition to the Oscar since it was created is a minor streamlining of the base. The original Oscar mold was cast in 1928 at the C.W. Shumway & Sons Foundry in Batavia, Illinois, which also contributed to casting the molds for the Vince Lombardi Trophy and Emmy Award's statuettes. From 1983 to 2015,[18] approximately 50 Oscars in a tin alloy with gold plating were made each year in Chicago by Illinois manufacturer R.S. Owens & Company.[19] It would take between three and four weeks to manufacture 50 statuettes.[20] In 2016, the Academy returned to bronze as the core metal of the statuettes, handing manufacturing duties to Walden, New York-based Polich Tallix Fine Art Foundry.[21][22] While based on a digital scan of an original 1929 Oscar, the statuettes retain their modern-era dimensions and black pedestal. Cast in liquid bronze from 3D-printed ceramic molds and polished, they are then electroplated in 24-karat gold by Brooklyn, New York–based Epner Technology. The time required to produce 50 such statuettes is roughly three months.[23] R.S. Owens is expected to continue producing other awards for the Academy and service existing Oscars that need replating.[24]
|
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|
41 |
+
The Academy officially adopted the name "Oscar" for the trophies in 1939. However, the origin of the nickname is disputed.[25]
|
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+
|
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+
One biography of Bette Davis, who was a president of the Academy in 1941, claims she named the award after her first husband, band leader Harmon Oscar Nelson. A frequently mentioned originator is Margaret Herrick, the Academy executive secretary, who, when she first saw the award in 1931, said the statuette reminded her of "Uncle Oscar", a nickname for her cousin Oscar Pierce.[26]
|
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+
|
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Columnist Sidney Skolsky, who was present during Herrick's naming in 1931, wrote that "Employees have affectionately dubbed their famous statuette 'Oscar.'"[27] The Academy credits Skolsky with "the first confirmed newspaper reference" to Oscar in his column on March 16, 1934, which was written about that year's 6th Academy Awards.[28] The 1934 awards appeared again in another early media mention of Oscar: a Time magazine story.[29] In the ceremonies that year, Walt Disney was the first to thank the Academy for his "Oscar" during his acceptance speech.[30]
|
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|
47 |
+
To prevent information identifying the Oscar winners from leaking ahead of the ceremony, Oscar statuettes presented at the ceremony have blank baseplates. Until 2010, winners returned their statuettes to the Academy and had to wait several weeks to have their names inscribed on their respective Oscars. Since 2010, winners have had the option of having engraved nameplates applied to their statuettes at an inscription-processing station at the Governor's Ball, a party held immediately after the Oscar ceremony. The R.S. Owens company has engraved nameplates made before the ceremony, bearing the name of every potential winner. The nameplates for the non-winning nominees are later recycled.[31][32]
|
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|
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Prior to 1950 Oscar statuettes were (and remain) the property of the recipient.[33] Since then the statuettes have been legally encumbered by the requirement that the statuette be first offered for sale back to the Academy for US$1. If a winner refuses to agree to this stipulation, then the Academy keeps the statuette. Academy Awards predating this agreement have been sold in public auctions and private deals for six-figure sums.[34]
|
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|
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In 1989 Michael Todd's grandson tried to sell Todd's Best Picture Oscar for his 1956 production of Around the World in 80 Days to a movie prop collector. The Academy earned enforcement of its statuette contract by gaining a permanent injunction against the sale.
|
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+
|
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In 1992, Harold Russell consigned his 1946 Oscar for Best Supporting Actor for The Best Years of Our Lives to auction in order to raise money for his wife's medical expenses. Though his decision caused controversy, the first-ever Oscar to be sold passed to a private collector on August 6, 1992 for $60,500 ($110,200 today). Russell defended his action, saying, "I don't know why anybody would be critical. My wife's health is much more important than sentimental reasons. The movie will be here, even if Oscar isn't." .[35]
|
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In December 2011, Orson Welles' 1941 Oscar for Citizen Kane (Academy Award for Best Original Screenplay) was put up for auction, after his heirs won a 2004 court decision contending that Welles did not sign any agreement to return the statue to the Academy.[36] On December 20, 2011, it sold in an online auction for US$861,542 ($0.98 million today).[37]
|
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|
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+
Some buyers have subsequently returned the statuettes to the Academy, which keeps them in its treasury.[34]
|
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In addition to the Academy Award of Merit (Oscar award), there are nine honorary (non-competitive) awards presented by the Academy from time to time (except for the Academy Honorary Award, the Technical Achievement Award, and the Student Academy Awards, which are presented annually):[38]
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|
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The Academy also awards Nicholl Fellowships in Screenwriting.
|
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|
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+
Since 2004, Academy Award nomination results have been announced to the public in mid-January. Prior to that, the results were announced in early February.
|
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|
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+
The Academy of Motion Picture Arts and Sciences (AMPAS), a professional honorary organization, maintains a voting membership of over 7,000 as of 2018[update].[39]
|
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+
Academy membership is divided into different branches, with each representing a different discipline in film production. Actors constitute the largest voting bloc, numbering 1,311 members (22 percent) of the Academy's composition. Votes have been certified by the auditing firm PricewaterhouseCoopers (and its predecessor Price Waterhouse) since the 7th Academy Awards in 1935. The firm mails the ballots of eligible nominees to members of the Academy in December to reflect the previous eligible year with a due date sometime in January of the next year, then tabulates the votes in a process that takes thousands of hours.[40][41][42]
|
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|
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All AMPAS members must be invited to join by the Board of Governors, on behalf of Academy Branch Executive Committees. Membership eligibility may be achieved by a competitive nomination or a member may submit a name based on other significant contributions to the field of motion pictures.
|
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New membership proposals are considered annually. The Academy does not publicly disclose its membership, although as recently as 2007 press releases have announced the names of those who have been invited to join. The 2007 release also stated that it has just under 6,000 voting members. While the membership had been growing, stricter policies have kept its size steady since then.[43]
|
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In 2012, the results of a study conducted by the Los Angeles Times were published describing the demographic breakdown of approximately 88% of AMPAS' voting membership. Of the 5,100+ active voters confirmed, 94% were Caucasian, 77% were male, and 54% were found to be over the age of 60. 33% of voting members are former nominees (14%) and winners (19%).[44]
|
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In May 2011, the Academy sent a letter advising its 6,000 or so voting members that an online system for Oscar voting would be implemented in 2013.[45]
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According to Rules 2 and 3 of the official Academy Awards Rules, a film must open in the previous calendar year, from midnight at the start of January 1 to midnight at the end of December 31, in Los Angeles County, California, and play for seven consecutive days, to qualify (except for the Best International Feature Film, Best Documentary Feature, and awards in short film categories). Additionally, the film must be shown at least three times on each day of its qualifying run, with at least one of the daily showings starting between 6 pm and 10 pm local time.[46][47]
|
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The Best International Feature Film award does not require a U.S. release. It requires the film to be submitted as its country's official selection.
|
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The Best Documentary Feature award requires either week-long releases in both Los Angeles County and New York City[a] during the previous calendar year, or a qualifying award at a competitive film festival from the Documentary Feature Qualifying Festival list (regardless of any public exhibition or distribution), or a submission in the International Feature Film category as its country's official selection. The qualifying theatrical runs must meet the same requirements as those for non-documentary films regarding numbers and times of screenings. Additionally, a film must have been reviewed by a critic from The New York Times, Time Out New York, the Los Angeles Times, or LA Weekly.[49]
|
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Awards in short film categories (Best Documentary Short Subject, Best Animated Short Film, and Best Live Action Short Film) have noticeably different eligibility rules from most other competitive awards. First, the qualifying period for release does not coincide with a calendar year, instead covering a one-year period starting on October 1 and ending on September 30 of the calendar year before the ceremony. Second, there are multiple methods of qualification. The main method is a week-long theatrical release in either Los Angeles County or New York City during the eligibility period. Films also can qualify by winning specified awards at one of a number of competitive film festivals designated by the Academy, also without regard to prior public distribution. Finally, a film that is selected as a gold, silver, or bronze medal winner in an appropriate category of the immediately previous Student Academy Awards is also eligible (Documentary category for that award, and Animation, Narrative, Alternative, or International for the other awards). The requirements for the qualifying theatrical run are also different from those for other awards. Only one screening per day is required. For the Documentary award, the screening must start between noon and 10 pm local time; for other awards, no specific start time is required, but the film must appear in regular theater listings with dates and screening times.[49][50]
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For example, the 2009 Best Picture winner, The Hurt Locker, was actually first released in 2008, but did not qualify for the 2008 awards, as it did not play its Oscar-qualifying run in Los Angeles until mid-2009, thus qualifying for the 2009 awards. Foreign films must include English subtitles, and each country can submit only one film for consideration in the International Feature Film category per year.[51]
|
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Rule 2 states that a film must be feature-length, defined as a minimum of 40 minutes, except for short-subject awards, and it must exist either on a 35 mm or 70 mm film print or in 24 frame/s or 48 frame/s progressive scan digital cinema format with a minimum projector resolution of 2048 by 1080 pixels.[52] Since the 90th Academy Awards, presented in 2018, multi-part and limited series have been ineligible for the Best Documentary Feature award. This followed the win of O.J.: Made in America, an eight-hour presentation that was screened in a limited release before being broadcast in five parts on ABC and ESPN, in that category in 2017. The Academy's announcement of the new rule made no direct mention of that film.[53]
|
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Producers must submit an Official Screen Credits online form before the deadline; in case it is not submitted by the defined deadline, the film will be ineligible for Academy Awards in any year. The form includes the production credits for all related categories. Then, each form is checked and put in a Reminder List of Eligible Releases.
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In late December, ballots, and copies of the Reminder List of Eligible Releases are mailed to around 6,000 active members. For most categories, members from each of the branches vote to determine the nominees only in their respective categories (i.e. only directors vote for directors, writers for writers, actors for actors, etc.). In the special case of Best Picture, all voting members are eligible to select the nominees. In all major categories, a variant of the single transferable vote is used, with each member casting a ballot with up to five nominees (ten for Best Picture) ranked preferentially.[54][55][56] In certain categories, including International Feature Film, Documentary and Animated Feature, nominees are selected by special screening committees made up of members from all branches.
|
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|
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In most categories, the winner is selected from among the nominees by plurality voting of all members.[54][56] Since 2009, the Best Picture winner has been chosen by instant runoff voting.[56][57] Since 2013, re-weighted range voting has been used to select the nominees for the Best Visual Effects.[58][59]
|
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Film companies will spend as much as several million dollars on marketing to awards voters for a movie in the running for Best Picture, in attempts to improve chances of receiving Oscars and other movie awards conferred in Oscar season. The Academy enforces rules to limit overt campaigning by its members so as to try to eliminate excesses and prevent the process from becoming undignified. It has an awards czar on staff who advises members on allowed practices and levies penalties on offenders.[60] For example, a producer of the 2009 Best Picture nominee The Hurt Locker was disqualified as a producer in the category when he contacted associates urging them to vote for his film and not another that was seen as the front-runner (The Hurt Locker eventually won).
|
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The major awards are presented at a live televised ceremony, commonly in late February or early March following the relevant calendar year, and six weeks after the announcement of the nominees. It is the culmination of the film awards season, which usually begins during November or December of the previous year. This is an elaborate extravaganza, with the invited guests walking up the red carpet in the creations of the most prominent fashion designers of the day. Black tie dress is the most common outfit for men, although fashion may dictate not wearing a bow-tie, and musical performers sometimes do not adhere to this. (The artists who recorded the nominees for Best Original Song quite often perform those songs live at the awards ceremony, and the fact that they are performing is often used to promote the television broadcast.)
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The Academy Awards is the world's longest-running awards show televised live from the U.S. to all time zones in North America and worldwide, and gathers billions of viewers elsewhere throughout the world.[61] The Oscars were first televised in 1953 by NBC, which continued to broadcast the event until 1960, when ABC took over, televising the festivities (including the first color broadcast of the event in 1966) through 1970. NBC regained the rights for five years (1971–75), then ABC resumed broadcast duties in 1976 and its current contract with the Academy runs through 2028.[62] The Academy has also produced condensed versions of the ceremony for broadcast in international markets (especially those outside of the Americas) in more desirable local timeslots. The ceremony was broadcast live internationally for the first time via satellite since 1970, but only two South American countries, Chile and Brazil, purchased the rights to air the broadcast. By that time, the television rights to the Academy Awards had been sold in 50 countries. A decade later, the rights were already being sold to 60 countries, and by 1984, the TV rights to the Awards were licensed in 76 countries.
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The ceremonies were moved up from late March/early April to late February, since 2004, to help disrupt and shorten the intense lobbying and ad campaigns associated with Oscar season in the film industry. Another reason was because of the growing TV ratings success coinciding with the NCAA Basketball Tournament, which would cut into the Academy Awards audience. (In 1976 and 1977, ABC's regained Oscars were moved from Tuesday to Monday and went directly opposite NBC's NCAA title game.) The earlier date is also to the advantage of ABC, as it now usually occurs during the highly profitable and important February sweeps period. Some years, the ceremony is moved into the first Sunday of March to avoid a clash with the Winter Olympic Games. Another reason for the move to late February and early March is also to avoid the awards ceremony occurring so close to the religious holidays of Passover and Easter, which for decades had been a grievance from members and the general public.[63] Advertising is somewhat restricted, however, as traditionally no movie studios or competitors of official Academy Award sponsors may advertise during the telecast. The production of the Academy Awards telecast currently holds the distinction of winning the most Emmys in history, with 47 wins and 195 nominations overall since that award's own launch in 1949.[64]
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After many years of being held on Mondays at 9:00 pm Eastern/6:00 p.m Pacific, since the 1999 ceremonies, it was moved to Sundays at 8:30 pm ET/5:30 pm PT.[65] The reasons given for the move were that more viewers would tune in on Sundays, that Los Angeles rush-hour traffic jams could be avoided, and an earlier start time would allow viewers on the East Coast to go to bed earlier.[66] For many years the film industry opposed a Sunday broadcast because it would cut into the weekend box office.[67] In 2010, the Academy contemplated moving the ceremony even further back into January, citing TV viewers' fatigue with the film industry's long awards season. However, such an accelerated schedule would dramatically decrease the voting period for its members, to the point where some voters would only have time to view the contending films streamed on their computers (as opposed to traditionally receiving the films and ballots in the mail). Furthermore, a January ceremony on Sunday would clash with National Football League playoff games.[68] In 2018, the Academy announced that the ceremony would be moved from late February to mid February beginning with the 92nd Academy Awards in 2020.[69]
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Originally scheduled for April 8, 1968, the 40th Academy Awards ceremony was postponed for two days, because of the assassination of Dr. Martin Luther King, Jr.. On March 30, 1981, the 53rd Academy Awards was postponed for one day, after the shooting of President Ronald Reagan and others in Washington, D.C.[70]
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In 1993, an In Memoriam segment was introduced,[71] honoring those who had made a significant contribution to cinema who had died in the preceding 12 months, a selection compiled by a small committee of Academy members.[72] This segment has drawn criticism over the years for the omission of some names. Criticism was also levied for many years regarding another aspect, with the segment having a "popularity contest" feel as the audience varied their applause to those who had died by the subject's cultural impact; the applause has since been muted during the telecast, and the audience is discouraged from clapping during the segment and giving silent reflection instead. This segment was later followed by a commercial break.
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In terms of broadcast length, the ceremony generally averages three and a half hours. The first Oscars, in 1929, lasted 15 minutes. At the other end of the spectrum, the 2002 ceremony lasted four hours and twenty-three minutes.[73][74] In 2010, the organizers of the Academy Awards announced winners' acceptance speeches must not run past 45 seconds. This, according to organizer Bill Mechanic, was to ensure the elimination of what he termed "the single most hated thing on the show" – overly long and embarrassing displays of emotion.[75] In 2016, in a further effort to streamline speeches, winners' dedications were displayed on an on-screen ticker.[76] During the 2018 ceremony, host Jimmy Kimmel acknowledged how long the ceremony had become, by announcing that he would give a brand-new jet ski to whoever gave the shortest speech of the night (a reward won by Mark Bridges when accepting his Best Costume Design award for Phantom Thread).[77] The Wall Street Journal analyzed the average minutes spent across the 2014–2018 telecasts as follows: 14 on song performances; 25 on the hosts' speeches; 38 on prerecorded clips; and 78 on the awards themselves, broken into 24 on the introduction and announcement, 24 on winners walking to the stage, and 30 on their acceptance speeches.[78]
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Although still dominant in ratings, the viewership of the Academy Awards has steadily dropped; the 88th Academy Awards were the lowest-rated in the past eight years (although with increases in male and 18–49 viewership), while the show itself also faced mixed reception. Following the show, Variety reported that ABC was, in negotiating an extension to its contract to broadcast the Oscars, seeking to have more creative control over the broadcast itself. Currently and nominally, AMPAS is responsible for most aspects of the telecast, including the choice of production staff and hosting, although ABC is allowed to have some input on their decisions.[79] In August 2016, AMPAS extended its contract with ABC through 2028: the contract neither contains any notable changes nor gives ABC any further creative control over the telecast.[80]
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Historically, the telecast's viewership is higher when box-office hits are favored to win the Best Picture award. More than 57.25 million viewers tuned to the telecast for the 70th Academy Awards in 1998, the year of Titanic, which generated a box office haul during its initial 1997–98 run of US$600.8 million in the US, a box office record that would remain unsurpassed for years.[81] The 76th Academy Awards ceremony, in which The Lord of the Rings: The Return of the King (pre-telecast box office earnings of US$368 million) received 11 Awards including Best Picture, drew 43.56 million viewers.[82] The most watched ceremony based on Nielsen ratings to date, however, was the 42nd Academy Awards (Best Picture Midnight Cowboy) which drew a 43.4% household rating on April 7, 1970.[83]
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By contrast, ceremonies honoring films that have not performed well at the box office tend to show weaker ratings, despite how much critically acclaimed those films have been. The 78th Academy Awards which awarded low-budget independent film Crash (with a pre-Oscar gross of US$53.4 million) generated an audience of 38.64 million with a household rating of 22.91%.[84] In 2008, the 80th Academy Awards telecast was watched by 31.76 million viewers on average with an 18.66% household rating, the lowest-rated and least-watched ceremony at the time, in spite of celebrating 80 years of the Academy Awards.[85] The Best Picture winner of that particular ceremony was another independent film (No Country for Old Men).
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The 92nd Academy Awards drew an average of 23.6 million viewers -- a new low.
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The Academy Film Archive holds copies of every Academy Awards ceremony since the 1949 Oscars and material on many prior ceremonies, along with ancillary material related to more recent shows. Copies are held on a variety of film, video, and digital formats.[86]
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In 1929, the first Academy Awards were presented at a banquet dinner at the Hollywood Roosevelt Hotel. From 1930 to 1943, the ceremony alternated between two venues: the Ambassador Hotel on Wilshire Boulevard and the Biltmore Hotel in downtown Los Angeles.
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Grauman's Chinese Theatre in Hollywood then hosted the awards from 1944 to 1946, followed by the Shrine Auditorium in Los Angeles from 1947 to 1948. The 21st Academy Awards in 1949 were held at the Academy Award Theatre at what had been the Academy's headquarters on Melrose Avenue in Hollywood.[87]
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From 1950 to 1960, the awards were presented at Hollywood's Pantages Theatre. With the advent of television, the awards from 1953 to 1957 took place simultaneously in Hollywood and New York, first at the NBC International Theatre (1953) and then at the NBC Century Theatre, after which the ceremony took place solely in Los Angeles. The Oscars moved to the Santa Monica Civic Auditorium in Santa Monica, California, in 1961. By 1969, the Academy decided to move the ceremonies back to Downtown Los Angeles, this time to the Dorothy Chandler Pavilion at the Los Angeles County Music Center. In the late 1990s and early 2000s, the ceremony returned to the Shrine.
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In 2002, Hollywood's Dolby Theatre (previously known as the Kodak Theatre) became the presentation's current venue.[88]
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In the first year of the awards, the Best Directing award was split into two categories (Drama and Comedy). At times, the Best Original Score award has also been split into separate categories (Drama and Comedy/Musical). From the 1930s through the 1960s, the Art Direction (now Production Design), Cinematography, and Costume Design awards were likewise split into two categories (black-and-white films and color films). Prior to 2012, the Production Design award was called Art Direction, while the Makeup and Hairstyling award was called Makeup.
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In August 2018, the Academy announced that several categories would not be televised live, but rather be recorded during commercial breaks and aired later in the ceremony.[91]
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Following dissent from Academy members, they announced that they would indeed air all 24 categories live. This followed a number of proposals (among them, the introduction of a Popular Film category) that the Academy had announced but did not implement.[92]
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The Board of Governors meets each year and considers new award categories. To date, the following categories have been proposed:
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The Special Academy Awards are voted on by special committees, rather than by the Academy membership as a whole. They are not always presented on an annual basis.
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Due to the positive exposure and prestige of the Academy Awards, many studios spend millions of dollars and hire publicists specifically to promote their films during what is typically called the "Oscar season". This has generated accusations of the Academy Awards being influenced more by marketing than by quality. William Friedkin, an Academy Award-winning film director and former producer of the ceremony, expressed this sentiment at a conference in New York in 2009, describing it as "the greatest promotion scheme that any industry ever devised for itself".[99]
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Tim Dirks, editor of AMC's filmsite.org, has written of the Academy Awards:
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Unfortunately, the critical worth, artistic vision, cultural influence and innovative qualities of many films are not given the same voting weight. Especially since the 1980s, moneymaking "formula-made" blockbusters with glossy production values have often been crowd-pleasing titans (and Best Picture winners), but they haven't necessarily been great films with depth or critical acclaim by any measure.[100]
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A recent technique that has been claimed to be used during the Oscar season is the whisper campaign. These campaigns are intended to spread negative perceptions of other movies nominated and are believed to be perpetrated by those that were involved in creating the movie. Examples of whisper campaigns include the allegations against Zero Dark Thirty suggesting that it justifies torture and the claim that Lincoln distorts history.[101]
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Typical criticism of the Academy Awards for Best Picture is that among the winners and nominees there is an over-representation of romantic historical epics, biographical dramas, romantic dramedies and family melodramas, most of which are released in the U.S. the last three months of the calendar year. The Oscars have been infamously known for selecting specific genres of movies to be awarded. This has led to the coining of the term 'Oscar bait', describing such movies. This has led at times to more specific criticisms that the Academy is disconnected from the audience, e.g., by favoring 'Oscar bait' over audience favorites, or favoring historical melodramas over critically acclaimed movies that depict current life issues.[102]
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The Academy Awards have long received criticism over its lack of diversity among the nominees.[103][104][105] This criticism is based on the statistics from every Academy Awards since 1929, which shows us that only 6.4% of academy award nominees have been non-white and since 1991, 11.2% of nominees have been non-white, with the rate of winners being even more polarizing.[106] More white actresses have won Oscars for yellowface portrayals of Asian characters than actual Asian actresses.[107][108] The 88th awards ceremony became the target of a boycott, popularized on social media with the hashtag #OscarsSoWhite, based on critics' perception that its all-white acting nominee list reflected bias. In response, the Academy initiated "historic" changes in membership by the year 2020.[109][110]
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Acting prizes in certain years have been criticized for not recognizing superior performances so much as being awarded for personal popularity,[111] to make up for a "snub" for a performance/work that proved in time to be more popular and/or renowned than the one actually awarded, or presented as a "career honor" to recognize a distinguished nominee's entire body of work.[27]
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Following the 91st Academy Awards in February 2019 in which the Netflix-broadcast film Roma had been nominated for ten awards including the Best Picture category, Steven Spielberg and other members of the Academy discussed changing the requirements through the Board of Governors for films as to exclude those from Netflix and other media streaming services. Spielberg had been concerned that Netflix as a movie production and distribution studio could spend much more than typical Oscar-winning films and have much wider and earlier distribution than other Best Picture-nominated films, while still being able to meet the minimal theatrical-run status to qualify for an Oscar.[112] The United States Department of Justice, having heard of this potential rule change, wrote a letter to the Academy in March 2019, cautioning them that placing additional restrictions on films that originate from streaming media services without proper justification could raise anti-trust concerns against the Academy.[113] Following its April 2019 board meeting, the Academy Board of Governors agreed to retain the current rules that allow for streaming media films to be eligible for Oscars as long as they enjoy limited theatrical runs.[114]
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Some winners critical of the Academy Awards have boycotted the ceremonies and refused to accept their Oscars. The first to do so was screenwriter Dudley Nichols (Best Writing in 1935 for The Informer). Nichols boycotted the 8th Academy Awards ceremony because of conflicts between the Academy and the Writers' Guild.[115] Nichols eventually accepted the 1935 award three years later, at the 1938 ceremony. Nichols was nominated for three further Academy Awards during his career.
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George C. Scott became the second person to refuse his award (Best Actor in 1970 for Patton) at the 43rd Academy Awards ceremony. Scott described it as a "meat parade", saying, "I don't want any part of it."[116][117][118]
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The third person to refuse the award was Marlon Brando, who refused his award (Best Actor for 1972's The Godfather), citing the film industry's discrimination and mistreatment of Native Americans. At the 45th Academy Awards ceremony, Brando sent actress and civil rights activist Sacheen Littlefeather to read a 15-page speech, detailing his criticisms, which was booed by the audience.[119][115]
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At the 89th Academy Awards ceremony, Warren Beatty and Faye Dunaway mistakenly announced La La Land as the recipient of the Best Picture award, instead of Moonlight, the actual winner. Beatty had been given the wrong envelope and after hesitating during the announcement, handed the envelope to Dunaway, which listed Emma Stone as Best Actress for La La Land and led to the confusion. The proper winner was announced after the acceptance speeches by La La Land producers Fred Berger, Jordan Horowitz and Marc Platt.
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The following year, Beatty and Dunaway were invited back as presenters of the Best Picture award, which they accomplished without error.[120]
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Nine films have been disqualified before an official award ceremony because they violated the regulations:[123]
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One film was disqualified after winning the award, and had the winner return the Oscar:
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The following events are closely associated with the annual Academy Awards:
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It has become a tradition to give out gift bags to the presenters and performers at the Oscars. In recent years, these gifts have also been extended to award nominees and winners.[124] The value of each of these gift bags can reach into the tens of thousands of dollars. In 2014, the value was reported to be as high as US$80,000.[125] The value has risen to the point where the U.S. Internal Revenue Service issued a statement regarding the gifts and their taxable status.[126]
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Oscar gift bags have included vacation packages to Hawaii and Mexico and Japan, a private dinner party for the recipient and friends at a restaurant, videophones, a four-night stay at a hotel, watches, bracelets, spa treatments, bottles of vodka, maple salad dressing, weight-loss gummie candy and up to $25,000 worth of cosmetic treatments and rejuvenation procedures such as lip fillers and chemical peels from New York City facial plastic surgeon Konstantin Vasyukevich.[124][127][128][129][130] Some of the gifts have even had a "risque" element to them; in 2014, the adult products retailer Adam & Eve had a "Secret Room Gifting Suite". Celebrities visiting the gifting suite included Judith Hoag, Carolyn Hennesy, Kate Linder, Chris Mulkey, Jim O'Heir, and John Salley.[131]
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From 2006 onwards, results are Live+SD; all previous years are live viewing.[132]
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The term "Oscar" is a registered trademark of the AMPAS; however, in the Italian language, it is used generically to refer to any award or award ceremony, regardless of which field.[140][141]
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Green Book
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Parasite
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The Academy Awards, popularly known as the Oscars,[1] are awards for artistic and technical merit in the film industry. Given annually by the Academy of Motion Picture Arts and Sciences (AMPAS), the awards are an international recognition of excellence in cinematic achievements, as assessed by the Academy's voting membership. The various category winners are awarded a copy of a golden statuette, officially called the "Academy Award of Merit", although more commonly referred to by its nickname, the "Oscar". The statuette depicts a knight rendered in the Art Deco style.
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The award was originally sculpted by George Stanley from a design sketch by Cedric Gibbons.[2] AMPAS first presented it in 1929 at a private dinner hosted by Douglas Fairbanks in the Hollywood Roosevelt Hotel in what would become known as the 1st Academy Awards.[3][4] The Academy Awards ceremony was first broadcast by radio in 1930 and was televised for the first time in 1953. It is the oldest worldwide entertainment awards ceremony and is now televised live worldwide.[5] It is also the oldest of the four major annual American entertainment awards; its equivalents – the Emmy Awards for television, the Tony Awards for theater, and the Grammy Awards for music – are modeled after the Academy Awards.[6] They are widely cited as the most famous and prestigious competitive awards in the field of entertainment.
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The 92nd Academy Awards ceremony, honoring the best films of 2019, was held on February 9, 2020. As with the previous ceremony, there was no host. The ceremony was broadcast on ABC. It took place at the Dolby Theatre in Los Angeles, California, for the 18th consecutive year. This ceremony marked the first time in Oscar history that more than three films had received ten or more nominations.[7] A total of 3,140 Oscar statuettes have been awarded since its inception in 1929.[8]
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The first Academy Awards presentation was held on May 16, 1929, at a private dinner function at the Hollywood Roosevelt Hotel with an audience of about 270 people.[9]
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The post-awards party was held at the Mayfair Hotel.[5] The cost of guest tickets for that night's ceremony was $5 ($74 in 2020 dollars). Fifteen statuettes were awarded, honoring artists, directors and other participants in the film-making industry of the time, for their works during the 1927–28 period. The ceremony ran for 15 minutes.
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Winners were announced to media three months earlier. That was changed for the second ceremony in 1930. Since then, for the rest of the first decade, the results were given to newspapers for publication at 11:00 pm on the night of the awards.[5] This method was used until 1940, when the Los Angeles Times announced the winners before the ceremony began; as a result, the Academy has, since 1941, used a sealed envelope to reveal the names of the winners.[5]
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The first Best Actor awarded was Emil Jannings, for his performances in The Last Command and The Way of All Flesh. He had to return to Europe before the ceremony, so the Academy agreed to give him the prize earlier; this made him the first Academy Award winner in history. At that time, winners were recognized for the entirety of their work done in a certain category during the qualifying period; for example, Jannings received the award for two movies in which he starred during that period, and Janet Gaynor later won a single Oscar for performances in three films. With the fourth ceremony, however, the system changed, and professionals were honored for a specific performance in a single film. For the first six ceremonies, the eligibility period spanned two calendar years.[5]
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At the 29th ceremony, held in 1957, the Best Foreign Language Film category, now known as Best International Feature Film, was introduced. Until then, foreign-language films had been honored with the Special Achievement Award.
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Perhaps the most widely seen streaker in history was 34-year-old Robert Opel, who streaked across the stage of The Dorothy Chandler Pavilion in Los Angeles flashing a peace sign on national US television at the 46th Academy Awards in 1974. Bemused host David Niven quipped, "Isn't it fascinating to think that probably the only laugh that man will ever get in his life is by stripping off and showing his shortcomings?" Later, evidence arose suggesting that Opel's appearance was facilitated as a publicity stunt by the show's producer Jack Haley Jr. Robert Metzler, the show's business manager, believed that the incident had been planned in some way; during the dress rehearsal Niven had asked Metzler's wife to borrow a pen so he could write down the famous line, which was thus not the ad-lib it appeared to be.[10]
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The 74th Academy Awards, held in 2002, presented the first Academy Award for Best Animated Feature.[11]
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Since 1973, all Academy Awards ceremonies have ended with the Academy Award for Best Picture.
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Traditionally, the previous year's winner for Best Actor and Best Supporting Actor present the awards for Best Actress and Best Supporting Actress, while the previous year's winner for Best Actress and Best Supporting Actress present the awards for Best Actor and Best Supporting Actor.
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Parasite became the first foreign-language film to win Best Picture at the February 9, 2020 award ceremony.[12]
|
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Tom Hanks announced at the 2020 Oscar Ceremony, the opening of the Academy Museum of Motion Pictures on December 14, 2020.[13] The museum development started in 2017 under Kerry Brougher, but is now led by Bill Kramer.[14] The industry curated exhibits will be geared toward the history of motion picture, the art & science of film making, exhibiting trailblazing directors, actors, film makers, sound editors and more, and will house famous artifacts from acclaimed movies like Dorothy's Ruby Red Slippers.
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Because of COVID-19, Academy president David Rubin and CEO Dawn Hudson announced that for the 2021 Oscar Ceremony, streaming movies not shown in theaters would be eligible, though at some point the requirement that movies be shown in theaters would return.[15]
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The best known award is the Academy Award of Merit, more popularly known as the Oscar statuette.[8] Made of gold-plated bronze on a black metal base, it is 13.5 in (34.3 cm) tall, weighs 8.5 lb (3.856 kg), and depicts a knight rendered in Art Deco style holding a crusader's sword standing on a reel of film with five spokes. The five spokes represent the original branches of the Academy: Actors, Writers, Directors, Producers, and Technicians.[16]
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Sculptor George Stanley (who also did the Muse Fountain at the Hollywood Bowl) sculpted Cedric Gibbons' design. The statuettes presented at the initial ceremonies were gold-plated solid bronze. Within a few years, the bronze was abandoned in favor of Britannia metal, a pewter-like alloy which is then plated in copper, nickel silver, and finally, 24-karat gold.[8] Due to a metal shortage during World War II, Oscars were made of painted plaster for three years. Following the war, the Academy invited recipients to redeem the plaster figures for gold-plated metal ones.[17] The only addition to the Oscar since it was created is a minor streamlining of the base. The original Oscar mold was cast in 1928 at the C.W. Shumway & Sons Foundry in Batavia, Illinois, which also contributed to casting the molds for the Vince Lombardi Trophy and Emmy Award's statuettes. From 1983 to 2015,[18] approximately 50 Oscars in a tin alloy with gold plating were made each year in Chicago by Illinois manufacturer R.S. Owens & Company.[19] It would take between three and four weeks to manufacture 50 statuettes.[20] In 2016, the Academy returned to bronze as the core metal of the statuettes, handing manufacturing duties to Walden, New York-based Polich Tallix Fine Art Foundry.[21][22] While based on a digital scan of an original 1929 Oscar, the statuettes retain their modern-era dimensions and black pedestal. Cast in liquid bronze from 3D-printed ceramic molds and polished, they are then electroplated in 24-karat gold by Brooklyn, New York–based Epner Technology. The time required to produce 50 such statuettes is roughly three months.[23] R.S. Owens is expected to continue producing other awards for the Academy and service existing Oscars that need replating.[24]
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The Academy officially adopted the name "Oscar" for the trophies in 1939. However, the origin of the nickname is disputed.[25]
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One biography of Bette Davis, who was a president of the Academy in 1941, claims she named the award after her first husband, band leader Harmon Oscar Nelson. A frequently mentioned originator is Margaret Herrick, the Academy executive secretary, who, when she first saw the award in 1931, said the statuette reminded her of "Uncle Oscar", a nickname for her cousin Oscar Pierce.[26]
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Columnist Sidney Skolsky, who was present during Herrick's naming in 1931, wrote that "Employees have affectionately dubbed their famous statuette 'Oscar.'"[27] The Academy credits Skolsky with "the first confirmed newspaper reference" to Oscar in his column on March 16, 1934, which was written about that year's 6th Academy Awards.[28] The 1934 awards appeared again in another early media mention of Oscar: a Time magazine story.[29] In the ceremonies that year, Walt Disney was the first to thank the Academy for his "Oscar" during his acceptance speech.[30]
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To prevent information identifying the Oscar winners from leaking ahead of the ceremony, Oscar statuettes presented at the ceremony have blank baseplates. Until 2010, winners returned their statuettes to the Academy and had to wait several weeks to have their names inscribed on their respective Oscars. Since 2010, winners have had the option of having engraved nameplates applied to their statuettes at an inscription-processing station at the Governor's Ball, a party held immediately after the Oscar ceremony. The R.S. Owens company has engraved nameplates made before the ceremony, bearing the name of every potential winner. The nameplates for the non-winning nominees are later recycled.[31][32]
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Prior to 1950 Oscar statuettes were (and remain) the property of the recipient.[33] Since then the statuettes have been legally encumbered by the requirement that the statuette be first offered for sale back to the Academy for US$1. If a winner refuses to agree to this stipulation, then the Academy keeps the statuette. Academy Awards predating this agreement have been sold in public auctions and private deals for six-figure sums.[34]
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In 1989 Michael Todd's grandson tried to sell Todd's Best Picture Oscar for his 1956 production of Around the World in 80 Days to a movie prop collector. The Academy earned enforcement of its statuette contract by gaining a permanent injunction against the sale.
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In 1992, Harold Russell consigned his 1946 Oscar for Best Supporting Actor for The Best Years of Our Lives to auction in order to raise money for his wife's medical expenses. Though his decision caused controversy, the first-ever Oscar to be sold passed to a private collector on August 6, 1992 for $60,500 ($110,200 today). Russell defended his action, saying, "I don't know why anybody would be critical. My wife's health is much more important than sentimental reasons. The movie will be here, even if Oscar isn't." .[35]
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In December 2011, Orson Welles' 1941 Oscar for Citizen Kane (Academy Award for Best Original Screenplay) was put up for auction, after his heirs won a 2004 court decision contending that Welles did not sign any agreement to return the statue to the Academy.[36] On December 20, 2011, it sold in an online auction for US$861,542 ($0.98 million today).[37]
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Some buyers have subsequently returned the statuettes to the Academy, which keeps them in its treasury.[34]
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In addition to the Academy Award of Merit (Oscar award), there are nine honorary (non-competitive) awards presented by the Academy from time to time (except for the Academy Honorary Award, the Technical Achievement Award, and the Student Academy Awards, which are presented annually):[38]
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The Academy also awards Nicholl Fellowships in Screenwriting.
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Since 2004, Academy Award nomination results have been announced to the public in mid-January. Prior to that, the results were announced in early February.
|
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The Academy of Motion Picture Arts and Sciences (AMPAS), a professional honorary organization, maintains a voting membership of over 7,000 as of 2018[update].[39]
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Academy membership is divided into different branches, with each representing a different discipline in film production. Actors constitute the largest voting bloc, numbering 1,311 members (22 percent) of the Academy's composition. Votes have been certified by the auditing firm PricewaterhouseCoopers (and its predecessor Price Waterhouse) since the 7th Academy Awards in 1935. The firm mails the ballots of eligible nominees to members of the Academy in December to reflect the previous eligible year with a due date sometime in January of the next year, then tabulates the votes in a process that takes thousands of hours.[40][41][42]
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All AMPAS members must be invited to join by the Board of Governors, on behalf of Academy Branch Executive Committees. Membership eligibility may be achieved by a competitive nomination or a member may submit a name based on other significant contributions to the field of motion pictures.
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New membership proposals are considered annually. The Academy does not publicly disclose its membership, although as recently as 2007 press releases have announced the names of those who have been invited to join. The 2007 release also stated that it has just under 6,000 voting members. While the membership had been growing, stricter policies have kept its size steady since then.[43]
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In 2012, the results of a study conducted by the Los Angeles Times were published describing the demographic breakdown of approximately 88% of AMPAS' voting membership. Of the 5,100+ active voters confirmed, 94% were Caucasian, 77% were male, and 54% were found to be over the age of 60. 33% of voting members are former nominees (14%) and winners (19%).[44]
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In May 2011, the Academy sent a letter advising its 6,000 or so voting members that an online system for Oscar voting would be implemented in 2013.[45]
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According to Rules 2 and 3 of the official Academy Awards Rules, a film must open in the previous calendar year, from midnight at the start of January 1 to midnight at the end of December 31, in Los Angeles County, California, and play for seven consecutive days, to qualify (except for the Best International Feature Film, Best Documentary Feature, and awards in short film categories). Additionally, the film must be shown at least three times on each day of its qualifying run, with at least one of the daily showings starting between 6 pm and 10 pm local time.[46][47]
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The Best International Feature Film award does not require a U.S. release. It requires the film to be submitted as its country's official selection.
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The Best Documentary Feature award requires either week-long releases in both Los Angeles County and New York City[a] during the previous calendar year, or a qualifying award at a competitive film festival from the Documentary Feature Qualifying Festival list (regardless of any public exhibition or distribution), or a submission in the International Feature Film category as its country's official selection. The qualifying theatrical runs must meet the same requirements as those for non-documentary films regarding numbers and times of screenings. Additionally, a film must have been reviewed by a critic from The New York Times, Time Out New York, the Los Angeles Times, or LA Weekly.[49]
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Awards in short film categories (Best Documentary Short Subject, Best Animated Short Film, and Best Live Action Short Film) have noticeably different eligibility rules from most other competitive awards. First, the qualifying period for release does not coincide with a calendar year, instead covering a one-year period starting on October 1 and ending on September 30 of the calendar year before the ceremony. Second, there are multiple methods of qualification. The main method is a week-long theatrical release in either Los Angeles County or New York City during the eligibility period. Films also can qualify by winning specified awards at one of a number of competitive film festivals designated by the Academy, also without regard to prior public distribution. Finally, a film that is selected as a gold, silver, or bronze medal winner in an appropriate category of the immediately previous Student Academy Awards is also eligible (Documentary category for that award, and Animation, Narrative, Alternative, or International for the other awards). The requirements for the qualifying theatrical run are also different from those for other awards. Only one screening per day is required. For the Documentary award, the screening must start between noon and 10 pm local time; for other awards, no specific start time is required, but the film must appear in regular theater listings with dates and screening times.[49][50]
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For example, the 2009 Best Picture winner, The Hurt Locker, was actually first released in 2008, but did not qualify for the 2008 awards, as it did not play its Oscar-qualifying run in Los Angeles until mid-2009, thus qualifying for the 2009 awards. Foreign films must include English subtitles, and each country can submit only one film for consideration in the International Feature Film category per year.[51]
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Rule 2 states that a film must be feature-length, defined as a minimum of 40 minutes, except for short-subject awards, and it must exist either on a 35 mm or 70 mm film print or in 24 frame/s or 48 frame/s progressive scan digital cinema format with a minimum projector resolution of 2048 by 1080 pixels.[52] Since the 90th Academy Awards, presented in 2018, multi-part and limited series have been ineligible for the Best Documentary Feature award. This followed the win of O.J.: Made in America, an eight-hour presentation that was screened in a limited release before being broadcast in five parts on ABC and ESPN, in that category in 2017. The Academy's announcement of the new rule made no direct mention of that film.[53]
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Producers must submit an Official Screen Credits online form before the deadline; in case it is not submitted by the defined deadline, the film will be ineligible for Academy Awards in any year. The form includes the production credits for all related categories. Then, each form is checked and put in a Reminder List of Eligible Releases.
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In late December, ballots, and copies of the Reminder List of Eligible Releases are mailed to around 6,000 active members. For most categories, members from each of the branches vote to determine the nominees only in their respective categories (i.e. only directors vote for directors, writers for writers, actors for actors, etc.). In the special case of Best Picture, all voting members are eligible to select the nominees. In all major categories, a variant of the single transferable vote is used, with each member casting a ballot with up to five nominees (ten for Best Picture) ranked preferentially.[54][55][56] In certain categories, including International Feature Film, Documentary and Animated Feature, nominees are selected by special screening committees made up of members from all branches.
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In most categories, the winner is selected from among the nominees by plurality voting of all members.[54][56] Since 2009, the Best Picture winner has been chosen by instant runoff voting.[56][57] Since 2013, re-weighted range voting has been used to select the nominees for the Best Visual Effects.[58][59]
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Film companies will spend as much as several million dollars on marketing to awards voters for a movie in the running for Best Picture, in attempts to improve chances of receiving Oscars and other movie awards conferred in Oscar season. The Academy enforces rules to limit overt campaigning by its members so as to try to eliminate excesses and prevent the process from becoming undignified. It has an awards czar on staff who advises members on allowed practices and levies penalties on offenders.[60] For example, a producer of the 2009 Best Picture nominee The Hurt Locker was disqualified as a producer in the category when he contacted associates urging them to vote for his film and not another that was seen as the front-runner (The Hurt Locker eventually won).
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The major awards are presented at a live televised ceremony, commonly in late February or early March following the relevant calendar year, and six weeks after the announcement of the nominees. It is the culmination of the film awards season, which usually begins during November or December of the previous year. This is an elaborate extravaganza, with the invited guests walking up the red carpet in the creations of the most prominent fashion designers of the day. Black tie dress is the most common outfit for men, although fashion may dictate not wearing a bow-tie, and musical performers sometimes do not adhere to this. (The artists who recorded the nominees for Best Original Song quite often perform those songs live at the awards ceremony, and the fact that they are performing is often used to promote the television broadcast.)
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The Academy Awards is the world's longest-running awards show televised live from the U.S. to all time zones in North America and worldwide, and gathers billions of viewers elsewhere throughout the world.[61] The Oscars were first televised in 1953 by NBC, which continued to broadcast the event until 1960, when ABC took over, televising the festivities (including the first color broadcast of the event in 1966) through 1970. NBC regained the rights for five years (1971–75), then ABC resumed broadcast duties in 1976 and its current contract with the Academy runs through 2028.[62] The Academy has also produced condensed versions of the ceremony for broadcast in international markets (especially those outside of the Americas) in more desirable local timeslots. The ceremony was broadcast live internationally for the first time via satellite since 1970, but only two South American countries, Chile and Brazil, purchased the rights to air the broadcast. By that time, the television rights to the Academy Awards had been sold in 50 countries. A decade later, the rights were already being sold to 60 countries, and by 1984, the TV rights to the Awards were licensed in 76 countries.
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The ceremonies were moved up from late March/early April to late February, since 2004, to help disrupt and shorten the intense lobbying and ad campaigns associated with Oscar season in the film industry. Another reason was because of the growing TV ratings success coinciding with the NCAA Basketball Tournament, which would cut into the Academy Awards audience. (In 1976 and 1977, ABC's regained Oscars were moved from Tuesday to Monday and went directly opposite NBC's NCAA title game.) The earlier date is also to the advantage of ABC, as it now usually occurs during the highly profitable and important February sweeps period. Some years, the ceremony is moved into the first Sunday of March to avoid a clash with the Winter Olympic Games. Another reason for the move to late February and early March is also to avoid the awards ceremony occurring so close to the religious holidays of Passover and Easter, which for decades had been a grievance from members and the general public.[63] Advertising is somewhat restricted, however, as traditionally no movie studios or competitors of official Academy Award sponsors may advertise during the telecast. The production of the Academy Awards telecast currently holds the distinction of winning the most Emmys in history, with 47 wins and 195 nominations overall since that award's own launch in 1949.[64]
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After many years of being held on Mondays at 9:00 pm Eastern/6:00 p.m Pacific, since the 1999 ceremonies, it was moved to Sundays at 8:30 pm ET/5:30 pm PT.[65] The reasons given for the move were that more viewers would tune in on Sundays, that Los Angeles rush-hour traffic jams could be avoided, and an earlier start time would allow viewers on the East Coast to go to bed earlier.[66] For many years the film industry opposed a Sunday broadcast because it would cut into the weekend box office.[67] In 2010, the Academy contemplated moving the ceremony even further back into January, citing TV viewers' fatigue with the film industry's long awards season. However, such an accelerated schedule would dramatically decrease the voting period for its members, to the point where some voters would only have time to view the contending films streamed on their computers (as opposed to traditionally receiving the films and ballots in the mail). Furthermore, a January ceremony on Sunday would clash with National Football League playoff games.[68] In 2018, the Academy announced that the ceremony would be moved from late February to mid February beginning with the 92nd Academy Awards in 2020.[69]
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Originally scheduled for April 8, 1968, the 40th Academy Awards ceremony was postponed for two days, because of the assassination of Dr. Martin Luther King, Jr.. On March 30, 1981, the 53rd Academy Awards was postponed for one day, after the shooting of President Ronald Reagan and others in Washington, D.C.[70]
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In 1993, an In Memoriam segment was introduced,[71] honoring those who had made a significant contribution to cinema who had died in the preceding 12 months, a selection compiled by a small committee of Academy members.[72] This segment has drawn criticism over the years for the omission of some names. Criticism was also levied for many years regarding another aspect, with the segment having a "popularity contest" feel as the audience varied their applause to those who had died by the subject's cultural impact; the applause has since been muted during the telecast, and the audience is discouraged from clapping during the segment and giving silent reflection instead. This segment was later followed by a commercial break.
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In terms of broadcast length, the ceremony generally averages three and a half hours. The first Oscars, in 1929, lasted 15 minutes. At the other end of the spectrum, the 2002 ceremony lasted four hours and twenty-three minutes.[73][74] In 2010, the organizers of the Academy Awards announced winners' acceptance speeches must not run past 45 seconds. This, according to organizer Bill Mechanic, was to ensure the elimination of what he termed "the single most hated thing on the show" – overly long and embarrassing displays of emotion.[75] In 2016, in a further effort to streamline speeches, winners' dedications were displayed on an on-screen ticker.[76] During the 2018 ceremony, host Jimmy Kimmel acknowledged how long the ceremony had become, by announcing that he would give a brand-new jet ski to whoever gave the shortest speech of the night (a reward won by Mark Bridges when accepting his Best Costume Design award for Phantom Thread).[77] The Wall Street Journal analyzed the average minutes spent across the 2014–2018 telecasts as follows: 14 on song performances; 25 on the hosts' speeches; 38 on prerecorded clips; and 78 on the awards themselves, broken into 24 on the introduction and announcement, 24 on winners walking to the stage, and 30 on their acceptance speeches.[78]
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Although still dominant in ratings, the viewership of the Academy Awards has steadily dropped; the 88th Academy Awards were the lowest-rated in the past eight years (although with increases in male and 18–49 viewership), while the show itself also faced mixed reception. Following the show, Variety reported that ABC was, in negotiating an extension to its contract to broadcast the Oscars, seeking to have more creative control over the broadcast itself. Currently and nominally, AMPAS is responsible for most aspects of the telecast, including the choice of production staff and hosting, although ABC is allowed to have some input on their decisions.[79] In August 2016, AMPAS extended its contract with ABC through 2028: the contract neither contains any notable changes nor gives ABC any further creative control over the telecast.[80]
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Historically, the telecast's viewership is higher when box-office hits are favored to win the Best Picture award. More than 57.25 million viewers tuned to the telecast for the 70th Academy Awards in 1998, the year of Titanic, which generated a box office haul during its initial 1997–98 run of US$600.8 million in the US, a box office record that would remain unsurpassed for years.[81] The 76th Academy Awards ceremony, in which The Lord of the Rings: The Return of the King (pre-telecast box office earnings of US$368 million) received 11 Awards including Best Picture, drew 43.56 million viewers.[82] The most watched ceremony based on Nielsen ratings to date, however, was the 42nd Academy Awards (Best Picture Midnight Cowboy) which drew a 43.4% household rating on April 7, 1970.[83]
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|
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By contrast, ceremonies honoring films that have not performed well at the box office tend to show weaker ratings, despite how much critically acclaimed those films have been. The 78th Academy Awards which awarded low-budget independent film Crash (with a pre-Oscar gross of US$53.4 million) generated an audience of 38.64 million with a household rating of 22.91%.[84] In 2008, the 80th Academy Awards telecast was watched by 31.76 million viewers on average with an 18.66% household rating, the lowest-rated and least-watched ceremony at the time, in spite of celebrating 80 years of the Academy Awards.[85] The Best Picture winner of that particular ceremony was another independent film (No Country for Old Men).
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The 92nd Academy Awards drew an average of 23.6 million viewers -- a new low.
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The Academy Film Archive holds copies of every Academy Awards ceremony since the 1949 Oscars and material on many prior ceremonies, along with ancillary material related to more recent shows. Copies are held on a variety of film, video, and digital formats.[86]
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In 1929, the first Academy Awards were presented at a banquet dinner at the Hollywood Roosevelt Hotel. From 1930 to 1943, the ceremony alternated between two venues: the Ambassador Hotel on Wilshire Boulevard and the Biltmore Hotel in downtown Los Angeles.
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Grauman's Chinese Theatre in Hollywood then hosted the awards from 1944 to 1946, followed by the Shrine Auditorium in Los Angeles from 1947 to 1948. The 21st Academy Awards in 1949 were held at the Academy Award Theatre at what had been the Academy's headquarters on Melrose Avenue in Hollywood.[87]
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From 1950 to 1960, the awards were presented at Hollywood's Pantages Theatre. With the advent of television, the awards from 1953 to 1957 took place simultaneously in Hollywood and New York, first at the NBC International Theatre (1953) and then at the NBC Century Theatre, after which the ceremony took place solely in Los Angeles. The Oscars moved to the Santa Monica Civic Auditorium in Santa Monica, California, in 1961. By 1969, the Academy decided to move the ceremonies back to Downtown Los Angeles, this time to the Dorothy Chandler Pavilion at the Los Angeles County Music Center. In the late 1990s and early 2000s, the ceremony returned to the Shrine.
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In 2002, Hollywood's Dolby Theatre (previously known as the Kodak Theatre) became the presentation's current venue.[88]
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In the first year of the awards, the Best Directing award was split into two categories (Drama and Comedy). At times, the Best Original Score award has also been split into separate categories (Drama and Comedy/Musical). From the 1930s through the 1960s, the Art Direction (now Production Design), Cinematography, and Costume Design awards were likewise split into two categories (black-and-white films and color films). Prior to 2012, the Production Design award was called Art Direction, while the Makeup and Hairstyling award was called Makeup.
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In August 2018, the Academy announced that several categories would not be televised live, but rather be recorded during commercial breaks and aired later in the ceremony.[91]
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Following dissent from Academy members, they announced that they would indeed air all 24 categories live. This followed a number of proposals (among them, the introduction of a Popular Film category) that the Academy had announced but did not implement.[92]
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The Board of Governors meets each year and considers new award categories. To date, the following categories have been proposed:
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The Special Academy Awards are voted on by special committees, rather than by the Academy membership as a whole. They are not always presented on an annual basis.
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Due to the positive exposure and prestige of the Academy Awards, many studios spend millions of dollars and hire publicists specifically to promote their films during what is typically called the "Oscar season". This has generated accusations of the Academy Awards being influenced more by marketing than by quality. William Friedkin, an Academy Award-winning film director and former producer of the ceremony, expressed this sentiment at a conference in New York in 2009, describing it as "the greatest promotion scheme that any industry ever devised for itself".[99]
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Tim Dirks, editor of AMC's filmsite.org, has written of the Academy Awards:
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Unfortunately, the critical worth, artistic vision, cultural influence and innovative qualities of many films are not given the same voting weight. Especially since the 1980s, moneymaking "formula-made" blockbusters with glossy production values have often been crowd-pleasing titans (and Best Picture winners), but they haven't necessarily been great films with depth or critical acclaim by any measure.[100]
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A recent technique that has been claimed to be used during the Oscar season is the whisper campaign. These campaigns are intended to spread negative perceptions of other movies nominated and are believed to be perpetrated by those that were involved in creating the movie. Examples of whisper campaigns include the allegations against Zero Dark Thirty suggesting that it justifies torture and the claim that Lincoln distorts history.[101]
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Typical criticism of the Academy Awards for Best Picture is that among the winners and nominees there is an over-representation of romantic historical epics, biographical dramas, romantic dramedies and family melodramas, most of which are released in the U.S. the last three months of the calendar year. The Oscars have been infamously known for selecting specific genres of movies to be awarded. This has led to the coining of the term 'Oscar bait', describing such movies. This has led at times to more specific criticisms that the Academy is disconnected from the audience, e.g., by favoring 'Oscar bait' over audience favorites, or favoring historical melodramas over critically acclaimed movies that depict current life issues.[102]
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The Academy Awards have long received criticism over its lack of diversity among the nominees.[103][104][105] This criticism is based on the statistics from every Academy Awards since 1929, which shows us that only 6.4% of academy award nominees have been non-white and since 1991, 11.2% of nominees have been non-white, with the rate of winners being even more polarizing.[106] More white actresses have won Oscars for yellowface portrayals of Asian characters than actual Asian actresses.[107][108] The 88th awards ceremony became the target of a boycott, popularized on social media with the hashtag #OscarsSoWhite, based on critics' perception that its all-white acting nominee list reflected bias. In response, the Academy initiated "historic" changes in membership by the year 2020.[109][110]
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Acting prizes in certain years have been criticized for not recognizing superior performances so much as being awarded for personal popularity,[111] to make up for a "snub" for a performance/work that proved in time to be more popular and/or renowned than the one actually awarded, or presented as a "career honor" to recognize a distinguished nominee's entire body of work.[27]
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Following the 91st Academy Awards in February 2019 in which the Netflix-broadcast film Roma had been nominated for ten awards including the Best Picture category, Steven Spielberg and other members of the Academy discussed changing the requirements through the Board of Governors for films as to exclude those from Netflix and other media streaming services. Spielberg had been concerned that Netflix as a movie production and distribution studio could spend much more than typical Oscar-winning films and have much wider and earlier distribution than other Best Picture-nominated films, while still being able to meet the minimal theatrical-run status to qualify for an Oscar.[112] The United States Department of Justice, having heard of this potential rule change, wrote a letter to the Academy in March 2019, cautioning them that placing additional restrictions on films that originate from streaming media services without proper justification could raise anti-trust concerns against the Academy.[113] Following its April 2019 board meeting, the Academy Board of Governors agreed to retain the current rules that allow for streaming media films to be eligible for Oscars as long as they enjoy limited theatrical runs.[114]
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Some winners critical of the Academy Awards have boycotted the ceremonies and refused to accept their Oscars. The first to do so was screenwriter Dudley Nichols (Best Writing in 1935 for The Informer). Nichols boycotted the 8th Academy Awards ceremony because of conflicts between the Academy and the Writers' Guild.[115] Nichols eventually accepted the 1935 award three years later, at the 1938 ceremony. Nichols was nominated for three further Academy Awards during his career.
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George C. Scott became the second person to refuse his award (Best Actor in 1970 for Patton) at the 43rd Academy Awards ceremony. Scott described it as a "meat parade", saying, "I don't want any part of it."[116][117][118]
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The third person to refuse the award was Marlon Brando, who refused his award (Best Actor for 1972's The Godfather), citing the film industry's discrimination and mistreatment of Native Americans. At the 45th Academy Awards ceremony, Brando sent actress and civil rights activist Sacheen Littlefeather to read a 15-page speech, detailing his criticisms, which was booed by the audience.[119][115]
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At the 89th Academy Awards ceremony, Warren Beatty and Faye Dunaway mistakenly announced La La Land as the recipient of the Best Picture award, instead of Moonlight, the actual winner. Beatty had been given the wrong envelope and after hesitating during the announcement, handed the envelope to Dunaway, which listed Emma Stone as Best Actress for La La Land and led to the confusion. The proper winner was announced after the acceptance speeches by La La Land producers Fred Berger, Jordan Horowitz and Marc Platt.
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The following year, Beatty and Dunaway were invited back as presenters of the Best Picture award, which they accomplished without error.[120]
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Nine films have been disqualified before an official award ceremony because they violated the regulations:[123]
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One film was disqualified after winning the award, and had the winner return the Oscar:
|
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The following events are closely associated with the annual Academy Awards:
|
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It has become a tradition to give out gift bags to the presenters and performers at the Oscars. In recent years, these gifts have also been extended to award nominees and winners.[124] The value of each of these gift bags can reach into the tens of thousands of dollars. In 2014, the value was reported to be as high as US$80,000.[125] The value has risen to the point where the U.S. Internal Revenue Service issued a statement regarding the gifts and their taxable status.[126]
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Oscar gift bags have included vacation packages to Hawaii and Mexico and Japan, a private dinner party for the recipient and friends at a restaurant, videophones, a four-night stay at a hotel, watches, bracelets, spa treatments, bottles of vodka, maple salad dressing, weight-loss gummie candy and up to $25,000 worth of cosmetic treatments and rejuvenation procedures such as lip fillers and chemical peels from New York City facial plastic surgeon Konstantin Vasyukevich.[124][127][128][129][130] Some of the gifts have even had a "risque" element to them; in 2014, the adult products retailer Adam & Eve had a "Secret Room Gifting Suite". Celebrities visiting the gifting suite included Judith Hoag, Carolyn Hennesy, Kate Linder, Chris Mulkey, Jim O'Heir, and John Salley.[131]
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From 2006 onwards, results are Live+SD; all previous years are live viewing.[132]
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The term "Oscar" is a registered trademark of the AMPAS; however, in the Italian language, it is used generically to refer to any award or award ceremony, regardless of which field.[140][141]
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Green Book
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Parasite
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The Academy Awards, popularly known as the Oscars,[1] are awards for artistic and technical merit in the film industry. Given annually by the Academy of Motion Picture Arts and Sciences (AMPAS), the awards are an international recognition of excellence in cinematic achievements, as assessed by the Academy's voting membership. The various category winners are awarded a copy of a golden statuette, officially called the "Academy Award of Merit", although more commonly referred to by its nickname, the "Oscar". The statuette depicts a knight rendered in the Art Deco style.
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The award was originally sculpted by George Stanley from a design sketch by Cedric Gibbons.[2] AMPAS first presented it in 1929 at a private dinner hosted by Douglas Fairbanks in the Hollywood Roosevelt Hotel in what would become known as the 1st Academy Awards.[3][4] The Academy Awards ceremony was first broadcast by radio in 1930 and was televised for the first time in 1953. It is the oldest worldwide entertainment awards ceremony and is now televised live worldwide.[5] It is also the oldest of the four major annual American entertainment awards; its equivalents – the Emmy Awards for television, the Tony Awards for theater, and the Grammy Awards for music – are modeled after the Academy Awards.[6] They are widely cited as the most famous and prestigious competitive awards in the field of entertainment.
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The 92nd Academy Awards ceremony, honoring the best films of 2019, was held on February 9, 2020. As with the previous ceremony, there was no host. The ceremony was broadcast on ABC. It took place at the Dolby Theatre in Los Angeles, California, for the 18th consecutive year. This ceremony marked the first time in Oscar history that more than three films had received ten or more nominations.[7] A total of 3,140 Oscar statuettes have been awarded since its inception in 1929.[8]
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The first Academy Awards presentation was held on May 16, 1929, at a private dinner function at the Hollywood Roosevelt Hotel with an audience of about 270 people.[9]
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The post-awards party was held at the Mayfair Hotel.[5] The cost of guest tickets for that night's ceremony was $5 ($74 in 2020 dollars). Fifteen statuettes were awarded, honoring artists, directors and other participants in the film-making industry of the time, for their works during the 1927–28 period. The ceremony ran for 15 minutes.
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Winners were announced to media three months earlier. That was changed for the second ceremony in 1930. Since then, for the rest of the first decade, the results were given to newspapers for publication at 11:00 pm on the night of the awards.[5] This method was used until 1940, when the Los Angeles Times announced the winners before the ceremony began; as a result, the Academy has, since 1941, used a sealed envelope to reveal the names of the winners.[5]
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The first Best Actor awarded was Emil Jannings, for his performances in The Last Command and The Way of All Flesh. He had to return to Europe before the ceremony, so the Academy agreed to give him the prize earlier; this made him the first Academy Award winner in history. At that time, winners were recognized for the entirety of their work done in a certain category during the qualifying period; for example, Jannings received the award for two movies in which he starred during that period, and Janet Gaynor later won a single Oscar for performances in three films. With the fourth ceremony, however, the system changed, and professionals were honored for a specific performance in a single film. For the first six ceremonies, the eligibility period spanned two calendar years.[5]
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At the 29th ceremony, held in 1957, the Best Foreign Language Film category, now known as Best International Feature Film, was introduced. Until then, foreign-language films had been honored with the Special Achievement Award.
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Perhaps the most widely seen streaker in history was 34-year-old Robert Opel, who streaked across the stage of The Dorothy Chandler Pavilion in Los Angeles flashing a peace sign on national US television at the 46th Academy Awards in 1974. Bemused host David Niven quipped, "Isn't it fascinating to think that probably the only laugh that man will ever get in his life is by stripping off and showing his shortcomings?" Later, evidence arose suggesting that Opel's appearance was facilitated as a publicity stunt by the show's producer Jack Haley Jr. Robert Metzler, the show's business manager, believed that the incident had been planned in some way; during the dress rehearsal Niven had asked Metzler's wife to borrow a pen so he could write down the famous line, which was thus not the ad-lib it appeared to be.[10]
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The 74th Academy Awards, held in 2002, presented the first Academy Award for Best Animated Feature.[11]
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Since 1973, all Academy Awards ceremonies have ended with the Academy Award for Best Picture.
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Traditionally, the previous year's winner for Best Actor and Best Supporting Actor present the awards for Best Actress and Best Supporting Actress, while the previous year's winner for Best Actress and Best Supporting Actress present the awards for Best Actor and Best Supporting Actor.
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Parasite became the first foreign-language film to win Best Picture at the February 9, 2020 award ceremony.[12]
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Tom Hanks announced at the 2020 Oscar Ceremony, the opening of the Academy Museum of Motion Pictures on December 14, 2020.[13] The museum development started in 2017 under Kerry Brougher, but is now led by Bill Kramer.[14] The industry curated exhibits will be geared toward the history of motion picture, the art & science of film making, exhibiting trailblazing directors, actors, film makers, sound editors and more, and will house famous artifacts from acclaimed movies like Dorothy's Ruby Red Slippers.
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Because of COVID-19, Academy president David Rubin and CEO Dawn Hudson announced that for the 2021 Oscar Ceremony, streaming movies not shown in theaters would be eligible, though at some point the requirement that movies be shown in theaters would return.[15]
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The best known award is the Academy Award of Merit, more popularly known as the Oscar statuette.[8] Made of gold-plated bronze on a black metal base, it is 13.5 in (34.3 cm) tall, weighs 8.5 lb (3.856 kg), and depicts a knight rendered in Art Deco style holding a crusader's sword standing on a reel of film with five spokes. The five spokes represent the original branches of the Academy: Actors, Writers, Directors, Producers, and Technicians.[16]
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Sculptor George Stanley (who also did the Muse Fountain at the Hollywood Bowl) sculpted Cedric Gibbons' design. The statuettes presented at the initial ceremonies were gold-plated solid bronze. Within a few years, the bronze was abandoned in favor of Britannia metal, a pewter-like alloy which is then plated in copper, nickel silver, and finally, 24-karat gold.[8] Due to a metal shortage during World War II, Oscars were made of painted plaster for three years. Following the war, the Academy invited recipients to redeem the plaster figures for gold-plated metal ones.[17] The only addition to the Oscar since it was created is a minor streamlining of the base. The original Oscar mold was cast in 1928 at the C.W. Shumway & Sons Foundry in Batavia, Illinois, which also contributed to casting the molds for the Vince Lombardi Trophy and Emmy Award's statuettes. From 1983 to 2015,[18] approximately 50 Oscars in a tin alloy with gold plating were made each year in Chicago by Illinois manufacturer R.S. Owens & Company.[19] It would take between three and four weeks to manufacture 50 statuettes.[20] In 2016, the Academy returned to bronze as the core metal of the statuettes, handing manufacturing duties to Walden, New York-based Polich Tallix Fine Art Foundry.[21][22] While based on a digital scan of an original 1929 Oscar, the statuettes retain their modern-era dimensions and black pedestal. Cast in liquid bronze from 3D-printed ceramic molds and polished, they are then electroplated in 24-karat gold by Brooklyn, New York–based Epner Technology. The time required to produce 50 such statuettes is roughly three months.[23] R.S. Owens is expected to continue producing other awards for the Academy and service existing Oscars that need replating.[24]
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The Academy officially adopted the name "Oscar" for the trophies in 1939. However, the origin of the nickname is disputed.[25]
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One biography of Bette Davis, who was a president of the Academy in 1941, claims she named the award after her first husband, band leader Harmon Oscar Nelson. A frequently mentioned originator is Margaret Herrick, the Academy executive secretary, who, when she first saw the award in 1931, said the statuette reminded her of "Uncle Oscar", a nickname for her cousin Oscar Pierce.[26]
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Columnist Sidney Skolsky, who was present during Herrick's naming in 1931, wrote that "Employees have affectionately dubbed their famous statuette 'Oscar.'"[27] The Academy credits Skolsky with "the first confirmed newspaper reference" to Oscar in his column on March 16, 1934, which was written about that year's 6th Academy Awards.[28] The 1934 awards appeared again in another early media mention of Oscar: a Time magazine story.[29] In the ceremonies that year, Walt Disney was the first to thank the Academy for his "Oscar" during his acceptance speech.[30]
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To prevent information identifying the Oscar winners from leaking ahead of the ceremony, Oscar statuettes presented at the ceremony have blank baseplates. Until 2010, winners returned their statuettes to the Academy and had to wait several weeks to have their names inscribed on their respective Oscars. Since 2010, winners have had the option of having engraved nameplates applied to their statuettes at an inscription-processing station at the Governor's Ball, a party held immediately after the Oscar ceremony. The R.S. Owens company has engraved nameplates made before the ceremony, bearing the name of every potential winner. The nameplates for the non-winning nominees are later recycled.[31][32]
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Prior to 1950 Oscar statuettes were (and remain) the property of the recipient.[33] Since then the statuettes have been legally encumbered by the requirement that the statuette be first offered for sale back to the Academy for US$1. If a winner refuses to agree to this stipulation, then the Academy keeps the statuette. Academy Awards predating this agreement have been sold in public auctions and private deals for six-figure sums.[34]
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In 1989 Michael Todd's grandson tried to sell Todd's Best Picture Oscar for his 1956 production of Around the World in 80 Days to a movie prop collector. The Academy earned enforcement of its statuette contract by gaining a permanent injunction against the sale.
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In 1992, Harold Russell consigned his 1946 Oscar for Best Supporting Actor for The Best Years of Our Lives to auction in order to raise money for his wife's medical expenses. Though his decision caused controversy, the first-ever Oscar to be sold passed to a private collector on August 6, 1992 for $60,500 ($110,200 today). Russell defended his action, saying, "I don't know why anybody would be critical. My wife's health is much more important than sentimental reasons. The movie will be here, even if Oscar isn't." .[35]
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In December 2011, Orson Welles' 1941 Oscar for Citizen Kane (Academy Award for Best Original Screenplay) was put up for auction, after his heirs won a 2004 court decision contending that Welles did not sign any agreement to return the statue to the Academy.[36] On December 20, 2011, it sold in an online auction for US$861,542 ($0.98 million today).[37]
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Some buyers have subsequently returned the statuettes to the Academy, which keeps them in its treasury.[34]
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In addition to the Academy Award of Merit (Oscar award), there are nine honorary (non-competitive) awards presented by the Academy from time to time (except for the Academy Honorary Award, the Technical Achievement Award, and the Student Academy Awards, which are presented annually):[38]
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The Academy also awards Nicholl Fellowships in Screenwriting.
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Since 2004, Academy Award nomination results have been announced to the public in mid-January. Prior to that, the results were announced in early February.
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The Academy of Motion Picture Arts and Sciences (AMPAS), a professional honorary organization, maintains a voting membership of over 7,000 as of 2018[update].[39]
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Academy membership is divided into different branches, with each representing a different discipline in film production. Actors constitute the largest voting bloc, numbering 1,311 members (22 percent) of the Academy's composition. Votes have been certified by the auditing firm PricewaterhouseCoopers (and its predecessor Price Waterhouse) since the 7th Academy Awards in 1935. The firm mails the ballots of eligible nominees to members of the Academy in December to reflect the previous eligible year with a due date sometime in January of the next year, then tabulates the votes in a process that takes thousands of hours.[40][41][42]
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All AMPAS members must be invited to join by the Board of Governors, on behalf of Academy Branch Executive Committees. Membership eligibility may be achieved by a competitive nomination or a member may submit a name based on other significant contributions to the field of motion pictures.
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New membership proposals are considered annually. The Academy does not publicly disclose its membership, although as recently as 2007 press releases have announced the names of those who have been invited to join. The 2007 release also stated that it has just under 6,000 voting members. While the membership had been growing, stricter policies have kept its size steady since then.[43]
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In 2012, the results of a study conducted by the Los Angeles Times were published describing the demographic breakdown of approximately 88% of AMPAS' voting membership. Of the 5,100+ active voters confirmed, 94% were Caucasian, 77% were male, and 54% were found to be over the age of 60. 33% of voting members are former nominees (14%) and winners (19%).[44]
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In May 2011, the Academy sent a letter advising its 6,000 or so voting members that an online system for Oscar voting would be implemented in 2013.[45]
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According to Rules 2 and 3 of the official Academy Awards Rules, a film must open in the previous calendar year, from midnight at the start of January 1 to midnight at the end of December 31, in Los Angeles County, California, and play for seven consecutive days, to qualify (except for the Best International Feature Film, Best Documentary Feature, and awards in short film categories). Additionally, the film must be shown at least three times on each day of its qualifying run, with at least one of the daily showings starting between 6 pm and 10 pm local time.[46][47]
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The Best International Feature Film award does not require a U.S. release. It requires the film to be submitted as its country's official selection.
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The Best Documentary Feature award requires either week-long releases in both Los Angeles County and New York City[a] during the previous calendar year, or a qualifying award at a competitive film festival from the Documentary Feature Qualifying Festival list (regardless of any public exhibition or distribution), or a submission in the International Feature Film category as its country's official selection. The qualifying theatrical runs must meet the same requirements as those for non-documentary films regarding numbers and times of screenings. Additionally, a film must have been reviewed by a critic from The New York Times, Time Out New York, the Los Angeles Times, or LA Weekly.[49]
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Awards in short film categories (Best Documentary Short Subject, Best Animated Short Film, and Best Live Action Short Film) have noticeably different eligibility rules from most other competitive awards. First, the qualifying period for release does not coincide with a calendar year, instead covering a one-year period starting on October 1 and ending on September 30 of the calendar year before the ceremony. Second, there are multiple methods of qualification. The main method is a week-long theatrical release in either Los Angeles County or New York City during the eligibility period. Films also can qualify by winning specified awards at one of a number of competitive film festivals designated by the Academy, also without regard to prior public distribution. Finally, a film that is selected as a gold, silver, or bronze medal winner in an appropriate category of the immediately previous Student Academy Awards is also eligible (Documentary category for that award, and Animation, Narrative, Alternative, or International for the other awards). The requirements for the qualifying theatrical run are also different from those for other awards. Only one screening per day is required. For the Documentary award, the screening must start between noon and 10 pm local time; for other awards, no specific start time is required, but the film must appear in regular theater listings with dates and screening times.[49][50]
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For example, the 2009 Best Picture winner, The Hurt Locker, was actually first released in 2008, but did not qualify for the 2008 awards, as it did not play its Oscar-qualifying run in Los Angeles until mid-2009, thus qualifying for the 2009 awards. Foreign films must include English subtitles, and each country can submit only one film for consideration in the International Feature Film category per year.[51]
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Rule 2 states that a film must be feature-length, defined as a minimum of 40 minutes, except for short-subject awards, and it must exist either on a 35 mm or 70 mm film print or in 24 frame/s or 48 frame/s progressive scan digital cinema format with a minimum projector resolution of 2048 by 1080 pixels.[52] Since the 90th Academy Awards, presented in 2018, multi-part and limited series have been ineligible for the Best Documentary Feature award. This followed the win of O.J.: Made in America, an eight-hour presentation that was screened in a limited release before being broadcast in five parts on ABC and ESPN, in that category in 2017. The Academy's announcement of the new rule made no direct mention of that film.[53]
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Producers must submit an Official Screen Credits online form before the deadline; in case it is not submitted by the defined deadline, the film will be ineligible for Academy Awards in any year. The form includes the production credits for all related categories. Then, each form is checked and put in a Reminder List of Eligible Releases.
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In late December, ballots, and copies of the Reminder List of Eligible Releases are mailed to around 6,000 active members. For most categories, members from each of the branches vote to determine the nominees only in their respective categories (i.e. only directors vote for directors, writers for writers, actors for actors, etc.). In the special case of Best Picture, all voting members are eligible to select the nominees. In all major categories, a variant of the single transferable vote is used, with each member casting a ballot with up to five nominees (ten for Best Picture) ranked preferentially.[54][55][56] In certain categories, including International Feature Film, Documentary and Animated Feature, nominees are selected by special screening committees made up of members from all branches.
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In most categories, the winner is selected from among the nominees by plurality voting of all members.[54][56] Since 2009, the Best Picture winner has been chosen by instant runoff voting.[56][57] Since 2013, re-weighted range voting has been used to select the nominees for the Best Visual Effects.[58][59]
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Film companies will spend as much as several million dollars on marketing to awards voters for a movie in the running for Best Picture, in attempts to improve chances of receiving Oscars and other movie awards conferred in Oscar season. The Academy enforces rules to limit overt campaigning by its members so as to try to eliminate excesses and prevent the process from becoming undignified. It has an awards czar on staff who advises members on allowed practices and levies penalties on offenders.[60] For example, a producer of the 2009 Best Picture nominee The Hurt Locker was disqualified as a producer in the category when he contacted associates urging them to vote for his film and not another that was seen as the front-runner (The Hurt Locker eventually won).
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The major awards are presented at a live televised ceremony, commonly in late February or early March following the relevant calendar year, and six weeks after the announcement of the nominees. It is the culmination of the film awards season, which usually begins during November or December of the previous year. This is an elaborate extravaganza, with the invited guests walking up the red carpet in the creations of the most prominent fashion designers of the day. Black tie dress is the most common outfit for men, although fashion may dictate not wearing a bow-tie, and musical performers sometimes do not adhere to this. (The artists who recorded the nominees for Best Original Song quite often perform those songs live at the awards ceremony, and the fact that they are performing is often used to promote the television broadcast.)
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The Academy Awards is the world's longest-running awards show televised live from the U.S. to all time zones in North America and worldwide, and gathers billions of viewers elsewhere throughout the world.[61] The Oscars were first televised in 1953 by NBC, which continued to broadcast the event until 1960, when ABC took over, televising the festivities (including the first color broadcast of the event in 1966) through 1970. NBC regained the rights for five years (1971–75), then ABC resumed broadcast duties in 1976 and its current contract with the Academy runs through 2028.[62] The Academy has also produced condensed versions of the ceremony for broadcast in international markets (especially those outside of the Americas) in more desirable local timeslots. The ceremony was broadcast live internationally for the first time via satellite since 1970, but only two South American countries, Chile and Brazil, purchased the rights to air the broadcast. By that time, the television rights to the Academy Awards had been sold in 50 countries. A decade later, the rights were already being sold to 60 countries, and by 1984, the TV rights to the Awards were licensed in 76 countries.
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The ceremonies were moved up from late March/early April to late February, since 2004, to help disrupt and shorten the intense lobbying and ad campaigns associated with Oscar season in the film industry. Another reason was because of the growing TV ratings success coinciding with the NCAA Basketball Tournament, which would cut into the Academy Awards audience. (In 1976 and 1977, ABC's regained Oscars were moved from Tuesday to Monday and went directly opposite NBC's NCAA title game.) The earlier date is also to the advantage of ABC, as it now usually occurs during the highly profitable and important February sweeps period. Some years, the ceremony is moved into the first Sunday of March to avoid a clash with the Winter Olympic Games. Another reason for the move to late February and early March is also to avoid the awards ceremony occurring so close to the religious holidays of Passover and Easter, which for decades had been a grievance from members and the general public.[63] Advertising is somewhat restricted, however, as traditionally no movie studios or competitors of official Academy Award sponsors may advertise during the telecast. The production of the Academy Awards telecast currently holds the distinction of winning the most Emmys in history, with 47 wins and 195 nominations overall since that award's own launch in 1949.[64]
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After many years of being held on Mondays at 9:00 pm Eastern/6:00 p.m Pacific, since the 1999 ceremonies, it was moved to Sundays at 8:30 pm ET/5:30 pm PT.[65] The reasons given for the move were that more viewers would tune in on Sundays, that Los Angeles rush-hour traffic jams could be avoided, and an earlier start time would allow viewers on the East Coast to go to bed earlier.[66] For many years the film industry opposed a Sunday broadcast because it would cut into the weekend box office.[67] In 2010, the Academy contemplated moving the ceremony even further back into January, citing TV viewers' fatigue with the film industry's long awards season. However, such an accelerated schedule would dramatically decrease the voting period for its members, to the point where some voters would only have time to view the contending films streamed on their computers (as opposed to traditionally receiving the films and ballots in the mail). Furthermore, a January ceremony on Sunday would clash with National Football League playoff games.[68] In 2018, the Academy announced that the ceremony would be moved from late February to mid February beginning with the 92nd Academy Awards in 2020.[69]
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Originally scheduled for April 8, 1968, the 40th Academy Awards ceremony was postponed for two days, because of the assassination of Dr. Martin Luther King, Jr.. On March 30, 1981, the 53rd Academy Awards was postponed for one day, after the shooting of President Ronald Reagan and others in Washington, D.C.[70]
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In 1993, an In Memoriam segment was introduced,[71] honoring those who had made a significant contribution to cinema who had died in the preceding 12 months, a selection compiled by a small committee of Academy members.[72] This segment has drawn criticism over the years for the omission of some names. Criticism was also levied for many years regarding another aspect, with the segment having a "popularity contest" feel as the audience varied their applause to those who had died by the subject's cultural impact; the applause has since been muted during the telecast, and the audience is discouraged from clapping during the segment and giving silent reflection instead. This segment was later followed by a commercial break.
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In terms of broadcast length, the ceremony generally averages three and a half hours. The first Oscars, in 1929, lasted 15 minutes. At the other end of the spectrum, the 2002 ceremony lasted four hours and twenty-three minutes.[73][74] In 2010, the organizers of the Academy Awards announced winners' acceptance speeches must not run past 45 seconds. This, according to organizer Bill Mechanic, was to ensure the elimination of what he termed "the single most hated thing on the show" – overly long and embarrassing displays of emotion.[75] In 2016, in a further effort to streamline speeches, winners' dedications were displayed on an on-screen ticker.[76] During the 2018 ceremony, host Jimmy Kimmel acknowledged how long the ceremony had become, by announcing that he would give a brand-new jet ski to whoever gave the shortest speech of the night (a reward won by Mark Bridges when accepting his Best Costume Design award for Phantom Thread).[77] The Wall Street Journal analyzed the average minutes spent across the 2014–2018 telecasts as follows: 14 on song performances; 25 on the hosts' speeches; 38 on prerecorded clips; and 78 on the awards themselves, broken into 24 on the introduction and announcement, 24 on winners walking to the stage, and 30 on their acceptance speeches.[78]
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Although still dominant in ratings, the viewership of the Academy Awards has steadily dropped; the 88th Academy Awards were the lowest-rated in the past eight years (although with increases in male and 18–49 viewership), while the show itself also faced mixed reception. Following the show, Variety reported that ABC was, in negotiating an extension to its contract to broadcast the Oscars, seeking to have more creative control over the broadcast itself. Currently and nominally, AMPAS is responsible for most aspects of the telecast, including the choice of production staff and hosting, although ABC is allowed to have some input on their decisions.[79] In August 2016, AMPAS extended its contract with ABC through 2028: the contract neither contains any notable changes nor gives ABC any further creative control over the telecast.[80]
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Historically, the telecast's viewership is higher when box-office hits are favored to win the Best Picture award. More than 57.25 million viewers tuned to the telecast for the 70th Academy Awards in 1998, the year of Titanic, which generated a box office haul during its initial 1997–98 run of US$600.8 million in the US, a box office record that would remain unsurpassed for years.[81] The 76th Academy Awards ceremony, in which The Lord of the Rings: The Return of the King (pre-telecast box office earnings of US$368 million) received 11 Awards including Best Picture, drew 43.56 million viewers.[82] The most watched ceremony based on Nielsen ratings to date, however, was the 42nd Academy Awards (Best Picture Midnight Cowboy) which drew a 43.4% household rating on April 7, 1970.[83]
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By contrast, ceremonies honoring films that have not performed well at the box office tend to show weaker ratings, despite how much critically acclaimed those films have been. The 78th Academy Awards which awarded low-budget independent film Crash (with a pre-Oscar gross of US$53.4 million) generated an audience of 38.64 million with a household rating of 22.91%.[84] In 2008, the 80th Academy Awards telecast was watched by 31.76 million viewers on average with an 18.66% household rating, the lowest-rated and least-watched ceremony at the time, in spite of celebrating 80 years of the Academy Awards.[85] The Best Picture winner of that particular ceremony was another independent film (No Country for Old Men).
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The 92nd Academy Awards drew an average of 23.6 million viewers -- a new low.
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The Academy Film Archive holds copies of every Academy Awards ceremony since the 1949 Oscars and material on many prior ceremonies, along with ancillary material related to more recent shows. Copies are held on a variety of film, video, and digital formats.[86]
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In 1929, the first Academy Awards were presented at a banquet dinner at the Hollywood Roosevelt Hotel. From 1930 to 1943, the ceremony alternated between two venues: the Ambassador Hotel on Wilshire Boulevard and the Biltmore Hotel in downtown Los Angeles.
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Grauman's Chinese Theatre in Hollywood then hosted the awards from 1944 to 1946, followed by the Shrine Auditorium in Los Angeles from 1947 to 1948. The 21st Academy Awards in 1949 were held at the Academy Award Theatre at what had been the Academy's headquarters on Melrose Avenue in Hollywood.[87]
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From 1950 to 1960, the awards were presented at Hollywood's Pantages Theatre. With the advent of television, the awards from 1953 to 1957 took place simultaneously in Hollywood and New York, first at the NBC International Theatre (1953) and then at the NBC Century Theatre, after which the ceremony took place solely in Los Angeles. The Oscars moved to the Santa Monica Civic Auditorium in Santa Monica, California, in 1961. By 1969, the Academy decided to move the ceremonies back to Downtown Los Angeles, this time to the Dorothy Chandler Pavilion at the Los Angeles County Music Center. In the late 1990s and early 2000s, the ceremony returned to the Shrine.
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In 2002, Hollywood's Dolby Theatre (previously known as the Kodak Theatre) became the presentation's current venue.[88]
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In the first year of the awards, the Best Directing award was split into two categories (Drama and Comedy). At times, the Best Original Score award has also been split into separate categories (Drama and Comedy/Musical). From the 1930s through the 1960s, the Art Direction (now Production Design), Cinematography, and Costume Design awards were likewise split into two categories (black-and-white films and color films). Prior to 2012, the Production Design award was called Art Direction, while the Makeup and Hairstyling award was called Makeup.
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In August 2018, the Academy announced that several categories would not be televised live, but rather be recorded during commercial breaks and aired later in the ceremony.[91]
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Following dissent from Academy members, they announced that they would indeed air all 24 categories live. This followed a number of proposals (among them, the introduction of a Popular Film category) that the Academy had announced but did not implement.[92]
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The Board of Governors meets each year and considers new award categories. To date, the following categories have been proposed:
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The Special Academy Awards are voted on by special committees, rather than by the Academy membership as a whole. They are not always presented on an annual basis.
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Due to the positive exposure and prestige of the Academy Awards, many studios spend millions of dollars and hire publicists specifically to promote their films during what is typically called the "Oscar season". This has generated accusations of the Academy Awards being influenced more by marketing than by quality. William Friedkin, an Academy Award-winning film director and former producer of the ceremony, expressed this sentiment at a conference in New York in 2009, describing it as "the greatest promotion scheme that any industry ever devised for itself".[99]
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Tim Dirks, editor of AMC's filmsite.org, has written of the Academy Awards:
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Unfortunately, the critical worth, artistic vision, cultural influence and innovative qualities of many films are not given the same voting weight. Especially since the 1980s, moneymaking "formula-made" blockbusters with glossy production values have often been crowd-pleasing titans (and Best Picture winners), but they haven't necessarily been great films with depth or critical acclaim by any measure.[100]
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A recent technique that has been claimed to be used during the Oscar season is the whisper campaign. These campaigns are intended to spread negative perceptions of other movies nominated and are believed to be perpetrated by those that were involved in creating the movie. Examples of whisper campaigns include the allegations against Zero Dark Thirty suggesting that it justifies torture and the claim that Lincoln distorts history.[101]
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Typical criticism of the Academy Awards for Best Picture is that among the winners and nominees there is an over-representation of romantic historical epics, biographical dramas, romantic dramedies and family melodramas, most of which are released in the U.S. the last three months of the calendar year. The Oscars have been infamously known for selecting specific genres of movies to be awarded. This has led to the coining of the term 'Oscar bait', describing such movies. This has led at times to more specific criticisms that the Academy is disconnected from the audience, e.g., by favoring 'Oscar bait' over audience favorites, or favoring historical melodramas over critically acclaimed movies that depict current life issues.[102]
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The Academy Awards have long received criticism over its lack of diversity among the nominees.[103][104][105] This criticism is based on the statistics from every Academy Awards since 1929, which shows us that only 6.4% of academy award nominees have been non-white and since 1991, 11.2% of nominees have been non-white, with the rate of winners being even more polarizing.[106] More white actresses have won Oscars for yellowface portrayals of Asian characters than actual Asian actresses.[107][108] The 88th awards ceremony became the target of a boycott, popularized on social media with the hashtag #OscarsSoWhite, based on critics' perception that its all-white acting nominee list reflected bias. In response, the Academy initiated "historic" changes in membership by the year 2020.[109][110]
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Acting prizes in certain years have been criticized for not recognizing superior performances so much as being awarded for personal popularity,[111] to make up for a "snub" for a performance/work that proved in time to be more popular and/or renowned than the one actually awarded, or presented as a "career honor" to recognize a distinguished nominee's entire body of work.[27]
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Following the 91st Academy Awards in February 2019 in which the Netflix-broadcast film Roma had been nominated for ten awards including the Best Picture category, Steven Spielberg and other members of the Academy discussed changing the requirements through the Board of Governors for films as to exclude those from Netflix and other media streaming services. Spielberg had been concerned that Netflix as a movie production and distribution studio could spend much more than typical Oscar-winning films and have much wider and earlier distribution than other Best Picture-nominated films, while still being able to meet the minimal theatrical-run status to qualify for an Oscar.[112] The United States Department of Justice, having heard of this potential rule change, wrote a letter to the Academy in March 2019, cautioning them that placing additional restrictions on films that originate from streaming media services without proper justification could raise anti-trust concerns against the Academy.[113] Following its April 2019 board meeting, the Academy Board of Governors agreed to retain the current rules that allow for streaming media films to be eligible for Oscars as long as they enjoy limited theatrical runs.[114]
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Some winners critical of the Academy Awards have boycotted the ceremonies and refused to accept their Oscars. The first to do so was screenwriter Dudley Nichols (Best Writing in 1935 for The Informer). Nichols boycotted the 8th Academy Awards ceremony because of conflicts between the Academy and the Writers' Guild.[115] Nichols eventually accepted the 1935 award three years later, at the 1938 ceremony. Nichols was nominated for three further Academy Awards during his career.
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George C. Scott became the second person to refuse his award (Best Actor in 1970 for Patton) at the 43rd Academy Awards ceremony. Scott described it as a "meat parade", saying, "I don't want any part of it."[116][117][118]
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The third person to refuse the award was Marlon Brando, who refused his award (Best Actor for 1972's The Godfather), citing the film industry's discrimination and mistreatment of Native Americans. At the 45th Academy Awards ceremony, Brando sent actress and civil rights activist Sacheen Littlefeather to read a 15-page speech, detailing his criticisms, which was booed by the audience.[119][115]
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At the 89th Academy Awards ceremony, Warren Beatty and Faye Dunaway mistakenly announced La La Land as the recipient of the Best Picture award, instead of Moonlight, the actual winner. Beatty had been given the wrong envelope and after hesitating during the announcement, handed the envelope to Dunaway, which listed Emma Stone as Best Actress for La La Land and led to the confusion. The proper winner was announced after the acceptance speeches by La La Land producers Fred Berger, Jordan Horowitz and Marc Platt.
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The following year, Beatty and Dunaway were invited back as presenters of the Best Picture award, which they accomplished without error.[120]
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Nine films have been disqualified before an official award ceremony because they violated the regulations:[123]
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One film was disqualified after winning the award, and had the winner return the Oscar:
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The following events are closely associated with the annual Academy Awards:
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It has become a tradition to give out gift bags to the presenters and performers at the Oscars. In recent years, these gifts have also been extended to award nominees and winners.[124] The value of each of these gift bags can reach into the tens of thousands of dollars. In 2014, the value was reported to be as high as US$80,000.[125] The value has risen to the point where the U.S. Internal Revenue Service issued a statement regarding the gifts and their taxable status.[126]
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Oscar gift bags have included vacation packages to Hawaii and Mexico and Japan, a private dinner party for the recipient and friends at a restaurant, videophones, a four-night stay at a hotel, watches, bracelets, spa treatments, bottles of vodka, maple salad dressing, weight-loss gummie candy and up to $25,000 worth of cosmetic treatments and rejuvenation procedures such as lip fillers and chemical peels from New York City facial plastic surgeon Konstantin Vasyukevich.[124][127][128][129][130] Some of the gifts have even had a "risque" element to them; in 2014, the adult products retailer Adam & Eve had a "Secret Room Gifting Suite". Celebrities visiting the gifting suite included Judith Hoag, Carolyn Hennesy, Kate Linder, Chris Mulkey, Jim O'Heir, and John Salley.[131]
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From 2006 onwards, results are Live+SD; all previous years are live viewing.[132]
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The term "Oscar" is a registered trademark of the AMPAS; however, in the Italian language, it is used generically to refer to any award or award ceremony, regardless of which field.[140][141]
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1 |
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A bone is a rigid organ that constitutes part of the vertebrate skeleton in animals. Bones protect the various organs of the body, produce red and white blood cells, store minerals, provide structure and support for the body, and enable mobility. Bones come in a variety of shapes and sizes and have a complex internal and external structure. They are lightweight yet strong and hard, and serve multiple functions.
|
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|
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Bone tissue (osseous tissue) is a hard tissue, a type of dense connective tissue. It has a honeycomb-like matrix internally, which helps to give the bone rigidity. Bone tissue is made up of different types of bone cells. Osteoblasts and osteocytes are involved in the formation and mineralization of bone; osteoclasts are involved in the resorption of bone tissue. Modified (flattened) osteoblasts become the lining cells that form a protective layer on the bone surface. The mineralised matrix of bone tissue has an organic component of mainly collagen called ossein and an inorganic component of bone mineral made up of various salts. Bone tissue is a mineralized tissue of two types, cortical bone and cancellous bone. Other types of tissue found in bones include bone marrow, endosteum, periosteum, nerves, blood vessels and cartilage.
|
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In the human body at birth, there are approximately 270 bones present; many of these fuse together during development, leaving a total of 206 separate bones in the adult, not counting numerous small sesamoid bones.[1][2] The largest bone in the body is the femur or thigh-bone, and the smallest is the stapes in the middle ear.
|
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The Greek word for bone is ὀστέον ("osteon"), hence the many terms that use it as a prefix—such as osteopathy.
|
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|
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+
Bone is not uniformly solid, but consists of a flexible matrix (about 30%) and bound minerals (about 70%) which are intricately woven and endlessly remodeled by a group of specialized bone cells. Their unique composition and design allows bones to be relatively hard and strong, while remaining lightweight.
|
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Bone matrix is 90 to 95% composed of elastic collagen fibers, also known as ossein,[3] and the remainder is ground substance.[4] The elasticity of collagen improves fracture resistance.[5] The matrix is hardened by the binding of inorganic mineral salt, calcium phosphate, in a chemical arrangement known as calcium hydroxylapatite. It is the bone mineralization that give bones rigidity.
|
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Bone is actively constructed and remodeled throughout life by special bone cells known as osteoblasts and osteoclasts. Within any single bone, the tissue is woven into two main patterns, known as cortical and cancellous bone, and each with different appearance and characteristics.
|
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The hard outer layer of bones is composed of cortical bone, which is also called compact bone as it is much denser than cancellous bone. It forms the hard exterior (cortex) of bones. The cortical bone gives bone its smooth, white, and solid appearance, and accounts for 80% of the total bone mass of an adult human skeleton.[6] It facilitates bone's main functions - to support the whole body, to protect organs, to provide levers for movement, and to store and release chemical elements, mainly calcium. It consists of multiple microscopic columns, each called an osteon or Haversian system. Each column is multiple layers of osteoblasts and osteocytes around a central canal called the haversian canal. Volkmann's canals at right angles connect the osteons together. The columns are metabolically active, and as bone is reabsorbed and created the nature and location of the cells within the osteon will change. Cortical bone is covered by a periosteum on its outer surface, and an endosteum on its inner surface. The endosteum is the boundary between the cortical bone and the cancellous bone.[7] The primary anatomical and functional unit of cortical bone is the osteon.
|
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|
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Cancellous bone, also called trabecular or spongy bone,[7] is the internal tissue of the skeletal bone and is an open cell porous network. Cancellous bone has a higher surface-area-to-volume ratio than cortical bone and it is less dense. This makes it weaker and more flexible. The greater surface area also makes it suitable for metabolic activities such as the exchange of calcium ions. Cancellous bone is typically found at the ends of long bones, near joints and in the interior of vertebrae. Cancellous bone is highly vascular and often contains red bone marrow where hematopoiesis, the production of blood cells, occurs. The primary anatomical and functional unit of cancellous bone is the trabecula. The trabeculae are aligned towards the mechanical load distribution that a bone experiences within long bones such as the femur. As far as short bones are concerned, trabecular alignment has been studied in the vertebral pedicle.[8] Thin formations of osteoblasts covered in endosteum create an irregular network of spaces,[9] known as trabeculae. Within these spaces are bone marrow and hematopoietic stem cells that give rise to platelets, red blood cells and white blood cells.[9] Trabecular marrow is composed of a network of rod- and plate-like elements that make the overall organ lighter and allow room for blood vessels and marrow. Trabecular bone accounts for the remaining 20% of total bone mass but has nearly ten times the surface area of compact bone.[10]
|
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The words cancellous and trabecular refer to the tiny lattice-shaped units (trabeculae) that form the tissue. It was first illustrated accurately in the engravings of Crisóstomo Martinez.[11]
|
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+
|
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+
Bone marrow, also known as myeloid tissue in red bone marrow, can be found in almost any bone that holds cancellous tissue. In newborns, all such bones are filled exclusively with red marrow or hematopoietic marrow, but as the child ages the hematopoietic fraction decreases in quantity and the fatty/ yellow fraction called marrow adipose tissue (MAT) increases in quantity. In adults, red marrow is mostly found in the bone marrow of the femur, the ribs, the vertebrae and pelvic bones.[12]
|
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|
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+
Bone is a metabolically active tissue composed of several types of cells. These cells include osteoblasts, which are involved in the creation and mineralization of bone tissue, osteocytes, and osteoclasts, which are involved in the reabsorption of bone tissue. Osteoblasts and osteocytes are derived from osteoprogenitor cells, but osteoclasts are derived from the same cells that differentiate to form macrophages and monocytes.[13] Within the marrow of the bone there are also hematopoietic stem cells. These cells give rise to other cells, including white blood cells, red blood cells, and platelets.[14]
|
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+
Osteoblasts are mononucleate bone-forming cells. They are located on the surface of osteon seams and make a protein mixture known as osteoid, which mineralizes to become bone.[15] The osteoid seam is a narrow region of newly formed organic matrix, not yet mineralized, located on the surface of a bone. Osteoid is primarily composed of Type I collagen. Osteoblasts also manufacture hormones, such as prostaglandins, to act on the bone itself. The osteoblast creates and repairs new bone by actually building around itself. First, the osteoblast puts up collagen fibers. These collagen fibers are used as a framework for the osteoblasts' work. The osteoblast then deposits calcium phosphate which is hardened by hydroxide and bicarbonate ions. The brand new bone created by the osteoblast is called osteoid.[16] Once the osteoblast is finished working it is actually trapped inside the bone once it hardens. When the osteoblast becomes trapped, it becomes known as an osteocyte.[17] Other osteoblasts remain on the top of the new bone and are used to protect the underlying bone, these become known as lining cells.
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Osteocytes are mostly inactive osteoblasts.[13] Osteocytes originate from osteoblasts that have migrated into and become trapped and surrounded by bone matrix that they themselves produced.[7] The spaces they occupy are known as lacunae. Osteocytes have many processes that reach out to meet osteoblasts and other osteocytes probably for the purposes of communication.[18] Osteocytes remain in contact with other cells in the bone through gap junctions—coupled cell processes—which pass through small channels in the bone matrix called the canaliculi.
|
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|
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+
Osteoclasts are very large multinucleate cells that are responsible for the breakdown of bones by the process of bone resorption. New bone is then formed by the osteoblasts. Bone is constantly remodelled by the resorption of osteoclasts and created by osteoblasts.[13] Osteoclasts are large cells with multiple nuclei located on bone surfaces in what are called Howship's lacunae (or resorption pits). These lacunae are the result of surrounding bone tissue that has been reabsorbed.[19] Because the osteoclasts are derived from a monocyte stem-cell lineage, they are equipped with phagocytic-like mechanisms similar to circulating macrophages.[13] Osteoclasts mature and/or migrate to discrete bone surfaces. Upon arrival, active enzymes, such as tartrate resistant acid phosphatase, are secreted against the mineral substrate.[citation needed] The reabsorption of bone by osteoclasts also plays a role in calcium homeostasis.[19]
|
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Bones consist of living cells embedded in a mineralized organic matrix. This matrix consists of organic components, mainly type I collagen – "organic" referring to materials produced as a result of the human body – and inorganic components, primarily hydroxyapatite and other salts of calcium and phosphate. Above 30% of the acellular part of bone consists of the organic components, and 70% of salts.[20] The collagen fibers give bone its tensile strength, and the interspersed crystals of hydroxyapatite give bone its compressive strength. These effects are synergistic.[20]
|
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The inorganic composition of bone (bone mineral) is primarily formed from salts of calcium and phosphate, the major salt being hydroxyapatite (Ca10(PO4)6(OH)2).[20] The exact composition of the matrix may be subject to change over time due to nutrition and biomineralization, with the ratio of calcium to phosphate varying between 1.3 and 2.0 (per weight), and trace minerals such as magnesium, sodium, potassium and carbonate also being found.[20]
|
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Type I collagen composes 90–95% of the organic matrix, with remainder of the matrix being a homogenous liquid called ground substance consisting of proteoglycans such as hyaluronic acid and chondroitin sulfate,[20] as well as non-collagenous proteins such as osteocalcin, osteopontin or bone sialoprotein. Collagen consists of strands of repeating units, which give bone tensile strength, and are arranged in an overlapping fashion that prevents shear stress. The function of ground substance is not fully known.[20] Two types of bone can be identified microscopically according to the arrangement of collagen: woven and lamellar.
|
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Woven bone is produced when osteoblasts produce osteoid rapidly, which occurs initially in all fetal bones, but is later replaced by more resilient lamellar bone. In adults woven bone is created after fractures or in Paget's disease. Woven bone is weaker, with a smaller number of randomly oriented collagen fibers, but forms quickly; it is for this appearance of the fibrous matrix that the bone is termed woven. It is soon replaced by lamellar bone, which is highly organized in concentric sheets with a much lower proportion of osteocytes to surrounding tissue. Lamellar bone, which makes its first appearance in humans in the fetus during the third trimester,[22] is stronger and filled with many collagen fibers parallel to other fibers in the same layer (these parallel columns are called osteons). In cross-section, the fibers run in opposite directions in alternating layers, much like in plywood, assisting in the bone's ability to resist torsion forces. After a fracture, woven bone forms initially and is gradually replaced by lamellar bone during a process known as "bony substitution." Compared to woven bone, lamellar bone formation takes place more slowly. The orderly deposition of collagen fibers restricts the formation of osteoid to about 1 to 2 µm per day. Lamellar bone also requires a relatively flat surface to lay the collagen fibers in parallel or concentric layers.[23]
|
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+
The extracellular matrix of bone is laid down by osteoblasts, which secrete both collagen and ground substance. These synthesise collagen within the cell, and then secrete collagen fibrils. The collagen fibers rapidly polymerise to form collagen strands. At this stage they are not yet mineralised, and are called "osteoid". Around the strands calcium and phosphate precipitate on the surface of these strands, within days to weeks becoming crystals of hydroxyapatite.[20]
|
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|
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+
In order to mineralise the bone, the osteoblasts secrete vesicles containing alkaline phosphatase. This cleaves the phosphate groups and acts as the foci for calcium and phosphate deposition. The vesicles then rupture and act as a centre for crystals to grow on. More particularly, bone mineral is formed from globular and plate structures.[24][25]
|
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+
|
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+
There are five types of bones in the human body: long, short, flat, irregular, and sesamoid.[26]
|
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+
|
47 |
+
In the study of anatomy, anatomists use a number of anatomical terms to describe the appearance, shape and function of bones. Other anatomical terms are also used to describe the location of bones. Like other anatomical terms, many of these derive from Latin and Greek. Some anatomists still use Latin to refer to bones. The term "osseous", and the prefix "osteo-", referring to things related to bone, are still used commonly today.
|
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|
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+
Some examples of terms used to describe bones include the term "foramen" to describe a hole through which something passes, and a "canal" or "meatus" to describe a tunnel-like structure. A protrusion from a bone can be called a number of terms, including a "condyle", "crest", "spine", "eminence", "tubercle" or "tuberosity", depending on the protrusion's shape and location. In general, long bones are said to have a "head", "neck", and "body".
|
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When two bones join together, they are said to "articulate". If the two bones have a fibrous connection and are relatively immobile, then the joint is called a "suture".
|
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The formation of bone is called ossification. During the fetal stage of development this occurs by two processes: intramembranous ossification and endochondral ossification.[31] Intramembranous ossification involves the formation of bone from connective tissue whereas endochondral ossification involves the formation of bone from cartilage.
|
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Intramembranous ossification mainly occurs during formation of the flat bones of the skull but also the mandible, maxilla, and clavicles; the bone is formed from connective tissue such as mesenchyme tissue rather than from cartilage. The process includes: the development of the ossification center, calcification, trabeculae formation and the development of the periosteum.[32]
|
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Endochondral ossification occurs in long bones and most other bones in the body; it involves the development of bone from cartilage. This process includes the development of a cartilage model, its growth and development, development of the primary and secondary ossification centers, and the formation of articular cartilage and the epiphyseal plates.[33]
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Endochondral ossification begins with points in the cartilage called "primary ossification centers." They mostly appear during fetal development, though a few short bones begin their primary ossification after birth. They are responsible for the formation of the diaphyses of long bones, short bones and certain parts of irregular bones. Secondary ossification occurs after birth, and forms the epiphyses of long bones and the extremities of irregular and flat bones. The diaphysis and both epiphyses of a long bone are separated by a growing zone of cartilage (the epiphyseal plate). At skeletal maturity (18 to 25 years of age), all of the cartilage is replaced by bone, fusing the diaphysis and both epiphyses together (epiphyseal closure).[34] In the upper limbs, only the diaphyses of the long bones and scapula are ossified. The epiphyses, carpal bones, coracoid process, medial border of the scapula, and acromion are still cartilaginous.[35]
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The following steps are followed in the conversion of cartilage to bone:
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Bones have a variety of functions:
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Bones serve a variety of mechanical functions. Together the bones in the body form the skeleton. They provide a frame to keep the body supported, and an attachment point for skeletal muscles, tendons, ligaments and joints, which function together to generate and transfer forces so that individual body parts or the whole body can be manipulated in three-dimensional space (the interaction between bone and muscle is studied in biomechanics).
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Bones protect internal organs, such as the skull protecting the brain or the ribs protecting the heart and lungs. Because of the way that bone is formed, bone has a high compressive strength of about 170 MPa (1800 kgf/cm²),[5] poor tensile strength of 104–121 MPa, and a very low shear stress strength (51.6 MPa).[37][38] This means that bone resists pushing (compressional) stress well, resist pulling (tensional) stress less well, but only poorly resists shear stress (such as due to torsional loads). While bone is essentially brittle, bone does have a significant degree of elasticity, contributed chiefly by collagen.
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Mechanically, bones also have a special role in hearing. The ossicles are three small bones in the middle ear which are involved in sound transduction.
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+
The cancellous part of bones contain bone marrow. Bone marrow produces blood cells in a process called hematopoiesis.[39] Blood cells that are created in bone marrow include red blood cells, platelets and white blood cells.[40] Progenitor cells such as the hematopoietic stem cell divide in a process called mitosis to produce precursor cells. These include precursors which eventually give rise to white blood cells, and erythroblasts which give rise to red blood cells.[41] Unlike red and white blood cells, created by mitosis, platelets are shed from very large cells called megakaryocytes.[42] This process of progressive differentiation occurs within the bone marrow. After the cells are matured, they enter the circulation.[43] Every day, over 2.5 billion red blood cells and platelets, and 50–100 billion granulocytes are produced in this way.[14]
|
72 |
+
|
73 |
+
As well as creating cells, bone marrow is also one of the major sites where defective or aged red blood cells are destroyed.[14]
|
74 |
+
|
75 |
+
°Determined by the species, age, and the type of bone, bone cells make up to 15 percent of the bone. Growth factor storage – mineralized bone matrix stores important growth factors such as insulin-like growth factors, transforming growth factor, bone morphogenetic proteins and others.[46]
|
76 |
+
|
77 |
+
Bone is constantly being created and replaced in a process known as remodeling. This ongoing turnover of bone is a process of resorption followed by replacement of bone with little change in shape. This is accomplished through osteoblasts and osteoclasts. Cells are stimulated by a variety of signals, and together referred to as a remodeling unit. Approximately 10% of the skeletal mass of an adult is remodelled each year.[52] The purpose of remodeling is to regulate calcium homeostasis, repair microdamaged bones from everyday stress, and to shape the skeleton during growth.[citation needed] Repeated stress, such as weight-bearing exercise or bone healing, results in the bone thickening at the points of maximum stress (Wolff's law). It has been hypothesized that this is a result of bone's piezoelectric properties, which cause bone to generate small electrical potentials under stress.[53]
|
78 |
+
|
79 |
+
The action of osteoblasts and osteoclasts are controlled by a number of chemical enzymes that either promote or inhibit the activity of the bone remodeling cells, controlling the rate at which bone is made, destroyed, or changed in shape. The cells also use paracrine signalling to control the activity of each other.[citation needed] For example, the rate at which osteoclasts resorb bone is inhibited by calcitonin and osteoprotegerin. Calcitonin is produced by parafollicular cells in the thyroid gland, and can bind to receptors on osteoclasts to directly inhibit osteoclast activity. Osteoprotegerin is secreted by osteoblasts and is able to bind RANK-L, inhibiting osteoclast stimulation.[54]
|
80 |
+
|
81 |
+
Osteoblasts can also be stimulated to increase bone mass through increased secretion of osteoid and by inhibiting the ability of osteoclasts to break down osseous tissue.[citation needed] Increased secretion of osteoid is stimulated by the secretion of growth hormone by the pituitary, thyroid hormone and the sex hormones (estrogens and androgens). These hormones also promote increased secretion of osteoprotegerin.[54] Osteoblasts can also be induced to secrete a number of cytokines that promote reabsorption of bone by stimulating osteoclast activity and differentiation from progenitor cells. Vitamin D, parathyroid hormone and stimulation from osteocytes induce osteoblasts to increase secretion of RANK-ligand and interleukin 6, which cytokines then stimulate increased reabsorption of bone by osteoclasts. These same compounds also increase secretion of macrophage colony-stimulating factor by osteoblasts, which promotes the differentiation of progenitor cells into osteoclasts, and decrease secretion of osteoprotegerin.[citation needed]
|
82 |
+
|
83 |
+
Bone volume is determined by the rates of bone formation and bone resorption. Recent research has suggested that certain growth factors may work to locally alter bone formation by increasing osteoblast activity. Numerous bone-derived growth factors have been isolated and classified via bone cultures. These factors include insulin-like growth factors I and II, transforming growth factor-beta, fibroblast growth factor, platelet-derived growth factor, and bone morphogenetic proteins.[55] Evidence suggests that bone cells produce growth factors for extracellular storage in the bone matrix. The release of these growth factors from the bone matrix could cause the proliferation of osteoblast precursors. Essentially, bone growth factors may act as potential determinants of local bone formation.[55] Research has suggested that cancellous bone volume in postmenopausal osteoporosis may be determined by the relationship between the total bone forming surface and the percent of surface resorption.[56]
|
84 |
+
|
85 |
+
A number of diseases can affect bone, including arthritis, fractures, infections, osteoporosis and tumours. Conditions relating to bone can be managed by a variety of doctors, including rheumatologists for joints, and orthopedic surgeons, who may conduct surgery to fix broken bones. Other doctors, such as rehabilitation specialists may be involved in recovery, radiologists in interpreting the findings on imaging, and pathologists in investigating the cause of the disease, and family doctors may play a role in preventing complications of bone disease such as osteoporosis.
|
86 |
+
|
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+
When a doctor sees a patient, a history and exam will be taken. Bones are then often imaged, called radiography. This might include ultrasound X-ray, CT scan, MRI scan and other imaging such as a Bone scan, which may be used to investigate cancer.[57] Other tests such as a blood test for autoimmune markers may be taken, or a synovial fluid aspirate may be taken.[57]
|
88 |
+
|
89 |
+
In normal bone, fractures occur when there is significant force applied, or repetitive trauma over a long time. Fractures can also occur when a bone is weakened, such as with osteoporosis, or when there is a structural problem, such as when the bone remodels excessively (such as Paget's disease) or is the site of the growth of cancer.[58] Common fractures include wrist fractures and hip fractures, associated with osteoporosis, vertebral fractures associated with high-energy trauma and cancer, and fractures of long-bones. Not all fractures are painful.[58] When serious, depending on the fractures type and location, complications may include flail chest, compartment syndromes or fat embolism.
|
90 |
+
Compound fractures involve the bone's penetration through the skin. Some complex fractures can be treated by the use of bone grafting procedures that replace missing bone portions.
|
91 |
+
|
92 |
+
Fractures and their underlying causes can be investigated by X-rays, CT scans and MRIs.[58] Fractures are described by their location and shape, and several classification systems exist, depending on the location of the fracture. A common long bone fracture in children is a Salter–Harris fracture.[59] When fractures are managed, pain relief is often given, and the fractured area is often immobilised. This is to promote bone healing. In addition, surgical measures such as internal fixation may be used. Because of the immobilisation, people with fractures are often advised to undergo rehabilitation.[58]
|
93 |
+
|
94 |
+
There are several types of tumour that can affect bone; examples of benign bone tumours include osteoma, osteoid osteoma, osteochondroma, osteoblastoma, enchondroma, giant cell tumour of bone, and aneurysmal bone cyst.[60]
|
95 |
+
|
96 |
+
Cancer can arise in bone tissue, and bones are also a common site for other cancers to spread (metastasise) to.[61] Cancers that arise in bone are called "primary" cancers, although such cancers are rare.[61] Metastases within bone are "secondary" cancers, with the most common being breast cancer, lung cancer, prostate cancer, thyroid cancer, and kidney cancer.[61] Secondary cancers that affect bone can either destroy bone (called a "lytic" cancer) or create bone (a "sclerotic" cancer). Cancers of the bone marrow inside the bone can also affect bone tissue, examples including leukemia and multiple myeloma. Bone may also be affected by cancers in other parts of the body. Cancers in other parts of the body may release parathyroid hormone or parathyroid hormone-related peptide. This increases bone reabsorption, and can lead to bone fractures.
|
97 |
+
|
98 |
+
Bone tissue that is destroyed or altered as a result of cancers is distorted, weakened, and more prone to fracture. This may lead to compression of the spinal cord, destruction of the marrow resulting in bruising, bleeding and immunosuppression, and is one cause of bone pain. If the cancer is metastatic, then there might be other symptoms depending on the site of the original cancer. Some bone cancers can also be felt.
|
99 |
+
|
100 |
+
Cancers of the bone are managed according to their type, their stage, prognosis, and what symptoms they cause. Many primary cancers of bone are treated with radiotherapy. Cancers of bone marrow may be treated with chemotherapy, and other forms of targeted therapy such as immunotherapy may be used.[62] Palliative care, which focuses on maximising a person's quality of life, may play a role in management, particularly if the likelihood of survival within five years is poor.
|
101 |
+
|
102 |
+
Osteoporosis is a disease of bone where there is reduced bone mineral density, increasing the likelihood of fractures.[63] Osteoporosis is defined in women by the World Health Organization as a bone mineral density of 2.5 standard deviations below peak bone mass, relative to the age and sex-matched average. This density is measured using dual energy X-ray absorptiometry (DEXA), with the term "established osteoporosis" including the presence of a fragility fracture.[64] Osteoporosis is most common in women after menopause, when it is called "postmenopausal osteoporosis", but may develop in men and premenopausal women in the presence of particular hormonal disorders and other chronic diseases or as a result of smoking and medications, specifically glucocorticoids.[63] Osteoporosis usually has no symptoms until a fracture occurs.[63] For this reason, DEXA scans are often done in people with one or more risk factors, who have developed osteoporosis and are at risk of fracture.[63]
|
103 |
+
|
104 |
+
Osteoporosis treatment includes advice to stop smoking, decrease alcohol consumption, exercise regularly, and have a healthy diet. Calcium and trace mineral supplements may also be advised, as may Vitamin D. When medication is used, it may include bisphosphonates, Strontium ranelate, and hormone replacement therapy.[65]
|
105 |
+
|
106 |
+
Osteopathic medicine is a school of medical thought originally developed based on the idea of the link between the musculoskeletal system and overall health, but now very similar to mainstream medicine. As of 2012[update], over 77,000 physicians in the United States are trained in osteopathic medical schools.[66]
|
107 |
+
|
108 |
+
The study of bones and teeth is referred to as osteology. It is frequently used in anthropology, archeology and forensic science for a variety of tasks. This can include determining the nutritional, health, age or injury status of the individual the bones were taken from. Preparing fleshed bones for these types of studies can involve the process of maceration.
|
109 |
+
|
110 |
+
Typically anthropologists and archeologists study bone tools made by Homo sapiens and Homo neanderthalensis. Bones can serve a number of uses such as projectile points or artistic pigments, and can also be made from external bones such as antlers.
|
111 |
+
|
112 |
+
Bird skeletons are very lightweight. Their bones are smaller and thinner, to aid flight. Among mammals, bats come closest to birds in terms of bone density, suggesting that small dense bones are a flight adaptation. Many bird bones have little marrow due to their being hollow.[67]
|
113 |
+
|
114 |
+
A bird's beak is primarily made of bone as projections of the mandibles which are covered in keratin.
|
115 |
+
|
116 |
+
A deer's antlers are composed of bone which is an unusual example of bone being outside the skin of the animal once the velvet is shed.[68]
|
117 |
+
|
118 |
+
The extinct predatory fish Dunkleosteus had sharp edges of hard exposed bone along its jaws.[69][70]
|
119 |
+
|
120 |
+
Many animals possess an exoskeleton that is not made of bone. These include insects and crustaceans.
|
121 |
+
|
122 |
+
The proportion of cortical bone that is 80% in the human skeleton may be much lower in other animals, especially in marine mammals and marine turtles, or in various Mesozoic marine reptiles, such as ichthyosaurs,[71] among others.[72]
|
123 |
+
|
124 |
+
Many animals, particularly herbivores, practice osteophagy – the eating of bones. This is presumably carried out in order to replenish lacking phosphate.
|
125 |
+
|
126 |
+
Many bone diseases that affect humans also affect other vertebrates – an example of one disorder is skeletal fluorosis.
|
127 |
+
|
128 |
+
Bones from slaughtered animals have a number of uses. In prehistoric times, they have been used for making bone tools.[73] They have further been used in bone carving, already important in prehistoric art, and also in modern time as crafting materials for buttons, beads, handles, bobbins, calculation aids, head nuts, dice, poker chips, pick-up sticks, ornaments, etc. A special genre is scrimshaw.[citation needed]
|
129 |
+
|
130 |
+
Bone glue can be made by prolonged boiling of ground or cracked bones, followed by filtering and evaporation to thicken the resulting fluid. Historically once important, bone glue and other animal glues today have only a few specialized uses, such as in antiques restoration. Essentially the same process, with further refinement, thickening and drying, is used to make gelatin.
|
131 |
+
|
132 |
+
Broth is made by simmering several ingredients for a long time, traditionally including bones.
|
133 |
+
|
134 |
+
Bone char, a porous, black, granular material primarily used for filtration and also as a black pigment, is produced by charring mammal bones.
|
135 |
+
|
136 |
+
Oracle bone script was a writing system used in Ancient China based on inscriptions in bones. Its name originates from oracle bones, which were mainly ox clavicle. The Ancient Chinese (mainly in the Shang Dynasty), would write their questions on the Oracle Bone, and burn the bone, and where the bone cracked would be the answer for the questions.
|
137 |
+
|
138 |
+
To point the bone at someone is considered bad luck in some cultures, such as Australian aborigines, such as by the Kurdaitcha.
|
139 |
+
|
140 |
+
The wishbones of fowl have been used for divination, and are still customarily used in a tradition to determine which one of two people pulling on either prong of the bone may make a wish.
|
141 |
+
|
142 |
+
Various cultures throughout history have adopted the custom of shaping an infant's head by the practice of artificial cranial deformation. A widely practised
|
143 |
+
custom in China was that of foot binding to limit the normal growth of the foot.
|
144 |
+
|
145 |
+
Cells in bone marrow
|
146 |
+
|
147 |
+
Scanning electron microscope of bone at 100x magnification
|
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+
|
149 |
+
Structure detail of an animal bone
|
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1 |
+
|
2 |
+
|
3 |
+
An operating system (OS) is system software that manages computer hardware, software resources, and provides common services for computer programs.
|
4 |
+
|
5 |
+
Time-sharing operating systems schedule tasks for efficient use of the system and may also include accounting software for cost allocation of processor time, mass storage, printing, and other resources.
|
6 |
+
|
7 |
+
For hardware functions such as input and output and memory allocation, the operating system acts as an intermediary between programs and the computer hardware,[1][2] although the application code is usually executed directly by the hardware and frequently makes system calls to an OS function or is interrupted by it. Operating systems are found on many devices that contain a computer – from cellular phones and video game consoles to web servers and supercomputers.
|
8 |
+
|
9 |
+
The dominant desktop operating system is Microsoft Windows with a market share of around 82.74%. macOS by Apple Inc. is in second place (13.23%), and the varieties of Linux are collectively in third place (1.57%).[3] In the mobile sector (including smartphones and tablets), Android's share is up to 70% in the year 2017.[4] According to third quarter 2016 data, Android's share on smartphones is dominant with 87.5 percent with also a growth rate of 10.3 percent per year, followed by Apple's iOS with 12.1 percent with per year decrease in market share of 5.2 percent, while other operating systems amount to just 0.3 percent.[5] Linux distributions are dominant in the server and supercomputing sectors. Other specialized classes of operating systems, such as embedded and real-time systems, exist for many applications.
|
10 |
+
|
11 |
+
A single-tasking system can only run one program at a time, while a multi-tasking operating system allows more than one program to be running in concurrency. This is achieved by time-sharing, where the available processor time is divided between multiple processes. These processes are each interrupted repeatedly in time slices by a task-scheduling subsystem of the operating system. Multi-tasking may be characterized in preemptive and co-operative types. In preemptive multitasking, the operating system slices the CPU time and dedicates a slot to each of the programs. Unix-like operating systems, such as Solaris and Linux—as well as non-Unix-like, such as AmigaOS—support preemptive multitasking. Cooperative multitasking is achieved by relying on each process to provide time to the other processes in a defined manner. 16-bit versions of Microsoft Windows used cooperative multi-tasking; 32-bit versions of both Windows NT and Win9x used preemptive multi-tasking.
|
12 |
+
|
13 |
+
Single-user operating systems have no facilities to distinguish users, but may allow multiple programs to run in tandem.[6] A multi-user operating system extends the basic concept of multi-tasking with facilities that identify processes and resources, such as disk space, belonging to multiple users, and the system permits multiple users to interact with the system at the same time. Time-sharing operating systems schedule tasks for efficient use of the system and may also include accounting software for cost allocation of processor time, mass storage, printing, and other resources to multiple users.
|
14 |
+
|
15 |
+
A distributed operating system manages a group of distinct, networked computers and makes them appear to be a single computer, as all computations are distributed (divided amongst the constituent computers).[7]
|
16 |
+
|
17 |
+
In the distributed and cloud computing context of an OS, templating refers to creating a single virtual machine image as a guest operating system, then saving it as a tool for multiple running virtual machines. The technique is used both in virtualization and cloud computing management, and is common in large server warehouses.[8]
|
18 |
+
|
19 |
+
Embedded operating systems are designed to be used in embedded computer systems. They are designed to operate on small machines with less autonomy (e.g. PDAs). They are very compact and extremely efficient by design, and are able to operate with a limited amount of resources. Windows CE and Minix 3 are some examples of embedded operating systems.
|
20 |
+
|
21 |
+
A real-time operating system is an operating system that guarantees to process events or data by a specific moment in time. A real-time operating system may be single- or multi-tasking, but when multitasking, it uses specialized scheduling algorithms so that a deterministic nature of behavior is achieved. Such an event-driven system switches between tasks based on their priorities or external events, whereas time-sharing operating systems switch tasks based on clock interrupts.
|
22 |
+
|
23 |
+
A library operating system is one in which the services that a typical operating system provides, such as networking, are provided in the form of libraries and composed with the application and configuration code to construct a unikernel: a specialized, single address space, machine image that can be deployed to cloud or embedded environments.
|
24 |
+
|
25 |
+
Early computers were built to perform a series of single tasks, like a calculator. Basic operating system features were developed in the 1950s, such as resident monitor functions that could automatically run different programs in succession to speed up processing. Operating systems did not exist in their modern and more complex forms until the early 1960s.[9] Hardware features were added, that enabled use of runtime libraries, interrupts, and parallel processing. When personal computers became popular in the 1980s, operating systems were made for them similar in concept to those used on larger computers.
|
26 |
+
|
27 |
+
In the 1940s, the earliest electronic digital systems had no operating systems. Electronic systems of this time were programmed on rows of mechanical switches or by jumper wires on plugboards. These were special-purpose systems that, for example, generated ballistics tables for the military or controlled the printing of payroll checks from data on punched paper cards. After programmable general purpose computers were invented, machine languages (consisting of strings of the binary digits 0 and 1 on punched paper tape) were introduced that sped up the programming process (Stern, 1981).[full citation needed]
|
28 |
+
|
29 |
+
In the early 1950s, a computer could execute only one program at a time. Each user had sole use of the computer for a limited period and would arrive at a scheduled time with their program and data on punched paper cards or punched tape. The program would be loaded into the machine, and the machine would be set to work until the program completed or crashed. Programs could generally be debugged via a front panel using toggle switches and panel lights. It is said that Alan Turing was a master of this on the early Manchester Mark 1 machine, and he was already deriving the primitive conception of an operating system from the principles of the universal Turing machine.[9]
|
30 |
+
|
31 |
+
Later machines came with libraries of programs, which would be linked to a user's program to assist in operations such as input and output and compiling (generating machine code from human-readable symbolic code). This was the genesis of the modern-day operating system. However, machines still ran a single job at a time. At Cambridge University in England, the job queue was at one time a washing line (clothes line) from which tapes were hung with different colored clothes-pegs to indicate job priority.[citation needed]
|
32 |
+
|
33 |
+
An improvement was the Atlas Supervisor. Introduced with the Manchester Atlas in 1962, it is considered by many to be the first recognisable modern operating system.[10] Brinch Hansen described it as "the most significant breakthrough in the history of operating systems."[11]
|
34 |
+
|
35 |
+
Through the 1950s, many major features were pioneered in the field of operating systems on mainframe computers, including batch processing, input/output interrupting, buffering, multitasking, spooling, runtime libraries, link-loading, and programs for sorting records in files. These features were included or not included in application software at the option of application programmers, rather than in a separate operating system used by all applications. In 1959, the SHARE Operating System was released as an integrated utility for the IBM 704, and later in the 709 and 7090 mainframes, although it was quickly supplanted by IBSYS/IBJOB on the 709, 7090 and 7094.
|
36 |
+
|
37 |
+
During the 1960s, IBM's OS/360 introduced the concept of a single OS spanning an entire product line, which was crucial for the success of the System/360 machines. IBM's current mainframe operating systems are distant descendants of this original system and modern machines are backwards-compatible with applications written for OS/360.[citation needed]
|
38 |
+
|
39 |
+
OS/360 also pioneered the concept that the operating system keeps track of all of the system resources that are used, including program and data space allocation in main memory and file space in secondary storage, and file locking during updates. When a process is terminated for any reason, all of these resources are re-claimed by the operating system.
|
40 |
+
|
41 |
+
The alternative CP-67 system for the S/360-67 started a whole line of IBM operating systems focused on the concept of virtual machines. Other operating systems used on IBM S/360 series mainframes included systems developed by IBM: COS/360 (Compatibility Operating System), DOS/360 (Disk Operating System), TSS/360 (Time Sharing System), TOS/360 (Tape Operating System), BOS/360 (Basic Operating System), and ACP (Airline Control Program), as well as a few non-IBM systems: MTS (Michigan Terminal System), MUSIC (Multi-User System for Interactive Computing), and ORVYL (Stanford Timesharing System).
|
42 |
+
|
43 |
+
Control Data Corporation developed the SCOPE operating system in the 1960s, for batch processing. In cooperation with the University of Minnesota, the Kronos and later the NOS operating systems were developed during the 1970s, which supported simultaneous batch and timesharing use. Like many commercial timesharing systems, its interface was an extension of the Dartmouth BASIC operating systems, one of the pioneering efforts in timesharing and programming languages. In the late 1970s, Control Data and the University of Illinois developed the PLATO operating system, which used plasma panel displays and long-distance time sharing networks. Plato was remarkably innovative for its time, featuring real-time chat, and multi-user graphical games.
|
44 |
+
|
45 |
+
In 1961, Burroughs Corporation introduced the B5000 with the MCP (Master Control Program) operating system. The B5000 was a stack machine designed to exclusively support high-level languages with no machine language or assembler; indeed, the MCP was the first OS to be written exclusively in a high-level language (ESPOL, a dialect of ALGOL). MCP also introduced many other ground-breaking innovations, such as being the first commercial implementation of virtual memory. During development of the AS/400, IBM made an approach to Burroughs to license MCP to run on the AS/400 hardware. This proposal was declined by Burroughs management to protect its existing hardware production. MCP is still in use today in the Unisys company's ClearPath/MCP line of computers.
|
46 |
+
|
47 |
+
UNIVAC, the first commercial computer manufacturer, produced a series of EXEC operating systems[citation needed]. Like all early main-frame systems, this batch-oriented system managed magnetic drums, disks, card readers and line printers. In the 1970s, UNIVAC produced the Real-Time Basic (RTB) system to support large-scale time sharing, also patterned after the Dartmouth BC system.
|
48 |
+
|
49 |
+
General Electric and MIT developed General Electric Comprehensive Operating Supervisor (GECOS), which introduced the concept of ringed security privilege levels. After acquisition by Honeywell it was renamed General Comprehensive Operating System (GCOS).
|
50 |
+
|
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+
Digital Equipment Corporation developed many operating systems for its various computer lines, including TOPS-10 and TOPS-20 time sharing systems for the 36-bit PDP-10 class systems. Before the widespread use of UNIX, TOPS-10 was a particularly popular system in universities, and in the early ARPANET community. RT-11 was a single-user real-time OS for the PDP-11 class minicomputer, and RSX-11 was the corresponding multi-user OS.
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From the late 1960s through the late 1970s, several hardware capabilities evolved that allowed similar or ported software to run on more than one system. Early systems had utilized microprogramming to implement features on their systems in order to permit different underlying computer architectures to appear to be the same as others in a series. In fact, most 360s after the 360/40 (except the 360/165 and 360/168) were microprogrammed implementations.
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The enormous investment in software for these systems made since the 1960s caused most of the original computer manufacturers to continue to develop compatible operating systems along with the hardware. Notable supported mainframe operating systems include:
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The first microcomputers did not have the capacity or need for the elaborate operating systems that had been developed for mainframes and minis; minimalistic operating systems were developed, often loaded from ROM and known as monitors. One notable early disk operating system was CP/M, which was supported on many early microcomputers and was closely imitated by Microsoft's MS-DOS, which became widely popular as the operating system chosen for the IBM PC (IBM's version of it was called IBM DOS or PC DOS). In the 1980s, Apple Computer Inc. (now Apple Inc.) abandoned its popular Apple II series of microcomputers to introduce the Apple Macintosh computer with an innovative graphical user interface (GUI) to the Mac OS.
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The introduction of the Intel 80386 CPU chip in October 1985,[12] with 32-bit architecture and paging capabilities, provided personal computers with the ability to run multitasking operating systems like those of earlier minicomputers and mainframes. Microsoft responded to this progress by hiring Dave Cutler, who had developed the VMS operating system for Digital Equipment Corporation. He would lead the development of the Windows NT operating system, which continues to serve as the basis for Microsoft's operating systems line. Steve Jobs, a co-founder of Apple Inc., started NeXT Computer Inc., which developed the NEXTSTEP operating system. NEXTSTEP would later be acquired by Apple Inc. and used, along with code from FreeBSD as the core of Mac OS X (macOS after latest name change).
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The GNU Project was started by activist and programmer Richard Stallman with the goal of creating a complete free software replacement to the proprietary UNIX operating system. While the project was highly successful in duplicating the functionality of various parts of UNIX, development of the GNU Hurd kernel proved to be unproductive. In 1991, Finnish computer science student Linus Torvalds, with cooperation from volunteers collaborating over the Internet, released the first version of the Linux kernel. It was soon merged with the GNU user space components and system software to form a complete operating system. Since then, the combination of the two major components has usually been referred to as simply "Linux" by the software industry, a naming convention that Stallman and the Free Software Foundation remain opposed to, preferring the name GNU/Linux. The Berkeley Software Distribution, known as BSD, is the UNIX derivative distributed by the University of California, Berkeley, starting in the 1970s. Freely distributed and ported to many minicomputers, it eventually also gained a following for use on PCs, mainly as FreeBSD, NetBSD and OpenBSD.
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Unix was originally written in assembly language.[13] Ken Thompson wrote B, mainly based on BCPL, based on his experience in the MULTICS project. B was replaced by C, and Unix, rewritten in C, developed into a large, complex family of inter-related operating systems which have been influential in every modern operating system (see History).
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The Unix-like family is a diverse group of operating systems, with several major sub-categories including System V, BSD, and Linux. The name "UNIX" is a trademark of The Open Group which licenses it for use with any operating system that has been shown to conform to their definitions. "UNIX-like" is commonly used to refer to the large set of operating systems which resemble the original UNIX.
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Unix-like systems run on a wide variety of computer architectures. They are used heavily for servers in business, as well as workstations in academic and engineering environments. Free UNIX variants, such as Linux and BSD, are popular in these areas.
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Four operating systems are certified by The Open Group (holder of the Unix trademark) as Unix. HP's HP-UX and IBM's AIX are both descendants of the original System V Unix and are designed to run only on their respective vendor's hardware. In contrast, Sun Microsystems's Solaris can run on multiple types of hardware, including x86 and Sparc servers, and PCs. Apple's macOS, a replacement for Apple's earlier (non-Unix) Mac OS, is a hybrid kernel-based BSD variant derived from NeXTSTEP, Mach, and FreeBSD.
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Unix interoperability was sought by establishing the POSIX standard. The POSIX standard can be applied to any operating system, although it was originally created for various Unix variants.
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A subgroup of the Unix family is the Berkeley Software Distribution family, which includes FreeBSD, NetBSD, and OpenBSD. These operating systems are most commonly found on webservers, although they can also function as a personal computer OS. The Internet owes much of its existence to BSD, as many of the protocols now commonly used by computers to connect, send and receive data over a network were widely implemented and refined in BSD. The World Wide Web was also first demonstrated on a number of computers running an OS based on BSD called NeXTSTEP.
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In 1974, University of California, Berkeley installed its first Unix system. Over time, students and staff in the computer science department there began adding new programs to make things easier, such as text editors. When Berkeley received new VAX computers in 1978 with Unix installed, the school's undergraduates modified Unix even more in order to take advantage of the computer's hardware possibilities. The Defense Advanced Research Projects Agency of the US Department of Defense took interest, and decided to fund the project. Many schools, corporations, and government organizations took notice and started to use Berkeley's version of Unix instead of the official one distributed by AT&T.
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Steve Jobs, upon leaving Apple Inc. in 1985, formed NeXT Inc., a company that manufactured high-end computers running on a variation of BSD called NeXTSTEP. One of these computers was used by Tim Berners-Lee as the first webserver to create the World Wide Web.
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Developers like Keith Bostic encouraged the project to replace any non-free code that originated with Bell Labs. Once this was done, however, AT&T sued. After two years of legal disputes, the BSD project spawned a number of free derivatives, such as NetBSD and FreeBSD (both in 1993), and OpenBSD (from NetBSD in 1995).
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macOS (formerly "Mac OS X" and later "OS X") is a line of open core graphical operating systems developed, marketed, and sold by Apple Inc., the latest of which is pre-loaded on all currently shipping Macintosh computers. macOS is the successor to the original classic Mac OS, which had been Apple's primary operating system since 1984. Unlike its predecessor, macOS is a UNIX operating system built on technology that had been developed at NeXT through the second half of the 1980s and up until Apple purchased the company in early 1997.
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The operating system was first released in 1999 as Mac OS X Server 1.0, followed in March 2001 by a client version (Mac OS X v10.0 "Cheetah"). Since then, six more distinct "client" and "server" editions of macOS have been released, until the two were merged in OS X 10.7 "Lion".
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Prior to its merging with macOS, the server edition – macOS Server – was architecturally identical to its desktop counterpart and usually ran on Apple's line of Macintosh server hardware. macOS Server included work group management and administration software tools that provide simplified access to key network services, including a mail transfer agent, a Samba server, an LDAP server, a domain name server, and others. With Mac OS X v10.7 Lion, all server aspects of Mac OS X Server have been integrated into the client version and the product re-branded as "OS X" (dropping "Mac" from the name). The server tools are now offered as an application.[14]
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The Linux kernel originated in 1991, as a project of Linus Torvalds, while a university student in Finland. He posted information about his project on a newsgroup for computer students and programmers, and received support and assistance from volunteers who succeeded in creating a complete and functional kernel.
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Linux is Unix-like, but was developed without any Unix code, unlike BSD and its variants. Because of its open license model, the Linux kernel code is available for study and modification, which resulted in its use on a wide range of computing machinery from supercomputers to smart-watches. Although estimates suggest that Linux is used on only 1.82% of all "desktop" (or laptop) PCs,[15] it has been widely adopted for use in servers[16] and embedded systems[17] such as cell phones. Linux has superseded Unix on many platforms and is used on most supercomputers including the top 385.[18] Many of the same computers are also on Green500 (but in different order), and Linux runs on the top 10. Linux is also commonly used on other small energy-efficient computers, such as smartphones and smartwatches. The Linux kernel is used in some popular distributions, such as Red Hat, Debian, Ubuntu, Linux Mint and Google's Android, Chrome OS, and Chromium OS.
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Microsoft Windows is a family of proprietary operating systems designed by Microsoft Corporation and primarily targeted to Intel architecture based computers, with an estimated 88.9 percent total usage share on Web connected computers.[15][19][20][21] The latest version is Windows 10.
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In 2011, Windows 7 overtook Windows XP as most common version in use.[22][23][24]
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Microsoft Windows was first released in 1985, as an operating environment running on top of MS-DOS, which was the standard operating system shipped on most Intel architecture personal computers at the time. In 1995, Windows 95 was released which only used MS-DOS as a bootstrap. For backwards compatibility, Win9x could run real-mode MS-DOS[25][26] and 16-bit Windows 3.x[27] drivers. Windows ME, released in 2000, was the last version in the Win9x family. Later versions have all been based on the Windows NT kernel. Current client versions of Windows run on IA-32, x86-64 and 32-bit ARM microprocessors.[28] In addition Itanium is still supported in older server version Windows Server 2008 R2. In the past, Windows NT supported additional architectures.
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Server editions of Windows are widely used. In recent years, Microsoft has expended significant capital in an effort to promote the use of Windows as a server operating system. However, Windows' usage on servers is not as widespread as on personal computers as Windows competes against Linux and BSD for server market share.[29][30]
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ReactOS is a Windows-alternative operating system, which is being developed on the principles of Windows – without using any of Microsoft's code.
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There have been many operating systems that were significant in their day but are no longer so, such as AmigaOS; OS/2 from IBM and Microsoft; classic Mac OS, the non-Unix precursor to Apple's macOS; BeOS; XTS-300; RISC OS; MorphOS; Haiku; BareMetal and FreeMint. Some are still used in niche markets and continue to be developed as minority platforms for enthusiast communities and specialist applications. OpenVMS, formerly from DEC, is still under active development by Hewlett-Packard. Yet other operating systems are used almost exclusively in academia, for operating systems education or to do research on operating system concepts. A typical example of a system that fulfills both roles is MINIX, while for example Singularity is used purely for research. Another example is the Oberon System designed at ETH Zürich by Niklaus Wirth, Jürg Gutknecht and a group of students at the former Computer Systems Institute in the 1980s. It was used mainly for research, teaching, and daily work in Wirth's group.
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Other operating systems have failed to win significant market share, but have introduced innovations that have influenced mainstream operating systems, not least Bell Labs' Plan 9.
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The components of an operating system all exist in order to make the different parts of a computer work together. All user software needs to go through the operating system in order to use any of the hardware, whether it be as simple as a mouse or keyboard or as complex as an Internet component.
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With the aid of the firmware and device drivers, the kernel provides the most basic level of control over all of the computer's hardware devices. It manages memory access for programs in the RAM, it determines which programs get access to which hardware resources, it sets up or resets the CPU's operating states for optimal operation at all times, and it organizes the data for long-term non-volatile storage with file systems on such media as disks, tapes, flash memory, etc.
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The operating system provides an interface between an application program and the computer hardware, so that an application program can interact with the hardware only by obeying rules and procedures programmed into the operating system. The operating system is also a set of services which simplify development and execution of application programs. Executing an application program involves the creation of a process by the operating system kernel which assigns memory space and other resources, establishes a priority for the process in multi-tasking systems, loads program binary code into memory, and initiates execution of the application program which then interacts with the user and with hardware devices.
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Interrupts are central to operating systems, as they provide an efficient way for the operating system to interact with and react to its environment. The alternative – having the operating system "watch" the various sources of input for events (polling) that require action – can be found in older systems with very small stacks (50 or 60 bytes) but is unusual in modern systems with large stacks. Interrupt-based programming is directly supported by most modern CPUs. Interrupts provide a computer with a way of automatically saving local register contexts, and running specific code in response to events. Even very basic computers support hardware interrupts, and allow the programmer to specify code which may be run when that event takes place.
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When an interrupt is received, the computer's hardware automatically suspends whatever program is currently running, saves its status, and runs computer code previously associated with the interrupt; this is analogous to placing a bookmark in a book in response to a phone call. In modern operating systems, interrupts are handled by the operating system's kernel. Interrupts may come from either the computer's hardware or the running program.
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When a hardware device triggers an interrupt, the operating system's kernel decides how to deal with this event, generally by running some processing code. The amount of code being run depends on the priority of the interrupt (for example: a person usually responds to a smoke detector alarm before answering the phone). The processing of hardware interrupts is a task that is usually delegated to software called a device driver, which may be part of the operating system's kernel, part of another program, or both. Device drivers may then relay information to a running program by various means.
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A program may also trigger an interrupt to the operating system. If a program wishes to access hardware, for example, it may interrupt the operating system's kernel, which causes control to be passed back to the kernel. The kernel then processes the request. If a program wishes additional resources (or wishes to shed resources) such as memory, it triggers an interrupt to get the kernel's attention.
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Modern microprocessors (CPU or MPU) support multiple modes of operation. CPUs with this capability offer at least two modes: user mode and supervisor mode. In general terms, supervisor mode operation allows unrestricted access to all machine resources, including all MPU instructions. User mode operation sets limits on instruction use and typically disallows direct access to machine resources. CPUs might have other modes similar to user mode as well, such as the virtual modes in order to emulate older processor types, such as 16-bit processors on a 32-bit one, or 32-bit processors on a 64-bit one.
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At power-on or reset, the system begins in supervisor mode. Once an operating system kernel has been loaded and started, the boundary between user mode and supervisor mode (also known as kernel mode) can be established.
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Supervisor mode is used by the kernel for low level tasks that need unrestricted access to hardware, such as controlling how memory is accessed, and communicating with devices such as disk drives and video display devices. User mode, in contrast, is used for almost everything else. Application programs, such as word processors and database managers, operate within user mode, and can only access machine resources by turning control over to the kernel, a process which causes a switch to supervisor mode. Typically, the transfer of control to the kernel is achieved by executing a software interrupt instruction, such as the Motorola 68000 TRAP instruction. The software interrupt causes the microprocessor to switch from user mode to supervisor mode and begin executing code that allows the kernel to take control.
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In user mode, programs usually have access to a restricted set of microprocessor instructions, and generally cannot execute any instructions that could potentially cause disruption to the system's operation. In supervisor mode, instruction execution restrictions are typically removed, allowing the kernel unrestricted access to all machine resources.
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The term "user mode resource" generally refers to one or more CPU registers, which contain information that the running program isn't allowed to alter. Attempts to alter these resources generally causes a switch to supervisor mode, where the operating system can deal with the illegal operation the program was attempting, for example, by forcibly terminating ("killing") the program).
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Among other things, a multiprogramming operating system kernel must be responsible for managing all system memory which is currently in use by programs. This ensures that a program does not interfere with memory already in use by another program. Since programs time share, each program must have independent access to memory.
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Cooperative memory management, used by many early operating systems, assumes that all programs make voluntary use of the kernel's memory manager, and do not exceed their allocated memory. This system of memory management is almost never seen any more, since programs often contain bugs which can cause them to exceed their allocated memory. If a program fails, it may cause memory used by one or more other programs to be affected or overwritten. Malicious programs or viruses may purposefully alter another program's memory, or may affect the operation of the operating system itself. With cooperative memory management, it takes only one misbehaved program to crash the system.
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Memory protection enables the kernel to limit a process' access to the computer's memory. Various methods of memory protection exist, including memory segmentation and paging. All methods require some level of hardware support (such as the 80286 MMU), which doesn't exist in all computers.
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In both segmentation and paging, certain protected mode registers specify to the CPU what memory address it should allow a running program to access. Attempts to access other addresses trigger an interrupt which cause the CPU to re-enter supervisor mode, placing the kernel in charge. This is called a segmentation violation or Seg-V for short, and since it is both difficult to assign a meaningful result to such an operation, and because it is usually a sign of a misbehaving program, the kernel generally resorts to terminating the offending program, and reports the error.
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Windows versions 3.1 through ME had some level of memory protection, but programs could easily circumvent the need to use it. A general protection fault would be produced, indicating a segmentation violation had occurred; however, the system would often crash anyway.
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The use of virtual memory addressing (such as paging or segmentation) means that the kernel can choose what memory each program may use at any given time, allowing the operating system to use the same memory locations for multiple tasks.
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If a program tries to access memory that isn't in its current range of accessible memory, but nonetheless has been allocated to it, the kernel is interrupted in the same way as it would if the program were to exceed its allocated memory. (See section on memory management.) Under UNIX this kind of interrupt is referred to as a page fault.
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When the kernel detects a page fault it generally adjusts the virtual memory range of the program which triggered it, granting it access to the memory requested. This gives the kernel discretionary power over where a particular application's memory is stored, or even whether or not it has actually been allocated yet.
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In modern operating systems, memory which is accessed less frequently can be temporarily stored on disk or other media to make that space available for use by other programs. This is called swapping, as an area of memory can be used by multiple programs, and what that memory area contains can be swapped or exchanged on demand.
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"Virtual memory" provides the programmer or the user with the perception that there is a much larger amount of RAM in the computer than is really there.[31]
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Multitasking refers to the running of multiple independent computer programs on the same computer; giving the appearance that it is performing the tasks at the same time. Since most computers can do at most one or two things at one time, this is generally done via time-sharing, which means that each program uses a share of the computer's time to execute.
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An operating system kernel contains a scheduling program which determines how much time each process spends executing, and in which order execution control should be passed to programs. Control is passed to a process by the kernel, which allows the program access to the CPU and memory. Later, control is returned to the kernel through some mechanism, so that another program may be allowed to use the CPU. This so-called passing of control between the kernel and applications is called a context switch.
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An early model which governed the allocation of time to programs was called cooperative multitasking. In this model, when control is passed to a program by the kernel, it may execute for as long as it wants before explicitly returning control to the kernel. This means that a malicious or malfunctioning program may not only prevent any other programs from using the CPU, but it can hang the entire system if it enters an infinite loop.
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Modern operating systems extend the concepts of application preemption to device drivers and kernel code, so that the operating system has preemptive control over internal run-times as well.
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The philosophy governing preemptive multitasking is that of ensuring that all programs are given regular time on the CPU. This implies that all programs must be limited in how much time they are allowed to spend on the CPU without being interrupted. To accomplish this, modern operating system kernels make use of a timed interrupt. A protected mode timer is set by the kernel which triggers a return to supervisor mode after the specified time has elapsed. (See above sections on Interrupts and Dual Mode Operation.)
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On many single user operating systems cooperative multitasking is perfectly adequate, as home computers generally run a small number of well tested programs. The AmigaOS is an exception, having preemptive multitasking from its very first version. Windows NT was the first version of Microsoft Windows which enforced preemptive multitasking, but it didn't reach the home user market until Windows XP (since Windows NT was targeted at professionals).
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Access to data stored on disks is a central feature of all operating systems. Computers store data on disks using files, which are structured in specific ways in order to allow for faster access, higher reliability, and to make better use of the drive's available space. The specific way in which files are stored on a disk is called a file system, and enables files to have names and attributes. It also allows them to be stored in a hierarchy of directories or folders arranged in a directory tree.
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Early operating systems generally supported a single type of disk drive and only one kind of file system. Early file systems were limited in their capacity, speed, and in the kinds of file names and directory structures they could use. These limitations often reflected limitations in the operating systems they were designed for, making it very difficult for an operating system to support more than one file system.
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While many simpler operating systems support a limited range of options for accessing storage systems, operating systems like UNIX and Linux support a technology known as a virtual file system or VFS. An operating system such as UNIX supports a wide array of storage devices, regardless of their design or file systems, allowing them to be accessed through a common application programming interface (API). This makes it unnecessary for programs to have any knowledge about the device they are accessing. A VFS allows the operating system to provide programs with access to an unlimited number of devices with an infinite variety of file systems installed on them, through the use of specific device drivers and file system drivers.
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A connected storage device, such as a hard drive, is accessed through a device driver. The device driver understands the specific language of the drive and is able to translate that language into a standard language used by the operating system to access all disk drives. On UNIX, this is the language of block devices.
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When the kernel has an appropriate device driver in place, it can then access the contents of the disk drive in raw format, which may contain one or more file systems. A file system driver is used to translate the commands used to access each specific file system into a standard set of commands that the operating system can use to talk to all file systems. Programs can then deal with these file systems on the basis of filenames, and directories/folders, contained within a hierarchical structure. They can create, delete, open, and close files, as well as gather various information about them, including access permissions, size, free space, and creation and modification dates.
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Various differences between file systems make supporting all file systems difficult. Allowed characters in file names, case sensitivity, and the presence of various kinds of file attributes makes the implementation of a single interface for every file system a daunting task. Operating systems tend to recommend using (and so support natively) file systems specifically designed for them; for example, NTFS in Windows and ext3 and ReiserFS in Linux. However, in practice, third party drivers are usually available to give support for the most widely used file systems in most general-purpose operating systems (for example, NTFS is available in Linux through NTFS-3g, and ext2/3 and ReiserFS are available in Windows through third-party software).
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Support for file systems is highly varied among modern operating systems, although there are several common file systems which almost all operating systems include support and drivers for. Operating systems vary on file system support and on the disk formats they may be installed on. Under Windows, each file system is usually limited in application to certain media; for example, CDs must use ISO 9660 or UDF, and as of Windows Vista, NTFS is the only file system which the operating system can be installed on. It is possible to install Linux onto many types of file systems. Unlike other operating systems, Linux and UNIX allow any file system to be used regardless of the media it is stored in, whether it is a hard drive, a disc (CD, DVD...), a USB flash drive, or even contained within a file located on another file system.
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A device driver is a specific type of computer software developed to allow interaction with hardware devices. Typically this constitutes an interface for communicating with the device, through the specific computer bus or communications subsystem that the hardware is connected to, providing commands to and/or receiving data from the device, and on the other end, the requisite interfaces to the operating system and software applications. It is a specialized hardware-dependent computer program which is also operating system specific that enables another program, typically an operating system or applications software package or computer program running under the operating system kernel, to interact transparently with a hardware device, and usually provides the requisite interrupt handling necessary for any necessary asynchronous time-dependent hardware interfacing needs.
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The key design goal of device drivers is abstraction. Every model of hardware (even within the same class of device) is different. Newer models also are released by manufacturers that provide more reliable or better performance and these newer models are often controlled differently. Computers and their operating systems cannot be expected to know how to control every device, both now and in the future. To solve this problem, operating systems essentially dictate how every type of device should be controlled. The function of the device driver is then to translate these operating system mandated function calls into device specific calls. In theory a new device, which is controlled in a new manner, should function correctly if a suitable driver is available. This new driver ensures that the device appears to operate as usual from the operating system's point of view.
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Under versions of Windows before Vista and versions of Linux before 2.6, all driver execution was co-operative, meaning that if a driver entered an infinite loop it would freeze the system. More recent revisions of these operating systems incorporate kernel preemption, where the kernel interrupts the driver to give it tasks, and then separates itself from the process until it receives a response from the device driver, or gives it more tasks to do.
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Currently most operating systems support a variety of networking protocols, hardware, and applications for using them. This means that computers running dissimilar operating systems can participate in a common network for sharing resources such as computing, files, printers, and scanners using either wired or wireless connections. Networks can essentially allow a computer's operating system to access the resources of a remote computer to support the same functions as it could if those resources were connected directly to the local computer. This includes everything from simple communication, to using networked file systems or even sharing another computer's graphics or sound hardware. Some network services allow the resources of a computer to be accessed transparently, such as SSH which allows networked users direct access to a computer's command line interface.
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Client/server networking allows a program on a computer, called a client, to connect via a network to another computer, called a server. Servers offer (or host) various services to other network computers and users. These services are usually provided through ports or numbered access points beyond the server's IP address. Each port number is usually associated with a maximum of one running program, which is responsible for handling requests to that port. A daemon, being a user program, can in turn access the local hardware resources of that computer by passing requests to the operating system kernel.
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Many operating systems support one or more vendor-specific or open networking protocols as well, for example, SNA on IBM systems, DECnet on systems from Digital Equipment Corporation, and Microsoft-specific protocols (SMB) on Windows. Specific protocols for specific tasks may also be supported such as NFS for file access. Protocols like ESound, or esd can be easily extended over the network to provide sound from local applications, on a remote system's sound hardware.
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A computer being secure depends on a number of technologies working properly. A modern operating system provides access to a number of resources, which are available to software running on the system, and to external devices like networks via the kernel.[citation needed]
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The operating system must be capable of distinguishing between requests which should be allowed to be processed, and others which should not be processed. While some systems may simply distinguish between "privileged" and "non-privileged", systems commonly have a form of requester identity, such as a user name. To establish identity there may be a process of authentication. Often a username must be quoted, and each username may have a password. Other methods of authentication, such as magnetic cards or biometric data, might be used instead. In some cases, especially connections from the network, resources may be accessed with no authentication at all (such as reading files over a network share). Also covered by the concept of requester identity is authorization; the particular services and resources accessible by the requester once logged into a system are tied to either the requester's user account or to the variously configured groups of users to which the requester belongs.[citation needed]
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In addition to the allow or disallow model of security, a system with a high level of security also offers auditing options. These would allow tracking of requests for access to resources (such as, "who has been reading this file?"). Internal security, or security from an already running program is only possible if all possibly harmful requests must be carried out through interrupts to the operating system kernel. If programs can directly access hardware and resources, they cannot be secured.[citation needed]
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External security involves a request from outside the computer, such as a login at a connected console or some kind of network connection. External requests are often passed through device drivers to the operating system's kernel, where they can be passed onto applications, or carried out directly. Security of operating systems has long been a concern because of highly sensitive data held on computers, both of a commercial and military nature. The United States Government Department of Defense (DoD) created the Trusted Computer System Evaluation Criteria (TCSEC) which is a standard that sets basic requirements for assessing the effectiveness of security. This became of vital importance to operating system makers, because the TCSEC was used to evaluate, classify and select trusted operating systems being considered for the processing, storage and retrieval of sensitive or classified information.
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Network services include offerings such as file sharing, print services, email, web sites, and file transfer protocols (FTP), most of which can have compromised security. At the front line of security are hardware devices known as firewalls or intrusion detection/prevention systems. At the operating system level, there are a number of software firewalls available, as well as intrusion detection/prevention systems. Most modern operating systems include a software firewall, which is enabled by default. A software firewall can be configured to allow or deny network traffic to or from a service or application running on the operating system. Therefore, one can install and be running an insecure service, such as Telnet or FTP, and not have to be threatened by a security breach because the firewall would deny all traffic trying to connect to the service on that port.
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An alternative strategy, and the only sandbox strategy available in systems that do not meet the Popek and Goldberg virtualization requirements, is where the operating system is not running user programs as native code, but instead either emulates a processor or provides a host for a p-code based system such as Java.
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Internal security is especially relevant for multi-user systems; it allows each user of the system to have private files that the other users cannot tamper with or read. Internal security is also vital if auditing is to be of any use, since a program can potentially bypass the operating system, inclusive of bypassing auditing.
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Every computer that is to be operated by an individual requires a user interface. The user interface is usually referred to as a shell and is essential if human interaction is to be supported. The user interface views the directory structure and requests services from the operating system that will acquire data from input hardware devices, such as a keyboard, mouse or credit card reader, and requests operating system services to display prompts, status messages and such on output hardware devices, such as a video monitor or printer. The two most common forms of a user interface have historically been the command-line interface, where computer commands are typed out line-by-line, and the graphical user interface, where a visual environment (most commonly a WIMP) is present.
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Most of the modern computer systems support graphical user interfaces (GUI), and often include them. In some computer systems, such as the original implementation of the classic Mac OS, the GUI is integrated into the kernel.
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While technically a graphical user interface is not an operating system service, incorporating support for one into the operating system kernel can allow the GUI to be more responsive by reducing the number of context switches required for the GUI to perform its output functions. Other operating systems are modular, separating the graphics subsystem from the kernel and the Operating System. In the 1980s UNIX, VMS and many others had operating systems that were built this way. Linux and macOS are also built this way. Modern releases of Microsoft Windows such as Windows Vista implement a graphics subsystem that is mostly in user-space; however the graphics drawing routines of versions between Windows NT 4.0 and Windows Server 2003 exist mostly in kernel space. Windows 9x had very little distinction between the interface and the kernel.
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Many computer operating systems allow the user to install or create any user interface they desire. The X Window System in conjunction with GNOME or KDE Plasma 5 is a commonly found setup on most Unix and Unix-like (BSD, Linux, Solaris) systems. A number of Windows shell replacements have been released for Microsoft Windows, which offer alternatives to the included Windows shell, but the shell itself cannot be separated from Windows.
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Numerous Unix-based GUIs have existed over time, most derived from X11. Competition among the various vendors of Unix (HP, IBM, Sun) led to much fragmentation, though an effort to standardize in the 1990s to COSE and CDE failed for various reasons, and were eventually eclipsed by the widespread adoption of GNOME and K Desktop Environment. Prior to free software-based toolkits and desktop environments, Motif was the prevalent toolkit/desktop combination (and was the basis upon which CDE was developed).
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Graphical user interfaces evolve over time. For example, Windows has modified its user interface almost every time a new major version of Windows is released, and the Mac OS GUI changed dramatically with the introduction of Mac OS X in 1999.[32]
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A real-time operating system (RTOS) is an operating system intended for applications with fixed deadlines (real-time computing). Such applications include some small embedded systems, automobile engine controllers, industrial robots, spacecraft, industrial control, and some large-scale computing systems.
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An early example of a large-scale real-time operating system was Transaction Processing Facility developed by American Airlines and IBM for the Sabre Airline Reservations System.
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Embedded systems that have fixed deadlines use a real-time operating system such as VxWorks, PikeOS, eCos, QNX, MontaVista Linux and RTLinux. Windows CE is a real-time operating system that shares similar APIs to desktop Windows but shares none of desktop Windows' codebase.[33] Symbian OS also has an RTOS kernel (EKA2) starting with version 8.0b.
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Some embedded systems use operating systems such as Palm OS, BSD, and Linux, although such operating systems do not support real-time computing.
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Operating system development is one of the most complicated activities in which a computing hobbyist may engage.[citation needed] A hobby operating system may be classified as one whose code has not been directly derived from an existing operating system, and has few users and active developers.[34]
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In some cases, hobby development is in support of a "homebrew" computing device, for example, a simple single-board computer powered by a 6502 microprocessor. Or, development may be for an architecture already in widespread use. Operating system development may come from entirely new concepts, or may commence by modeling an existing operating system. In either case, the hobbyist is his/her own developer, or may interact with a small and sometimes unstructured group of individuals who have like interests.
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Examples of a hobby operating system include Syllable and TempleOS.
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Application software is generally written for use on a specific operating system, and sometimes even for specific hardware.[citation needed] When porting the application to run on another OS, the functionality required by that application may be implemented differently by that OS (the names of functions, meaning of arguments, etc.) requiring the application to be adapted, changed, or otherwise maintained.
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Unix was the first operating system not written in assembly language, making it very portable to systems different from its native PDP-11.[35]
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This cost in supporting operating systems diversity can be avoided by instead writing applications against software platforms such as Java or Qt. These abstractions have already borne the cost of adaptation to specific operating systems and their system libraries.
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Another approach is for operating system vendors to adopt standards. For example, POSIX and OS abstraction layers provide commonalities that reduce porting costs.
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Osmium (from Greek ὀσμή osme, "smell") is a chemical element with the symbol Os and atomic number 76. It is a hard, brittle, bluish-white transition metal in the platinum group that is found as a trace element in alloys, mostly in platinum ores. Osmium is the densest naturally occurring element, with an experimentally measured (using x-ray crystallography) density of 22.59 g/cm3. Manufacturers use its alloys with platinum, iridium, and other platinum-group metals to make fountain pen nib tipping, electrical contacts, and in other applications that require extreme durability and hardness.[3] The element's abundance in the Earth's crust is among the rarest.[4][5]
|
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|
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+
Osmium has a blue-gray tint and is the densest stable element; it is approximately twice as dense as lead[3] and slightly denser than iridium.[6] Calculations of density from the X-ray diffraction data may produce the most reliable data for these elements, giving a value of 22.587±0.009 g/cm3 for osmium, slightly denser than the 22.562±0.009 g/cm3 of iridium; both metals are nearly 23 times as dense as water, and 1 1⁄6 times as dense as gold.[7]
|
6 |
+
|
7 |
+
Osmium is a hard but brittle metal that remains lustrous even at high temperatures. It has a very low compressibility. Correspondingly, its bulk modulus is extremely high, reported between 395 and 462 GPa, which rivals that of diamond (443 GPa). The hardness of osmium is moderately high at 4 GPa.[8][9][10] Because of its hardness, brittleness, low vapor pressure (the lowest of the platinum-group metals), and very high melting point (the third highest of all elements, after only tungsten, and rhenium), solid osmium is difficult to machine, form, or work.
|
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+
|
9 |
+
Osmium forms compounds with oxidation states ranging from −2 to +8. The most common oxidation states are +2, +3, +4, and +8. The +8 oxidation state is notable for being the highest attained by any chemical element aside from iridium's +9[11] and is encountered only in xenon,[12][13] ruthenium,[14] hassium,[15] and iridium.[dead link][16] The oxidation states −1 and −2 represented by the two reactive compounds Na2[Os4(CO)13] and Na2[Os(CO)4] are used in the synthesis of osmium cluster compounds.[17][18]
|
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+
|
11 |
+
The most common compound exhibiting the +8 oxidation state is osmium tetroxide. This toxic compound is formed when powdered osmium is exposed to air. It is a very volatile, water-soluble, pale yellow, crystalline solid with a strong smell. Osmium powder has the characteristic smell of osmium tetroxide.[19] Osmium tetroxide forms red osmates OsO4(OH)2−2 upon reaction with a base. With ammonia, it forms the nitrido-osmates OsO3N−.[20][21][22] Osmium tetroxide boils at 130 °C and is a powerful oxidizing agent. By contrast, osmium dioxide (OsO2) is black, non-volatile, and much less reactive and toxic.
|
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+
|
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+
Only two osmium compounds have major applications: osmium tetroxide for staining tissue in electron microscopy and for the oxidation of alkenes in organic synthesis, and the non-volatile osmates for organic oxidation reactions.[23]
|
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+
|
15 |
+
Osmium pentafluoride (OsF5) is known, but osmium trifluoride (OsF3) has not yet been synthesized. The lower oxidation states are stabilized by the larger halogens, so that the trichloride, tribromide, triiodide, and even diiodide are known. The oxidation state +1 is known only for osmium iodide (OsI), whereas several carbonyl complexes of osmium, such as triosmium dodecacarbonyl (Os3(CO)12), represent oxidation state 0.[20][21][24][25]
|
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+
|
17 |
+
In general, the lower oxidation states of osmium are stabilized by ligands that are good σ-donors (such as amines) and π-acceptors (heterocycles containing nitrogen). The higher oxidation states are stabilized by strong σ- and π-donors, such as O2− and N3−.[26]
|
18 |
+
|
19 |
+
Despite its broad range of compounds in numerous oxidation states, osmium in bulk form at ordinary temperatures and pressures resists attack by all acids, including aqua regia, but is attacked by fused alkalis.[27]
|
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+
|
21 |
+
Osmium has seven naturally occurring isotopes, six of which are stable: 184Os, 187Os, 188Os, 189Os, 190Os, and (most abundant) 192Os. 186Os undergoes alpha decay with such a long half-life (2.0±1.1)×1015 years, approximately 140000 times the age of the universe, that for practical purposes it can be considered stable. Alpha decay is predicted for all seven naturally occurring isotopes, but it has been observed only for 186Os, presumably due to very long half-lives. It is predicted that 184Os and 192Os can undergo double beta decay but this radioactivity has not been observed yet.[28]
|
22 |
+
|
23 |
+
187Os is the descendant of 187Re (half-life 4.56×1010 years) and is used extensively in dating terrestrial as well as meteoric rocks (see rhenium-osmium dating). It has also been used to measure the intensity of continental weathering over geologic time and to fix minimum ages for stabilization of the mantle roots of continental cratons. This decay is a reason why rhenium-rich minerals are abnormally rich in 187Os.[29] However, the most notable application of osmium isotopes in geology has been in conjunction with the abundance of iridium, to characterise the layer of shocked quartz along the Cretaceous–Paleogene boundary that marks the extinction of the non-avian dinosaurs 65 million years ago.[30]
|
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+
|
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+
Osmium was discovered in 1803 by Smithson Tennant and William Hyde Wollaston in London, England.[31] The discovery of osmium is intertwined with that of platinum and the other metals of the platinum group. Platinum reached Europe as platina ("small silver"), first encountered in the late 17th century in silver mines around the Chocó Department, in Colombia.[32] The discovery that this metal was not an alloy, but a distinct new element, was published in 1748.[33]
|
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+
Chemists who studied platinum dissolved it in aqua regia (a mixture of hydrochloric and nitric acids) to create soluble salts. They always observed a small amount of a dark, insoluble residue.[34] Joseph Louis Proust thought that the residue was graphite.[34] Victor Collet-Descotils, Antoine François, comte de Fourcroy, and Louis Nicolas Vauquelin also observed iridium in the black platinum residue in 1803, but did not obtain enough material for further experiments.[34] Later the two French chemists Antoine-François Fourcroy and Nicolas-Louis Vauquelin identified a metal in a platinum residue they called ‘ptène’.[35]
|
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+
|
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In 1803, Smithson Tennant analyzed the insoluble residue and concluded that it must contain a new metal. Vauquelin treated the powder alternately with alkali and acids[36] and obtained a volatile new oxide, which he believed was of this new metal—which he named ptene, from the Greek word πτηνος (ptènos) for winged.[37][38] However, Tennant, who had the advantage of a much larger amount of residue, continued his research and identified two previously undiscovered elements in the black residue, iridium and osmium.[34][36] He obtained a yellow solution (probably of cis–[Os(OH)2O4]2−) by reactions with sodium hydroxide at red heat. After acidification he was able to distill the formed OsO4.[37] He named it osmium after Greek osme meaning "a smell", because of the ashy and smoky smell of the volatile osmium tetroxide.[39] Discovery of the new elements was documented in a letter to the Royal Society on June 21, 1804.[34][40]
|
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|
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Uranium and osmium were early successful catalysts in the Haber process, the nitrogen fixation reaction of nitrogen and hydrogen to produce ammonia, giving enough yield to make the process economically successful. At the time, a group at BASF led by Carl Bosch bought most of the world's supply of osmium to use as a catalyst. Shortly thereafter, in 1908, cheaper catalysts based on iron and iron oxides were introduced by the same group for the first pilot plants, removing the need for the expensive and rare osmium.[41]
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Nowadays osmium is obtained primarily from the processing of platinum and nickel ores.[42]
|
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Osmium is one of the even-numbered elements, which puts it in the upper half of elements commonly found in space. It is, however, the least abundant stable element in Earth's crust, with an average mass fraction of 50 parts per trillion in the continental crust.[43]
|
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|
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Osmium is found in nature as an uncombined element or in natural alloys; especially the iridium–osmium alloys, osmiridium (osmium rich), and iridosmium (iridium rich).[36] In nickel and copper deposits, the platinum group metals occur as sulfides (i.e., (Pt,Pd)S)), tellurides (e.g., PtBiTe), antimonides (e.g., PdSb), and arsenides (e.g., PtAs2); in all these compounds platinum is exchanged by a small amount of iridium and osmium. As with all of the platinum group metals, osmium can be found naturally in alloys with nickel or copper.[44]
|
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Within Earth's crust, osmium, like iridium, is found at highest concentrations in three types of geologic structure: igneous deposits (crustal intrusions from below), impact craters, and deposits reworked from one of the former structures. The largest known primary reserves are in the Bushveld Igneous Complex in South Africa,[45] though the large copper–nickel deposits near Norilsk in Russia, and the Sudbury Basin in Canada are also significant sources of osmium. Smaller reserves can be found in the United States.[45] The alluvial deposits used by pre-Columbian people in the Chocó Department, Colombia are still a source for platinum group metals. The second large alluvial deposit was found in the Ural Mountains, Russia, which is still mined.[42][46]
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Osmium is obtained commercially as a by-product from nickel and copper mining and processing. During electrorefining of copper and nickel, noble metals such as silver, gold and the platinum group metals, together with non-metallic elements such as selenium and tellurium settle to the bottom of the cell as anode mud, which forms the starting material for their extraction.[47][48] Separating the metals requires that they first be brought into solution. Several methods can achieve this, depending on the separation process and the composition of the mixture. Two representative methods are fusion with sodium peroxide followed by dissolution in aqua regia, and dissolution in a mixture of chlorine with hydrochloric acid.[45][49] Osmium, ruthenium, rhodium and iridium can be separated from platinum, gold and base metals by their insolubility in aqua regia, leaving a solid residue. Rhodium can be separated from the residue by treatment with molten sodium bisulfate. The insoluble residue, containing Ru, Os and Ir, is treated with sodium oxide, in which Ir is insoluble, producing water-soluble Ru and Os salts. After oxidation to the volatile oxides, RuO4 is separated from OsO4 by precipitation of (NH4)3RuCl6 with ammonium chloride.
|
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After it is dissolved, osmium is separated from the other platinum group metals by distillation or extraction with organic solvents of the volatile osmium tetroxide.[50] The first method is similar to the procedure used by Tennant and Wollaston. Both methods are suitable for industrial scale production. In either case, the product is reduced using hydrogen, yielding the metal as a powder or sponge that can be treated using powder metallurgy techniques.[51]
|
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Neither the producers nor the United States Geological Survey published any production amounts for osmium. In 1971, estimations of the United States production of osmium as a byproduct of copper refining was 2000 troy ounces (62 kg).[52] In 2017, the estimated US import of osmium for consumption was 90 kg.[53]
|
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|
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Because of the volatility and extreme toxicity of its oxide, osmium is rarely used in its pure state, but is instead often alloyed with other metals for high-wear applications. Osmium alloys such as osmiridium are very hard and, along with other platinum-group metals, are used in the tips of fountain pens, instrument pivots, and electrical contacts, as they can resist wear from frequent operation. They were also used for the tips of phonograph styli during the late 78 rpm and early "LP" and "45" record era, circa 1945 to 1955. Osmium-alloy tips were significantly more durable than steel and chromium needle points, but wore out far more rapidly than competing, and costlier, sapphire and diamond tips, so they were discontinued.[54]
|
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|
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Osmium tetroxide has been used in fingerprint detection[55] and in staining fatty tissue for optical and electron microscopy. As a strong oxidant, it cross-links lipids mainly by reacting with unsaturated carbon–carbon bonds and thereby both fixes biological membranes in place in tissue samples and simultaneously stains them. Because osmium atoms are extremely electron-dense, osmium staining greatly enhances image contrast in transmission electron microscopy (TEM) studies of biological materials. Those carbon materials otherwise have very weak TEM contrast (see image).[23] Another osmium compound, osmium ferricyanide (OsFeCN), exhibits similar fixing and staining action.[56]
|
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The tetroxide and its derivative potassium osmate are important oxidants in organic synthesis. For the Sharpless asymmetric dihydroxylation, which uses osmate for the conversion of a double bond into a vicinal diol, Karl Barry Sharpless was awarded the Nobel Prize in Chemistry in 2001.[57][58] OsO4 is very expensive for this use, so KMnO4 is often used instead, even though the yields are less for this cheaper chemical reagent.
|
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|
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In 1898 an Austrian chemist Auer von Welsbach developed the Oslamp with a filament made of osmium, which he introduced commercially in 1902. After only a few years, osmium was replaced by the more stable metal tungsten. Tungsten has the highest melting point among all metals, and its use in light bulbs increases the luminous efficacy and life of incandescent lamps.[37]
|
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|
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The light bulb manufacturer Osram (founded in 1906, when three German companies, Auer-Gesellschaft, AEG and Siemens & Halske, combined their lamp production facilities) derived its name from the elements of osmium and Wolfram (the latter is German for tungsten).[59]
|
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|
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Like palladium, powdered osmium effectively absorbs hydrogen atoms. This could make osmium a potential candidate for a metal-hydride battery electrode. However, osmium is expensive and would react with potassium hydroxide, the most common battery electrolyte.[60]
|
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|
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Osmium has high reflectivity in the ultraviolet range of the electromagnetic spectrum; for example, at 600 Å osmium has a reflectivity twice that of gold.[61] This high reflectivity is desirable in space-based UV spectrometers, which have reduced mirror sizes due to space limitations. Osmium-coated mirrors were flown in several space missions aboard the Space Shuttle, but it soon became clear that the oxygen radicals in the low Earth orbit are abundant enough to significantly deteriorate the osmium layer.[62]
|
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|
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The only known clinical use of osmium is synovectomy in arthritic patients in Scandinavia.[63] It involves the local administration of osmium tetroxide (OsO4), which is a highly toxic compound. The lack of reports of long-term side effects suggest that osmium itself can be biocompatible, though this depends on the osmium compound administered. In 2011, osmium(VI)[64] and osmium(II)[65] compounds were reported to show anticancer activity in vivo, it indicated a promising future for using osmium compounds as anticancer drugs.[66]
|
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|
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The Sharpless dihydroxylation: RL = largest substituent; RM = medium-sized substituent; RS = smallest substituent
|
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|
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Post-flight appearance of Os, Ag, and Au mirrors from the front (left images) and rear panels of the Space Shuttle. Blackening reveals oxidation due to irradiation by oxygen atoms.[67][68]
|
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|
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Metallic osmium is harmless[69] but finely divided metallic osmium is pyrophoric[52] and reacts with oxygen at room temperature, forming volatile osmium tetroxide. Some osmium compounds are also converted to the tetroxide if oxygen is present.[52] This makes osmium tetroxide the main source of contact with the environment.
|
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|
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Osmium tetroxide is highly volatile and penetrates skin readily, and is very toxic by inhalation, ingestion, and skin contact.[70] Airborne low concentrations of osmium tetroxide vapor can cause lung congestion and skin or eye damage, and should therefore be used in a fume hood.[19] Osmium tetroxide is rapidly reduced to relatively inert compounds by e.g. ascorbic acid[71] or polyunsaturated vegetable oils (such as corn oil).[72]
|
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|
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Osmium is usually sold as a minimum 99.9% pure powder. Like other precious metals, it is measured by troy weight and by grams. The market price of osmium has not changed in decades, primarily because little change has occurred in supply and demand. In addition to so little of it being available, osmium is difficult to work with, has few uses, and is a challenge to store safely because of the toxic compound it produces when it oxidizes.
|
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|
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While the price of $400 per troy ounce has remained steady since the 1990s, inflation since that time has led to the metal losing about one-third of its value in the two decades prior to 2019.
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