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A mobile phone, cellular phone, cell phone, cellphone or hand phone, sometimes shortened to simply mobile, cell or just phone, is a portable telephone that can make and receive calls over a radio frequency link while the user is moving within a telephone service area. The radio frequency link establishes a connection to the switching systems of a mobile phone operator, which provides access to the public switched telephone network (PSTN). Modern mobile telephone services use a cellular network architecture and, therefore, mobile telephones are called cellular telephones or cell phones in North America. In addition to telephony, digital mobile phones (2G) support a variety of other services, such as text messaging, MMS, email, Internet access, short-range wireless communications (infrared, Bluetooth), business applications, video games and digital photography. Mobile phones offering only those capabilities are known as feature phones; mobile phones which offer greatly advanced computing capabilities are referred to as smartphones.
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The development of metal-oxide-semiconductor (MOS) large-scale integration (LSI) technology, information theory and cellular networking led to the development of affordable mobile communications.[1] The first handheld mobile phone was demonstrated by John F. Mitchell[2][3] and Martin Cooper of Motorola in 1973, using a handset weighing c. 2 kilograms (4.4 lbs).[4] In 1979, Nippon Telegraph and Telephone (NTT) launched the world's first cellular network in Japan.[citation needed] In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone. From 1983 to 2014, worldwide mobile phone subscriptions grew to over seven billion; enough to provide one for every person on Earth.[5] In the first quarter of 2016, the top smartphone developers worldwide were Samsung, Apple and Huawei; smartphone sales represented 78 percent of total mobile phone sales.[6] For feature phones (slang: “dumbphones”) as of 2016[update], the largest were Samsung, Nokia and Alcatel.[7]
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A handheld mobile radio telephone service was envisioned in the early stages of radio engineering. In 1917, Finnish inventor Eric Tigerstedt filed a patent for a "pocket-size folding telephone with a very thin carbon microphone". Early predecessors of cellular phones included analog radio communications from ships and trains. The race to create truly portable telephone devices began after World War II, with developments taking place in many countries. The advances in mobile telephony have been traced in successive "generations", starting with the early zeroth-generation (0G) services, such as Bell System's Mobile Telephone Service and its successor, the Improved Mobile Telephone Service. These 0G systems were not cellular, supported few simultaneous calls, and were very expensive.
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The development of metal-oxide-semiconductor (MOS) large-scale integration (LSI) technology, information theory and cellular networking led to the development of affordable mobile communications.[1] The first handheld cellular mobile phone was demonstrated by John F. Mitchell[2][3] and Martin Cooper of Motorola in 1973, using a handset weighing 2 kilograms (4.4 lb).[4] The first commercial automated cellular network (1G) analog was launched in Japan by Nippon Telegraph and Telephone in 1979. This was followed in 1981 by the simultaneous launch of the Nordic Mobile Telephone (NMT) system in Denmark, Finland, Norway, and Sweden.[8] Several other countries then followed in the early to mid-1980s. These first-generation (1G) systems could support far more simultaneous calls but still used analog cellular technology. In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone.
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Digital cellular networks appeared in the 1990s, enabled by the wide adoption of MOSFET-based RF power amplifiers (power MOSFET and LDMOS) and RF circuits (RF CMOS),[9][10][11] leading to the introduction of digital signal processing in wireless communications.[1] In 1991, the second-generation (2G) digital cellular technology was launched in Finland by Radiolinja on the GSM standard. This sparked competition in the sector as the new operators challenged the incumbent 1G network operators. The GSM standard is a European initiative expressed at the CEPT ("Conférence Européenne des Postes et Telecommunications", European Postal and Telecommunications conference). The Franco-German R&D cooperation demonstrated the technical feasibility, and in 1987 a Memorandum of Understanding was signed between 13 European countries who agreed to launch a commercial service by 1991. The first version of the GSM (=2G) standard had 6,000 pages. The IEEE/RSE awarded to Thomas Haug and Philippe Dupuis [fr] the 2018 James Clerk Maxwell medal for their contributions to the first digital mobile telephone standard.[12] In 2018, the GSM was used by over 5 billion people in over 220 countries. The GSM (2G) has evolved into 3G, 4G and 5G.
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The lithium-ion battery, an indispensable energy source for modern mobile phones,[13] was commercialized by Sony and Asahi Kasei in 1991.[14][15] In 2001, the third generation (3G) was launched in Japan by NTT DoCoMo on the WCDMA standard.[16] This was followed by 3.5G, 3G+ or turbo 3G enhancements based on the high-speed packet access (HSPA) family, allowing UMTS networks to have higher data transfer speeds and capacity.
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By 2009, it had become clear that, at some point, 3G networks would be overwhelmed by the growth of bandwidth-intensive applications, such as streaming media.[17] Consequently, the industry began looking to data-optimized fourth-generation technologies, with the promise of speed improvements up to ten-fold over existing 3G technologies. The first two commercially available technologies billed as 4G were the WiMAX standard, offered in North America by Sprint, and the LTE standard, first offered in Scandinavia by TeliaSonera.
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5G is a technology and term used in research papers and projects to denote the next major phase in mobile telecommunication standards beyond the 4G/IMT-Advanced standards. The term 5G is not officially used in any specification or official document yet made public by telecommunication companies or standardization bodies such as 3GPP, WiMAX Forum or ITU-R. New standards beyond 4G are currently being developed by standardization bodies, but they are at this time seen as under the 4G umbrella, not for a new mobile generation.
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Smartphones have a number of distinguishing features. The International Telecommunication Union measures those with Internet connection, which it calls Active Mobile-Broadband subscriptions (which includes tablets, etc.). In the developed world, smartphones have now overtaken the usage of earlier mobile systems. However, in the developing world, they account for around 50% of mobile telephony.
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Feature phone is a term typically used as a retronym to describe mobile phones which are limited in capabilities in contrast to a modern smartphone. Feature phones typically provide voice calling and text messaging functionality, in addition to basic multimedia and Internet capabilities, and other services offered by the user's wireless service provider. A feature phone has additional functions over and above a basic mobile phone which is only capable of voice calling and text messaging.[20][21] Feature phones and basic mobile phones tend to use a proprietary, custom-designed software and user interface. By contrast, smartphones generally use a mobile operating system that often shares common traits across devices.
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There are Orthodox Jewish religious restrictions which, by some interpretations, standard mobile telephones overstep. To deal with this problem, some rabbinical organizations have recommended that phones with text-messaging capability not be used by children.[22] Phones with restricted features are known as kosher phones and have rabbinical approval for use in Israel and elsewhere by observant Orthodox Jews. Although these phones are intended to prevent immodesty, some vendors report good sales to adults who prefer the simplicity of the devices. Some phones are approved for use by essential workers (such as health, security, and public service workers) on the sabbath (the use of any electrical device is generally prohibited during this time, other than to save lives, or reduce the risk of death or similar needs).[23]
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Mobile phones communicate with cell towers that are placed to give coverage across a telephone service area which is divided up into 'cells'. Each cell uses a different set of frequencies from neighboring cells, and will typically be covered by 3 towers placed at different locations. The cell towers are usually interconnected to each other and the phone network and the internet by wired connections. Due to bandwidth limitations each cell will have a maximum number of cell phones it can handle at once. The cells are therefore sized depending on the expected usage density, and may be much smaller in cities. In that case much lower transmitter powers are used to avoid broadcasting beyond the cell.
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In order to handle the high traffic, multiple towers can be setup in the same area(using different frequencies). This can be done permanently or temporarily such as at special events like at the Super Bowl, Taste of Chicago, State Fair, NYC New Year's Eve, hurricane hit cities, etc. where cell phone companies will bring a truck with equipment to host the abnormally high traffic with a portable cell.
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Cellular can greatly increase the capacity of simultaneous wireless phone calls. While a phone company for example, has a license to 1000 frequencies, each cell must use unique frequencies with each call using one of them when communicating. Because cells only slightly overlap, the same frequency can be reused. Example cell 1 uses frequency 1–500, next door cell uses frequency 501–1000, next door can reuse frequency 1–500. Cells 1 and 3 are not "touching" and do not overlap/communicate so each can reuse the same frequencies.[citation needed]
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This is even more greatly increased when phone companies implemented digital networks. With digital, one frequency can host multiple simultaneous calls increasing capacity even more.
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As a phone moves around, a phone will "hand off" - automatically disconnect and reconnect to the tower of another cell that gives the best reception.
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Additionally, short-range Wi-Fi infrastructure is often used by smartphones as much as possible as it offloads traffic from cell networks on to local area networks.
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The common components found on all mobile phones are:
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Low-end mobile phones are often referred to as feature phones and offer basic telephony. Handsets with more advanced computing ability through the use of native software applications are known as smartphones.
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Mobile phones have central processing units (CPUs), similar to those in computers, but optimised to operate in low power environments.
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Mobile CPU performance depends not only on the clock rate (generally given in multiples of hertz)[24] but also the memory hierarchy also greatly affects overall performance. Because of these problems, the performance of mobile phone CPUs is often more appropriately given by scores derived from various standardized tests to measure the real effective performance in commonly used applications.
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One of the main characteristics of phones is the screen. Depending on the device's type and design, the screen fills most or nearly all of the space on a device's front surface. Many smartphone displays have an aspect ratio of 16:9, but taller aspect ratios became more common in 2017.
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Screen sizes are often measured in diagonal inches or millimeters; feature phones generally have screen sizes below 90 millimetres (3.5 in). Phones with screens larger than 130 millimetres (5.2 in) are often called "phablets." Smartphones with screens over 115 millimetres (4.5 in) in size are commonly difficult to use with only a single hand, since most thumbs cannot reach the entire screen surface; they may need to be shifted around in the hand, held in one hand and manipulated by the other, or used in place with both hands. Due to design advances, some modern smartphones with large screen sizes and "edge-to-edge" designs have compact builds that improve their ergonomics, while the shift to taller aspect ratios have resulted in phones that have larger screen sizes whilst maintaining the ergonomics associated with smaller 16:9 displays.[25][26][27]
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Liquid-crystal displays are the most common; others are IPS, LED, OLED, and AMOLED displays. Some displays are integrated with pressure-sensitive digitizers, such as those developed by Wacom and Samsung,[28] and Apple's "3D Touch" system.
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In sound, smartphones and feature phones vary little. Some audio-quality enhancing features, such as Voice over LTE and HD Voice, have appeared and are often available on newer smartphones. Sound quality can remain a problem due to the design of the phone, the quality of the cellular network and compression algorithms used in long-distance calls.[29][30] Audio quality can be improved using a VoIP application over WiFi.[31] Cellphones have small speakers so that the user can use a speakerphone feature and talk to a person on the phone without holding it to their ear. The small speakers can also be used to listen to digital audio files of music or speech or watch videos with an audio component, without holding the phone close to the ear.
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The average phone battery lasts 2–3 years at best. Many of the wireless devices use a Lithium-Ion (Li-Ion) battery, which charges 500–2500 times, depending on how users take care of the battery and the charging techniques used.[32] It is only natural for these rechargeable batteries to chemically age, which is why the performance of the battery when used for a year or two will begin to deteriorate. Battery life can be extended by draining it regularly, not overcharging it, and keeping it away from heat.[33][34]
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Mobile phones require a small microchip called a Subscriber Identity Module or SIM card, in order to function. The SIM card is approximately the size of a small postage stamp and is usually placed underneath the battery in the rear of the unit. The SIM securely stores the service-subscriber key (IMSI) and the Ki used to identify and authenticate the user of the mobile phone. The SIM card allows users to change phones by simply removing the SIM card from one mobile phone and inserting it into another mobile phone or broadband telephony device, provided that this is not prevented by a SIM lock. The first SIM card was made in 1991 by Munich smart card maker Giesecke & Devrient for the Finnish wireless network operator Radiolinja.[citation needed]
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A hybrid mobile phone can hold up to four SIM cards, with a phone having an IMEI per SIM Card. SIM and R-UIM cards may be mixed together to allow both GSM and CDMA networks to be accessed. From 2010 onwards, such phones became popular in emerging markets,[35] and this was attributed to the desire to obtain the lowest on-net calling rate.
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Feature phones have basic software platforms.
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Smartphones have advanced software platforms.
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A mobile app is a computer program designed to run on a mobile device, such as a smartphone. The term "app" is a shortening of the term "software application".
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A common data application on mobile phones is Short Message Service (SMS) text messaging. The first SMS message was sent from a computer to a mobile phone in 1992 in the UK while the first person-to-person SMS from phone to phone was sent in Finland in 1993. The first mobile news service, delivered via SMS, was launched in Finland in 2000,[citation needed] and subsequently many organizations provided "on-demand" and "instant" news services by SMS. Multimedia Messaging Service (MMS) was introduced in 2001.[citation needed]
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The introduction of Apple's App Store for the iPhone and iPod Touch in July 2008 popularized manufacturer-hosted online distribution for third-party applications (software and computer programs) focused on a single platform. There are a huge variety of apps, including video games, music products and business tools. Up until that point, smartphone application distribution depended on third-party sources providing applications for multiple platforms, such as GetJar, Handango, Handmark, and PocketGear. Following the success of the App Store, other smartphone manufacturers launched application stores, such as Google's Android Market (later renamed to the Google Play Store) and RIM's BlackBerry App World and Android-related app stores like F-Droid. In February 2014, 93% of mobile developers were targeting smartphones first for mobile app development.[36]
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From 1983 to 1998, Motorola was market leader in mobile phones. Nokia was the market leader in mobile phones from 1998 to 2012.[38] In Q1 2012, Samsung surpassed Nokia, selling 93.5 million units as against Nokia's 82.7 million units. Samsung has retained its top position since then. In 2017, the top five manufacturers worldwide were Samsung (20.9%), Apple (14.0%), Huawei (9.8%), Oppo (5.7%), and Vivo (6.5%).[39] During Q2 2018, Huawei overtook Apple as the world's second-largest phone manufacturer.[40]
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The world's largest individual mobile operator by number of subscribers is China Mobile, which has over 902 million mobile phone subscribers as of June 2018[update].[41] Over 50 mobile operators have over ten million subscribers each, and over 150 mobile operators had at least one million subscribers by the end of 2009.[42] In 2014, there were more than seven billion mobile phone subscribers worldwide, a number that is expected to keep growing.
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Mobile phones are used for a variety of purposes, such as keeping in touch with family members, for conducting business, and in order to have access to a telephone in the event of an emergency. Some people carry more than one mobile phone for different purposes, such as for business and personal use. Multiple SIM cards may be used to take advantage of the benefits of different calling plans. For example, a particular plan might provide for cheaper local calls, long-distance calls, international calls, or roaming.
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The mobile phone has been used in a variety of diverse contexts in society. For example:
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In 1998, one of the first examples of distributing and selling media content through the mobile phone was the sale of ringtones by Radiolinja in Finland. Soon afterwards, other media content appeared, such as news, video games, jokes, horoscopes, TV content and advertising. Most early content for mobile phones tended to be copies of legacy media, such as banner advertisements or TV news highlight video clips. Recently, unique content for mobile phones has been emerging, from ringtones and ringback tones to mobisodes, video content that has been produced exclusively for mobile phones.
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In many countries, mobile phones are used to provide mobile banking services, which may include the ability to transfer cash payments by secure SMS text message. Kenya's M-PESA mobile banking service, for example, allows customers of the mobile phone operator Safaricom to hold cash balances which are recorded on their SIM cards. Cash can be deposited or withdrawn from M-PESA accounts at Safaricom retail outlets located throughout the country and can be transferred electronically from person to person and used to pay bills to companies.
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Branchless banking has also been successful in South Africa and the Philippines. A pilot project in Bali was launched in 2011 by the International Finance Corporation and an Indonesian bank, Bank Mandiri.[51]
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Another application of mobile banking technology is Zidisha, a US-based nonprofit micro-lending platform that allows residents of developing countries to raise small business loans from Web users worldwide. Zidisha uses mobile banking for loan disbursements and repayments, transferring funds from lenders in the United States to borrowers in rural Africa who have mobile phones and can use the Internet.[52]
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Mobile payments were first trialled in Finland in 1998 when two Coca-Cola vending machines in Espoo were enabled to work with SMS payments. Eventually, the idea spread and in 1999, the Philippines launched the country's first commercial mobile payments systems with mobile operators Globe and Smart.
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Some mobile phones can make mobile payments via direct mobile billing schemes, or through contactless payments if the phone and the point of sale support near field communication (NFC).[53] Enabling contactless payments through NFC-equipped mobile phones requires the co-operation of manufacturers, network operators, and retail merchants.[54][55]
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Mobile phones are commonly used to collect location data. While the phone is turned on, the geographical location of a mobile phone can be determined easily (whether it is being used or not) using a technique known as multilateration to calculate the differences in time for a signal to travel from the mobile phone to each of several cell towers near the owner of the phone.[56][57]
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The movements of a mobile phone user can be tracked by their service provider and if desired, by law enforcement agencies and their governments. Both the SIM card and the handset can be tracked.[56]
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China has proposed using this technology to track the commuting patterns of Beijing city residents.[58] In the UK and US, law enforcement and intelligence services use mobile phones to perform surveillance operations. They possess technology that enables them to activate the microphones in mobile phones remotely in order to listen to conversations which take place near the phone.[59][60]
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Hackers are able to track a phone's location, read messages, and record calls, just by knowing the phone number.[61]
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Mobile phone use while driving, including talking on the phone, texting, or operating other phone features, is common but controversial. It is widely considered dangerous due to distracted driving. Being distracted while operating a motor vehicle has been shown to increase the risk of accidents. In September 2010, the US National Highway Traffic Safety Administration (NHTSA) reported that 995 people were killed by drivers distracted by cell phones. In March 2011, a U.S. insurance company, State Farm Insurance, announced the results of a study which showed 19% of drivers surveyed accessed the Internet on a smartphone while driving.[62] Many jurisdictions prohibit the use of mobile phones while driving. In Egypt, Israel, Japan, Portugal, and Singapore, both handheld and hands-free use of a mobile phone (which uses a speakerphone) is banned. In other countries, including the UK and France and in many U.S. states, only handheld phone use is banned while hands-free use is permitted.
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A 2011 study reported that over 90% of college students surveyed text (initiate, reply or read) while driving.[63]
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The scientific literature on the dangers of driving while sending a text message from a mobile phone, or texting while driving, is limited. A simulation study at the University of Utah found a sixfold increase in distraction-related accidents when texting.[64]
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Due to the increasing complexity of mobile phones, they are often more like mobile computers in their available uses. This has introduced additional difficulties for law enforcement officials when attempting to distinguish one usage from another in drivers using their devices. This is more apparent in countries which ban both handheld and hands-free usage, rather than those which ban handheld use only, as officials cannot easily tell which function of the mobile phone is being used simply by looking at the driver. This can lead to drivers being stopped for using their device illegally for a phone call when, in fact, they were using the device legally, for example, when using the phone's incorporated controls for car stereo, GPS or satnav.
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A 2010 study reviewed the incidence of mobile phone use while cycling and its effects on behaviour and safety.[65] In 2013, a national survey in the US reported the number of drivers who reported using their cellphones to access the Internet while driving had risen to nearly one of four.[66] A study conducted by the University of Vienna examined approaches for reducing inappropriate and problematic use of mobile phones, such as using mobile phones while driving.[67]
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Accidents involving a driver being distracted by talking on a mobile phone have begun to be prosecuted as negligence similar to speeding. In the United Kingdom, from 27 February 2007, motorists who are caught using a hand-held mobile phone while driving will have three penalty points added to their license in addition to the fine of £60.[68] This increase was introduced to try to stem the increase in drivers ignoring the law.[69] Japan prohibits all mobile phone use while driving, including use of hands-free devices. New Zealand has banned hand-held cell phone use since 1 November 2009. Many states in the United States have banned texting on cell phones while driving. Illinois became the 17th American state to enforce this law.[70] As of July 2010[update], 30 states had banned texting while driving, with Kentucky becoming the most recent addition on 15 July.[71]
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Public Health Law Research maintains a list of distracted driving laws in the United States. This database of laws provides a comprehensive view of the provisions of laws that restrict the use of mobile communication devices while driving for all 50 states and the District of Columbia between 1992 when first law was passed, through 1 December 2010. The dataset contains information on 22 dichotomous, continuous or categorical variables including, for example, activities regulated (e.g., texting versus talking, hands-free versus handheld), targeted populations, and exemptions.[72]
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In 2010, an estimated 1500 pedestrians were injured in the US while using a cellphone and some jurisdictions have attempted to ban pedestrians from using their cellphones.[73][74]
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The effect of mobile phone radiation on human health is the subject of recent[when?] interest and study, as a result of the enormous increase in mobile phone usage throughout the world. Mobile phones use electromagnetic radiation in the microwave range, which some believe may be harmful to human health. A large body of research exists, both epidemiological and experimental, in non-human animals and in humans. The majority of this research shows no definite causative relationship between exposure to mobile phones and harmful biological effects in humans. This is often paraphrased simply as the balance of evidence showing no harm to humans from mobile phones, although a significant number of individual studies do suggest such a relationship, or are inconclusive. Other digital wireless systems, such as data communication networks, produce similar radiation.[citation needed]
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+
On 31 May 2011, the World Health Organization stated that mobile phone use may possibly represent a long-term health risk,[75][76] classifying mobile phone radiation as "possibly carcinogenic to humans" after a team of scientists reviewed studies on mobile phone safety.[77] The mobile phone is in category 2B, which ranks it alongside coffee and other possibly carcinogenic substances.[78][79]
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Some recent[when?] studies have found an association between mobile phone use and certain kinds of brain and salivary gland tumors. Lennart Hardell and other authors of a 2009 meta-analysis of 11 studies from peer-reviewed journals concluded that cell phone usage for at least ten years "approximately doubles the risk of being diagnosed with a brain tumor on the same ('ipsilateral') side of the head as that preferred for cell phone use".[80]
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One study of past mobile phone use cited in the report showed a "40% increased risk for gliomas (brain cancer) in the highest category of heavy users (reported average: 30 minutes per day over a 10‐year period)".[81] This is a reversal of the study's prior position that cancer was unlikely to be caused by cellular phones or their base stations and that reviews had found no convincing evidence for other health effects.[76][82] However, a study published 24 March 2012, in the British Medical Journal questioned these estimates because the increase in brain cancers has not paralleled the increase in mobile phone use.[83] Certain countries, including France, have warned against the use of mobile phones by minors in particular, due to health risk uncertainties.[84] Mobile pollution by transmitting electromagnetic waves can be decreased up to 90% by adopting the circuit as designed in mobile phone and mobile exchange.[85]
|
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+
In May 2016, preliminary findings of a long-term study by the U.S. government suggested that radio-frequency (RF) radiation, the type emitted by cell phones, can cause cancer.[86][87]
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+
A study by the London School of Economics found that banning mobile phones in schools could increase pupils' academic performance, providing benefits equal to one extra week of schooling per year.[88]
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Studies have shown that around 40–50% of the environmental impact of mobile phones occurs during the manufacture of their printed wiring boards and integrated circuits.[89]
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+
The average user replaces their mobile phone every 11 to 18 months,[90] and the discarded phones then contribute to electronic waste. Mobile phone manufacturers within Europe are subject to the WEEE directive, and Australia has introduced a mobile phone recycling scheme.[91]
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+
Apple Inc. had an advanced robotic disassembler and sorter called Liam specifically for recycling outdated or broken iPhones.[350]
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According to the Federal Communications Commission, one out of three robberies involve the theft of a cellular phone.[citation needed] Police data in San Francisco show that half of all robberies in 2012 were thefts of cellular phones.[citation needed] An online petition on Change.org, called Secure our Smartphones, urged smartphone manufacturers to install kill switches in their devices to make them unusable if stolen. The petition is part of a joint effort by New York Attorney General Eric Schneiderman and San Francisco District Attorney George Gascón and was directed to the CEOs of the major smartphone manufacturers and telecommunication carriers.[92] On 10 June 2013, Apple announced that it would install a "kill switch" on its next iPhone operating system, due to debut in October 2013.[93]
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All mobile phones have a unique identifier called IMEI. Anyone can report their phone as lost or stolen with their Telecom Carrier, and the IMEI would be blacklisted with a central registry.[94] Telecom carriers, depending upon local regulation can or must implement blocking of blacklisted phones in their network. There are, however, a number of ways to circumvent a blacklist. One method is to send the phone to a country where the telecom carriers are not required to implement the blacklisting and sell it there,[95] another involves altering the phone's IMEI number.[96] Even so, mobile phones typically have less value on the second-hand market if the phones original IMEI is blacklisted.
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An unusual example of a phone bill caused by theft (reported on 28 June 2018) was when a biological group in Poland put a GPS tracker on a white stork and released it; during autumn migration over the Blue Nile valley in eastern Sudan someone got hold of the stork's GPS tracker, and found in it a mobile-phone-type sim card, which he put in his mobile phone, and made 20 hours of calls on it, running up a bill of over 10,000 Polish zlotys (US$2,700) for the biological group.[97]
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Demand for metals used in mobile phones and other electronics fuelled the Second Congo War, which claimed almost 5.5 million lives.[98] In a 2012 news story, The Guardian reported: "In unsafe mines deep underground in eastern Congo, children are working to extract minerals essential for the electronics industry. The profits from the minerals finance the bloodiest conflict since the second world war; the war has lasted nearly 20 years and has recently flared up again. ... For the last 15 years, the Democratic Republic of the Congo has been a major source of natural resources for the mobile phone industry."[99] The company Fairphone has worked to develop a mobile phone that does not contain conflict minerals.
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A telephone is a telecommunications device that permits two or more users to conduct a conversation when they are too far apart to be heard directly. A telephone converts sound, typically and most efficiently the human voice, into electronic signals that are transmitted via cables and other communication channels to another telephone which reproduces the sound to the receiving user. The term is derived from Greek: τῆλε (tēle, far) and φωνή (phōnē, voice), together meaning distant voice. A common short form of the term is phone, which has been in use since the early 20th century.
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In 1876, Alexander Graham Bell was the first to be granted a United States patent for a device that produced clearly intelligible replication of the human voice. This instrument was further developed by many others, and became rapidly indispensable in business, government, and in households.
|
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|
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The essential elements of a telephone are a microphone (transmitter) to speak into and an earphone (receiver) which reproduces the voice in a distant location.[1] In addition, most telephones contain a ringer to announce an incoming telephone call, and a dial or keypad to enter a telephone number when initiating a call to another telephone. The receiver and transmitter are usually built into a handset which is held up to the ear and mouth during conversation. The dial may be located either on the handset or on a base unit to which the handset is connected. The transmitter converts the sound waves to electrical signals which are sent through a telephone network to the receiving telephone, which converts the signals into audible sound in the receiver or sometimes a loudspeaker. Telephones are duplex devices, meaning they permit transmission in both directions simultaneously.
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The first telephones were directly connected to each other from one customer's office or residence to another customer's location. Being impractical beyond just a few customers, these systems were quickly replaced by manually operated centrally located switchboards. These exchanges were soon connected together, eventually forming an automated, worldwide public switched telephone network. For greater mobility, various radio systems were developed for transmission between mobile stations on ships and automobiles in the mid-20th century. Hand-held mobile phones were introduced for personal service starting in 1973. In later decades their analog cellular system evolved into digital networks with greater capability and lower cost.
|
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+
Convergence has given most modern cell phones capabilities far beyond simple voice conversation. Most are smartphones, integrating all mobile communication and many computing needs.
|
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+
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+
A traditional landline telephone system, also known as plain old telephone service (POTS), commonly carries both control and audio signals on the same twisted pair (C in diagram) of insulated wires, the telephone line. The control and signaling equipment consists of three components, the ringer, the hookswitch, and a dial. The ringer, or beeper, light or other device (A7), alerts the user to incoming calls. The hookswitch signals to the central office that the user has picked up the handset to either answer a call or initiate a call. A dial, if present, is used by the subscriber to transmit a telephone number to the central office when initiating a call. Until the 1960s dials used almost exclusively the rotary technology, which was replaced by dual-tone multi-frequency signaling (DTMF) with pushbutton telephones (A4).
|
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|
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+
A major expense of wire-line telephone service is the outside wire plant. Telephones transmit both the incoming and outgoing speech signals on a single pair of wires. A twisted pair line rejects electromagnetic interference (EMI) and crosstalk better than a single wire or an untwisted pair. The strong outgoing speech signal from the microphone (transmitter) does not overpower the weaker incoming speaker (receiver) signal with sidetone because a hybrid coil (A3) and other components compensate the imbalance. The junction box (B) arrests lightning (B2) and adjusts the line's resistance (B1) to maximize the signal power for the line length. Telephones have similar adjustments for inside line lengths (A8). The line voltages are negative compared to earth, to reduce galvanic corrosion. Negative voltage attracts positive metal ions toward the wires.
|
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|
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+
The landline telephone contains a switchhook (A4) and an alerting device, usually a ringer (A7), that remains connected to the phone line whenever the phone is "on hook" (i.e. the switch (A4) is open), and other components which are connected when the phone is "off hook". The off-hook components include a transmitter (microphone, A2), a receiver (speaker, A1), and other circuits for dialing, filtering (A3), and amplification.
|
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|
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+
The place a telephone call, the calling party picks up the telephone's handset, thereby operating a lever which closes the hook switch (A4). This powers the telephone by connecting the transmission hybrid transformer, as well as the transmitter (microphone) and receiver (speaker) to the line. In this off-hook state, the telephone circuitry has a low resistance of typically than 300 ohms, which causes the flow of direct current (DC) in the line (C) from the telephone exchange. The exchange detects this current, attaches a digit receiver circuit to the line, and sends dial tone to indicate its readiness. On a modern push-button telephone, the caller then presses the number keys to send the telephone number of the destination, the called party. The keys control a tone generator circuit (not shown) that sends DTMF tones to the exchange. A rotary-dial telephone uses pulse dialing, sending electrical pulses, that the exchange counts to decode each digit of the telephone number. If the called party's line is available, the terminating exchange applies an intermittent alternating current (AC) ringing signal of 40 to 90 volts to alert the called party of the incoming call. If the called party's line is in use, however, the exchange returns a busy signal to the calling party. If the called party's line is in use but subscribes to call waiting service, the exchange sends an intermittent audible tone to the called party to indicate another call.
|
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+
The electromechanical ringer of a telephone (A7) is connected to the line through a capacitor (A6), which blocks direct current and passes the alternating current of the ringing power. The telephone draws no current when it is on hook, while a DC voltage is continually applied to the line. Exchange circuitry (D2) can send an alternating current down the line to activate the ringer and announce an incoming call. In manual service exchange areas, before dial service was installed, telephones had hand-cranked magneto generators to generate a ringing voltage back to the exchange or any other telephone on the same line. When a landline telephone is inactive (on hook), the circuitry at the telephone exchange detects the absence of direct current to indicate that the line is not in use.[2] When a party initiates a call to this line, the exchange sends the ringing signal. When the called party picks up the handset, they actuate a double-circuit switchhook (not shown) which may simultaneously disconnects the alerting device and connects the audio circuitry to the line. This, in turn, draws direct current through the line, confirming that the called phone is now active. The exchange circuitry turns off the ring signal, and both telephones are now active and connected through the exchange. The parties may now converse as long as both phones remain off hook. When a party hangs up, placing the handset back on the cradle or hook, direct current ceases in that line, signaling the exchange to disconnect the call.
|
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+
|
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+
Calls to parties beyond the local exchange are carried over trunk lines which establish connections between exchanges. In modern telephone networks, fiber-optic cable and digital technology are often employed in such connections. Satellite technology may be used for communication over very long distances.
|
24 |
+
|
25 |
+
In most landline telephones, the transmitter and receiver (microphone and speaker) are located in the handset, although in a speakerphone these components may be located in the base or in a separate enclosure. Powered by the line, the microphone (A2) produces a modulated electric current which varies its frequency and amplitude in response to the sound waves arriving at its diaphragm. The resulting current is transmitted along the telephone line to the local exchange then on to the other phone (via the local exchange or via a larger network), where it passes through the coil of the receiver (A3). The varying current in the coil produces a corresponding movement of the receiver's diaphragm, reproducing the original sound waves present at the transmitter.
|
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+
|
27 |
+
Along with the microphone and speaker, additional circuitry is incorporated to prevent the incoming speaker signal and the outgoing microphone signal from interfering with each other. This is accomplished through a hybrid coil (A3). The incoming audio signal passes through a resistor (A8) and the primary winding of the coil (A3) which passes it to the speaker (A1). Since the current path A8 – A3 has a far lower impedance than the microphone (A2), virtually all of the incoming signal passes through it and bypasses the microphone.
|
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+
|
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+
At the same time the DC voltage across the line causes a DC current which is split between the resistor-coil (A8-A3) branch and the microphone-coil (A2-A3) branch. The DC current through the resistor-coil branch has no effect on the incoming audio signal. But the DC current passing through the microphone is turned into AC (in response to voice sounds) which then passes through only the upper branch of the coil's (A3) primary winding, which has far fewer turns than the lower primary winding. This causes a small portion of the microphone output to be fed back to the speaker, while the rest of the AC goes out through the phone line.
|
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+
|
31 |
+
A lineman's handset is a telephone designed for testing the telephone network, and may be attached directly to aerial lines and other infrastructure components.
|
32 |
+
|
33 |
+
Before the development of the electric telephone, the term "telephone" was applied to other inventions, and not all early researchers of the electrical device called it "telephone". Perhaps the earliest use of the word for a communications system was the telephon created by Gottfried Huth in 1796. Huth proposed an alternative to the optical telegraph of Claude Chappe in which the operators in the signalling towers would shout to each other by means of what he called "speaking tubes", but would now be called giant megaphones.[4] A communication device for sailing vessels called a "telephone" was invented by the captain John Taylor in 1844. This instrument used four air horns to communicate with vessels in foggy weather.[5][6]
|
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+
|
35 |
+
Johann Philipp Reis used the term in reference to his invention, commonly known as the Reis telephone, in c. 1860. His device appears to be the first device based on conversion of sound into electrical impulses. The term telephone was adopted into the vocabulary of many languages. It is derived from the Greek: τῆλε, tēle, "far" and φωνή, phōnē, "voice", together meaning "distant voice".
|
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|
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+
Credit for the invention of the electric telephone is frequently disputed. As with other influential inventions such as radio, television, the light bulb, and the computer, several inventors pioneered experimental work on voice transmission over a wire and improved on each other's ideas. New controversies over the issue still arise from time to time. Charles Bourseul, Antonio Meucci, Johann Philipp Reis, Alexander Graham Bell, and Elisha Gray, amongst others, have all been credited with the invention of the telephone.[7][2]
|
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+
|
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+
Alexander Graham Bell was the first to be awarded a patent for the electric telephone by the United States Patent and Trademark Office (USPTO) in March 1876.[8] Before Bell's patent, the telephone transmitted sound in a way that was similar to the telegraph. This method used vibrations and circuits to send electrical pulses, but was missing key features. Bell found that this method produced a sound through intermittent currents, but in order for the telephone to work a fluctuating current reproduced sounds the best. The fluctuating currents became the basis for the working telephone, creating Bell's patent.[9] That first patent by Bell was the master patent of the telephone, from which other patents for electric telephone devices and features flowed.[10] The Bell patents were forensically victorious and commercially decisive.
|
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+
|
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+
In 1876, shortly after Bell's patent application, Hungarian engineer Tivadar Puskás proposed the telephone switch, which allowed for the formation of telephone exchanges, and eventually networks.[11]
|
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+
|
43 |
+
In the United Kingdom the blower is used as a slang term for a telephone. The term came from navy slang for a speaking tube.[12]
|
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+
|
45 |
+
Early telephones were technically diverse. Some used a water microphone, some had a metal diaphragm that induced current in an electromagnet wound around a permanent magnet, and some were dynamic – their diaphragm vibrated a coil of wire in the field of a permanent magnet or the coil vibrated the diaphragm. The sound-powered dynamic variants survived in small numbers through the 20th century in military and maritime applications, where its ability to create its own electrical power was crucial. Most, however, used the Edison/Berliner carbon transmitter, which was much louder than the other kinds, even though it required an induction coil which was an impedance matching transformer to make it compatible with the impedance of the line. The Edison patents kept the Bell monopoly viable into the 20th century, by which time the network was more important than the instrument.
|
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+
|
47 |
+
Early telephones were locally powered, using either a dynamic transmitter or by the powering of a transmitter with a local battery. One of the jobs of outside plant personnel was to visit each telephone periodically to inspect the battery. During the 20th century, telephones powered from the telephone exchange over the same wires that carried the voice signals became common.
|
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+
|
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+
Early telephones used a single wire for the subscriber's line, with ground return used to complete the circuit (as used in telegraphs). The earliest dynamic telephones also had only one port opening for sound, with the user alternately listening and speaking (or rather, shouting) into the same hole. Sometimes the instruments were operated in pairs at each end, making conversation more convenient but also more expensive.
|
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+
|
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+
At first, the benefits of a telephone exchange were not exploited. Instead telephones were leased in pairs to a subscriber, who had to arrange for a telegraph contractor to construct a line between them, for example between a home and a shop. Users who wanted the ability to speak to several different locations would need to obtain and set up three or four pairs of telephones. Western Union, already using telegraph exchanges, quickly extended the principle to its telephones in New York City and San Francisco, and Bell was not slow in appreciating the potential.
|
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|
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+
Signalling began in an appropriately primitive manner. The user alerted the other end, or the exchange operator, by whistling into the transmitter. Exchange operation soon resulted in telephones being equipped with a bell in a ringer box, first operated over a second wire, and later over the same wire, but with a condenser (capacitor) in series with the bell coil to allow the AC ringer signal through while still blocking DC (keeping the phone "on hook"). Telephones connected to the earliest Strowger switch automatic exchanges had seven wires, one for the knife switch, one for each telegraph key, one for the bell, one for the push-button and two for speaking. Large wall telephones in the early 20th century usually incorporated the bell, and separate bell boxes for desk phones dwindled away in the middle of the century.
|
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+
|
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+
Rural and other telephones that were not on a common battery exchange had a magneto hand-cranked generator to produce a high voltage alternating signal to ring the bells of other telephones on the line and to alert the operator. Some local farming communities that were not connected to the main networks set up barbed wire telephone lines that exploited the existing system of field fences to transmit the signal.
|
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+
|
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+
In the 1890s a new smaller style of telephone was introduced, packaged in three parts. The transmitter stood on a stand, known as a "candlestick" for its shape. When not in use, the receiver hung on a hook with a switch in it, known as a "switchhook". Previous telephones required the user to operate a separate switch to connect either the voice or the bell. With the new kind, the user was less likely to leave the phone "off the hook". In phones connected to magneto exchanges, the bell, induction coil, battery and magneto were in a separate bell box or "ringer box".[13] In phones connected to common battery exchanges, the ringer box was installed under a desk, or other out-of-the-way place, since it did not need a battery or magneto.
|
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+
|
59 |
+
Cradle designs were also used at this time, having a handle with the receiver and transmitter attached, now called a handset, separate from the cradle base that housed the magneto crank and other parts. They were larger than the "candlestick" and more popular.
|
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+
|
61 |
+
Disadvantages of single-wire operation such as crosstalk and hum from nearby AC power wires had already led to the use of twisted pairs and, for long-distance telephones, four-wire circuits. Users at the beginning of the 20th century did not place long-distance calls from their own telephones but made an appointment to use a special soundproofed long-distance telephone booth furnished with the latest technology.
|
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|
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+
What turned out to be the most popular and longest-lasting physical style of telephone was introduced in the early 20th century, including Bell's 202-type desk set. A carbon granule transmitter and electromagnetic receiver were united in a single molded plastic handle, which when not in use sat in a cradle in the base unit. The circuit diagram of the model 202 shows the direct connection of the transmitter to the line, while the receiver was induction coupled. In local battery configurations, when the local loop was too long to provide sufficient current from the exchange, the transmitter was powered by a local battery and inductively coupled, while the receiver was included in the local loop.[14] The coupling transformer and the ringer were mounted in a separate enclosure, called the subscriber set. The dial switch in the base interrupted the line current by repeatedly but very briefly disconnecting the line 1 to 10 times for each digit, and the hook switch (in the center of the circuit diagram) disconnected the line and the transmitter battery while the handset was on the cradle.
|
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|
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+
In the 1930s, telephone sets were developed that combined the bell and induction coil with the desk set, obviating a separate ringer box. The rotary dial becoming commonplace in the 1930s in many areas enabled customer-dialed service, but some magneto systems remained even into the 1960s. After World War II, the telephone networks saw rapid expansion and more efficient telephone sets, such as the model 500 telephone in the United States, were developed that permitted larger local networks centered around central offices. A breakthrough new technology was the introduction of Touch-Tone signaling using push-button telephones by American Telephone & Telegraph Company (AT&T) in 1963.
|
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|
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+
Ericsson DBH 1001 (ca. 1931), the first combined telephone made with a Bakelite housing and handset.
|
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|
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+
Telephone used by American soldiers (WWII, Minalin, Pampanga, Philippines)
|
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|
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+
Video shows the operation of an Ericofon
|
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|
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+
Modern sound-powered emergency telephone
|
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|
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+
One type of mobile phone, called a cell phone
|
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+
|
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+
The invention of the transistor in 1947 dramatically changed the technology used in telephone systems and in the long-distance transmission networks, over the next several decades. Along with the development of stored program control for electronic switching systems, and new transmission technologies, such as pulse-code modulation (PCM), telephony gradually evolved towards digital telephony, which improved the capacity, quality, and cost of the network.[15]
|
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|
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+
The development of digital data communications methods made it possible to digitize voice and transmit it as real-time data across computer networks and the Internet, giving rise to the field of Internet Protocol (IP) telephony, also known as voice over Internet Protocol (VoIP), a term that reflects the methodology memorably. VoIP has proven to be a disruptive technology that is rapidly replacing traditional telephone network infrastructure.
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|
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+
By January 2005, up to 10% of telephone subscribers in Japan and South Korea had switched to this digital telephone service. A January 2005 Newsweek article suggested that Internet telephony may be "the next big thing."[16] The technology has spawned a new industry comprising many VoIP companies that offer services to consumers and businesses.
|
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|
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+
IP telephony uses high-bandwidth Internet connections and specialized customer premises equipment to transmit telephone calls via the Internet, or any modern private data network. The customer equipment may be an analog telephone adapter (ATA) which translates the signals of a conventional analog telephone to packet-switched IP messages. IP Phones have these function combined in standalone device, and computer softphone applications use microphone and headset devices of a personal computer.
|
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+
|
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+
While traditional analog telephones are typically powered from the central office through the telephone line, digital telephones require a local power supply. Internet-based digital service also requires special provisions to provide the service location to the emergency services when an emergency telephone number is called.
|
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|
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+
In 2002, only 10% of the world's population used mobile phones and by 2005 that percentage had risen to 46%.[17] By the end of 2009, there were a total of nearly 6 billion mobile and fixed-line telephone subscribers worldwide. This included 1.26 billion fixed-line subscribers and 4.6 billion mobile subscribers.[18]
|
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The Unicode system provides various code points for graphic symbols used in designating telephone devices, services, or information, for print, signage, and other media.
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A mobile phone, cellular phone, cell phone, cellphone or hand phone, sometimes shortened to simply mobile, cell or just phone, is a portable telephone that can make and receive calls over a radio frequency link while the user is moving within a telephone service area. The radio frequency link establishes a connection to the switching systems of a mobile phone operator, which provides access to the public switched telephone network (PSTN). Modern mobile telephone services use a cellular network architecture and, therefore, mobile telephones are called cellular telephones or cell phones in North America. In addition to telephony, digital mobile phones (2G) support a variety of other services, such as text messaging, MMS, email, Internet access, short-range wireless communications (infrared, Bluetooth), business applications, video games and digital photography. Mobile phones offering only those capabilities are known as feature phones; mobile phones which offer greatly advanced computing capabilities are referred to as smartphones.
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The development of metal-oxide-semiconductor (MOS) large-scale integration (LSI) technology, information theory and cellular networking led to the development of affordable mobile communications.[1] The first handheld mobile phone was demonstrated by John F. Mitchell[2][3] and Martin Cooper of Motorola in 1973, using a handset weighing c. 2 kilograms (4.4 lbs).[4] In 1979, Nippon Telegraph and Telephone (NTT) launched the world's first cellular network in Japan.[citation needed] In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone. From 1983 to 2014, worldwide mobile phone subscriptions grew to over seven billion; enough to provide one for every person on Earth.[5] In the first quarter of 2016, the top smartphone developers worldwide were Samsung, Apple and Huawei; smartphone sales represented 78 percent of total mobile phone sales.[6] For feature phones (slang: “dumbphones”) as of 2016[update], the largest were Samsung, Nokia and Alcatel.[7]
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A handheld mobile radio telephone service was envisioned in the early stages of radio engineering. In 1917, Finnish inventor Eric Tigerstedt filed a patent for a "pocket-size folding telephone with a very thin carbon microphone". Early predecessors of cellular phones included analog radio communications from ships and trains. The race to create truly portable telephone devices began after World War II, with developments taking place in many countries. The advances in mobile telephony have been traced in successive "generations", starting with the early zeroth-generation (0G) services, such as Bell System's Mobile Telephone Service and its successor, the Improved Mobile Telephone Service. These 0G systems were not cellular, supported few simultaneous calls, and were very expensive.
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The development of metal-oxide-semiconductor (MOS) large-scale integration (LSI) technology, information theory and cellular networking led to the development of affordable mobile communications.[1] The first handheld cellular mobile phone was demonstrated by John F. Mitchell[2][3] and Martin Cooper of Motorola in 1973, using a handset weighing 2 kilograms (4.4 lb).[4] The first commercial automated cellular network (1G) analog was launched in Japan by Nippon Telegraph and Telephone in 1979. This was followed in 1981 by the simultaneous launch of the Nordic Mobile Telephone (NMT) system in Denmark, Finland, Norway, and Sweden.[8] Several other countries then followed in the early to mid-1980s. These first-generation (1G) systems could support far more simultaneous calls but still used analog cellular technology. In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone.
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Digital cellular networks appeared in the 1990s, enabled by the wide adoption of MOSFET-based RF power amplifiers (power MOSFET and LDMOS) and RF circuits (RF CMOS),[9][10][11] leading to the introduction of digital signal processing in wireless communications.[1] In 1991, the second-generation (2G) digital cellular technology was launched in Finland by Radiolinja on the GSM standard. This sparked competition in the sector as the new operators challenged the incumbent 1G network operators. The GSM standard is a European initiative expressed at the CEPT ("Conférence Européenne des Postes et Telecommunications", European Postal and Telecommunications conference). The Franco-German R&D cooperation demonstrated the technical feasibility, and in 1987 a Memorandum of Understanding was signed between 13 European countries who agreed to launch a commercial service by 1991. The first version of the GSM (=2G) standard had 6,000 pages. The IEEE/RSE awarded to Thomas Haug and Philippe Dupuis [fr] the 2018 James Clerk Maxwell medal for their contributions to the first digital mobile telephone standard.[12] In 2018, the GSM was used by over 5 billion people in over 220 countries. The GSM (2G) has evolved into 3G, 4G and 5G.
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The lithium-ion battery, an indispensable energy source for modern mobile phones,[13] was commercialized by Sony and Asahi Kasei in 1991.[14][15] In 2001, the third generation (3G) was launched in Japan by NTT DoCoMo on the WCDMA standard.[16] This was followed by 3.5G, 3G+ or turbo 3G enhancements based on the high-speed packet access (HSPA) family, allowing UMTS networks to have higher data transfer speeds and capacity.
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By 2009, it had become clear that, at some point, 3G networks would be overwhelmed by the growth of bandwidth-intensive applications, such as streaming media.[17] Consequently, the industry began looking to data-optimized fourth-generation technologies, with the promise of speed improvements up to ten-fold over existing 3G technologies. The first two commercially available technologies billed as 4G were the WiMAX standard, offered in North America by Sprint, and the LTE standard, first offered in Scandinavia by TeliaSonera.
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5G is a technology and term used in research papers and projects to denote the next major phase in mobile telecommunication standards beyond the 4G/IMT-Advanced standards. The term 5G is not officially used in any specification or official document yet made public by telecommunication companies or standardization bodies such as 3GPP, WiMAX Forum or ITU-R. New standards beyond 4G are currently being developed by standardization bodies, but they are at this time seen as under the 4G umbrella, not for a new mobile generation.
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Smartphones have a number of distinguishing features. The International Telecommunication Union measures those with Internet connection, which it calls Active Mobile-Broadband subscriptions (which includes tablets, etc.). In the developed world, smartphones have now overtaken the usage of earlier mobile systems. However, in the developing world, they account for around 50% of mobile telephony.
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Feature phone is a term typically used as a retronym to describe mobile phones which are limited in capabilities in contrast to a modern smartphone. Feature phones typically provide voice calling and text messaging functionality, in addition to basic multimedia and Internet capabilities, and other services offered by the user's wireless service provider. A feature phone has additional functions over and above a basic mobile phone which is only capable of voice calling and text messaging.[20][21] Feature phones and basic mobile phones tend to use a proprietary, custom-designed software and user interface. By contrast, smartphones generally use a mobile operating system that often shares common traits across devices.
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There are Orthodox Jewish religious restrictions which, by some interpretations, standard mobile telephones overstep. To deal with this problem, some rabbinical organizations have recommended that phones with text-messaging capability not be used by children.[22] Phones with restricted features are known as kosher phones and have rabbinical approval for use in Israel and elsewhere by observant Orthodox Jews. Although these phones are intended to prevent immodesty, some vendors report good sales to adults who prefer the simplicity of the devices. Some phones are approved for use by essential workers (such as health, security, and public service workers) on the sabbath (the use of any electrical device is generally prohibited during this time, other than to save lives, or reduce the risk of death or similar needs).[23]
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Mobile phones communicate with cell towers that are placed to give coverage across a telephone service area which is divided up into 'cells'. Each cell uses a different set of frequencies from neighboring cells, and will typically be covered by 3 towers placed at different locations. The cell towers are usually interconnected to each other and the phone network and the internet by wired connections. Due to bandwidth limitations each cell will have a maximum number of cell phones it can handle at once. The cells are therefore sized depending on the expected usage density, and may be much smaller in cities. In that case much lower transmitter powers are used to avoid broadcasting beyond the cell.
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In order to handle the high traffic, multiple towers can be setup in the same area(using different frequencies). This can be done permanently or temporarily such as at special events like at the Super Bowl, Taste of Chicago, State Fair, NYC New Year's Eve, hurricane hit cities, etc. where cell phone companies will bring a truck with equipment to host the abnormally high traffic with a portable cell.
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Cellular can greatly increase the capacity of simultaneous wireless phone calls. While a phone company for example, has a license to 1000 frequencies, each cell must use unique frequencies with each call using one of them when communicating. Because cells only slightly overlap, the same frequency can be reused. Example cell 1 uses frequency 1–500, next door cell uses frequency 501–1000, next door can reuse frequency 1–500. Cells 1 and 3 are not "touching" and do not overlap/communicate so each can reuse the same frequencies.[citation needed]
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This is even more greatly increased when phone companies implemented digital networks. With digital, one frequency can host multiple simultaneous calls increasing capacity even more.
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As a phone moves around, a phone will "hand off" - automatically disconnect and reconnect to the tower of another cell that gives the best reception.
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Additionally, short-range Wi-Fi infrastructure is often used by smartphones as much as possible as it offloads traffic from cell networks on to local area networks.
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The common components found on all mobile phones are:
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Low-end mobile phones are often referred to as feature phones and offer basic telephony. Handsets with more advanced computing ability through the use of native software applications are known as smartphones.
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Mobile phones have central processing units (CPUs), similar to those in computers, but optimised to operate in low power environments.
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Mobile CPU performance depends not only on the clock rate (generally given in multiples of hertz)[24] but also the memory hierarchy also greatly affects overall performance. Because of these problems, the performance of mobile phone CPUs is often more appropriately given by scores derived from various standardized tests to measure the real effective performance in commonly used applications.
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One of the main characteristics of phones is the screen. Depending on the device's type and design, the screen fills most or nearly all of the space on a device's front surface. Many smartphone displays have an aspect ratio of 16:9, but taller aspect ratios became more common in 2017.
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Screen sizes are often measured in diagonal inches or millimeters; feature phones generally have screen sizes below 90 millimetres (3.5 in). Phones with screens larger than 130 millimetres (5.2 in) are often called "phablets." Smartphones with screens over 115 millimetres (4.5 in) in size are commonly difficult to use with only a single hand, since most thumbs cannot reach the entire screen surface; they may need to be shifted around in the hand, held in one hand and manipulated by the other, or used in place with both hands. Due to design advances, some modern smartphones with large screen sizes and "edge-to-edge" designs have compact builds that improve their ergonomics, while the shift to taller aspect ratios have resulted in phones that have larger screen sizes whilst maintaining the ergonomics associated with smaller 16:9 displays.[25][26][27]
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Liquid-crystal displays are the most common; others are IPS, LED, OLED, and AMOLED displays. Some displays are integrated with pressure-sensitive digitizers, such as those developed by Wacom and Samsung,[28] and Apple's "3D Touch" system.
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In sound, smartphones and feature phones vary little. Some audio-quality enhancing features, such as Voice over LTE and HD Voice, have appeared and are often available on newer smartphones. Sound quality can remain a problem due to the design of the phone, the quality of the cellular network and compression algorithms used in long-distance calls.[29][30] Audio quality can be improved using a VoIP application over WiFi.[31] Cellphones have small speakers so that the user can use a speakerphone feature and talk to a person on the phone without holding it to their ear. The small speakers can also be used to listen to digital audio files of music or speech or watch videos with an audio component, without holding the phone close to the ear.
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The average phone battery lasts 2–3 years at best. Many of the wireless devices use a Lithium-Ion (Li-Ion) battery, which charges 500–2500 times, depending on how users take care of the battery and the charging techniques used.[32] It is only natural for these rechargeable batteries to chemically age, which is why the performance of the battery when used for a year or two will begin to deteriorate. Battery life can be extended by draining it regularly, not overcharging it, and keeping it away from heat.[33][34]
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Mobile phones require a small microchip called a Subscriber Identity Module or SIM card, in order to function. The SIM card is approximately the size of a small postage stamp and is usually placed underneath the battery in the rear of the unit. The SIM securely stores the service-subscriber key (IMSI) and the Ki used to identify and authenticate the user of the mobile phone. The SIM card allows users to change phones by simply removing the SIM card from one mobile phone and inserting it into another mobile phone or broadband telephony device, provided that this is not prevented by a SIM lock. The first SIM card was made in 1991 by Munich smart card maker Giesecke & Devrient for the Finnish wireless network operator Radiolinja.[citation needed]
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A hybrid mobile phone can hold up to four SIM cards, with a phone having an IMEI per SIM Card. SIM and R-UIM cards may be mixed together to allow both GSM and CDMA networks to be accessed. From 2010 onwards, such phones became popular in emerging markets,[35] and this was attributed to the desire to obtain the lowest on-net calling rate.
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Feature phones have basic software platforms.
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Smartphones have advanced software platforms.
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A mobile app is a computer program designed to run on a mobile device, such as a smartphone. The term "app" is a shortening of the term "software application".
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A common data application on mobile phones is Short Message Service (SMS) text messaging. The first SMS message was sent from a computer to a mobile phone in 1992 in the UK while the first person-to-person SMS from phone to phone was sent in Finland in 1993. The first mobile news service, delivered via SMS, was launched in Finland in 2000,[citation needed] and subsequently many organizations provided "on-demand" and "instant" news services by SMS. Multimedia Messaging Service (MMS) was introduced in 2001.[citation needed]
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The introduction of Apple's App Store for the iPhone and iPod Touch in July 2008 popularized manufacturer-hosted online distribution for third-party applications (software and computer programs) focused on a single platform. There are a huge variety of apps, including video games, music products and business tools. Up until that point, smartphone application distribution depended on third-party sources providing applications for multiple platforms, such as GetJar, Handango, Handmark, and PocketGear. Following the success of the App Store, other smartphone manufacturers launched application stores, such as Google's Android Market (later renamed to the Google Play Store) and RIM's BlackBerry App World and Android-related app stores like F-Droid. In February 2014, 93% of mobile developers were targeting smartphones first for mobile app development.[36]
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From 1983 to 1998, Motorola was market leader in mobile phones. Nokia was the market leader in mobile phones from 1998 to 2012.[38] In Q1 2012, Samsung surpassed Nokia, selling 93.5 million units as against Nokia's 82.7 million units. Samsung has retained its top position since then. In 2017, the top five manufacturers worldwide were Samsung (20.9%), Apple (14.0%), Huawei (9.8%), Oppo (5.7%), and Vivo (6.5%).[39] During Q2 2018, Huawei overtook Apple as the world's second-largest phone manufacturer.[40]
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The world's largest individual mobile operator by number of subscribers is China Mobile, which has over 902 million mobile phone subscribers as of June 2018[update].[41] Over 50 mobile operators have over ten million subscribers each, and over 150 mobile operators had at least one million subscribers by the end of 2009.[42] In 2014, there were more than seven billion mobile phone subscribers worldwide, a number that is expected to keep growing.
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Mobile phones are used for a variety of purposes, such as keeping in touch with family members, for conducting business, and in order to have access to a telephone in the event of an emergency. Some people carry more than one mobile phone for different purposes, such as for business and personal use. Multiple SIM cards may be used to take advantage of the benefits of different calling plans. For example, a particular plan might provide for cheaper local calls, long-distance calls, international calls, or roaming.
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The mobile phone has been used in a variety of diverse contexts in society. For example:
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In 1998, one of the first examples of distributing and selling media content through the mobile phone was the sale of ringtones by Radiolinja in Finland. Soon afterwards, other media content appeared, such as news, video games, jokes, horoscopes, TV content and advertising. Most early content for mobile phones tended to be copies of legacy media, such as banner advertisements or TV news highlight video clips. Recently, unique content for mobile phones has been emerging, from ringtones and ringback tones to mobisodes, video content that has been produced exclusively for mobile phones.
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In many countries, mobile phones are used to provide mobile banking services, which may include the ability to transfer cash payments by secure SMS text message. Kenya's M-PESA mobile banking service, for example, allows customers of the mobile phone operator Safaricom to hold cash balances which are recorded on their SIM cards. Cash can be deposited or withdrawn from M-PESA accounts at Safaricom retail outlets located throughout the country and can be transferred electronically from person to person and used to pay bills to companies.
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Branchless banking has also been successful in South Africa and the Philippines. A pilot project in Bali was launched in 2011 by the International Finance Corporation and an Indonesian bank, Bank Mandiri.[51]
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Another application of mobile banking technology is Zidisha, a US-based nonprofit micro-lending platform that allows residents of developing countries to raise small business loans from Web users worldwide. Zidisha uses mobile banking for loan disbursements and repayments, transferring funds from lenders in the United States to borrowers in rural Africa who have mobile phones and can use the Internet.[52]
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Mobile payments were first trialled in Finland in 1998 when two Coca-Cola vending machines in Espoo were enabled to work with SMS payments. Eventually, the idea spread and in 1999, the Philippines launched the country's first commercial mobile payments systems with mobile operators Globe and Smart.
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Some mobile phones can make mobile payments via direct mobile billing schemes, or through contactless payments if the phone and the point of sale support near field communication (NFC).[53] Enabling contactless payments through NFC-equipped mobile phones requires the co-operation of manufacturers, network operators, and retail merchants.[54][55]
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Mobile phones are commonly used to collect location data. While the phone is turned on, the geographical location of a mobile phone can be determined easily (whether it is being used or not) using a technique known as multilateration to calculate the differences in time for a signal to travel from the mobile phone to each of several cell towers near the owner of the phone.[56][57]
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The movements of a mobile phone user can be tracked by their service provider and if desired, by law enforcement agencies and their governments. Both the SIM card and the handset can be tracked.[56]
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China has proposed using this technology to track the commuting patterns of Beijing city residents.[58] In the UK and US, law enforcement and intelligence services use mobile phones to perform surveillance operations. They possess technology that enables them to activate the microphones in mobile phones remotely in order to listen to conversations which take place near the phone.[59][60]
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Hackers are able to track a phone's location, read messages, and record calls, just by knowing the phone number.[61]
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Mobile phone use while driving, including talking on the phone, texting, or operating other phone features, is common but controversial. It is widely considered dangerous due to distracted driving. Being distracted while operating a motor vehicle has been shown to increase the risk of accidents. In September 2010, the US National Highway Traffic Safety Administration (NHTSA) reported that 995 people were killed by drivers distracted by cell phones. In March 2011, a U.S. insurance company, State Farm Insurance, announced the results of a study which showed 19% of drivers surveyed accessed the Internet on a smartphone while driving.[62] Many jurisdictions prohibit the use of mobile phones while driving. In Egypt, Israel, Japan, Portugal, and Singapore, both handheld and hands-free use of a mobile phone (which uses a speakerphone) is banned. In other countries, including the UK and France and in many U.S. states, only handheld phone use is banned while hands-free use is permitted.
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A 2011 study reported that over 90% of college students surveyed text (initiate, reply or read) while driving.[63]
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The scientific literature on the dangers of driving while sending a text message from a mobile phone, or texting while driving, is limited. A simulation study at the University of Utah found a sixfold increase in distraction-related accidents when texting.[64]
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Due to the increasing complexity of mobile phones, they are often more like mobile computers in their available uses. This has introduced additional difficulties for law enforcement officials when attempting to distinguish one usage from another in drivers using their devices. This is more apparent in countries which ban both handheld and hands-free usage, rather than those which ban handheld use only, as officials cannot easily tell which function of the mobile phone is being used simply by looking at the driver. This can lead to drivers being stopped for using their device illegally for a phone call when, in fact, they were using the device legally, for example, when using the phone's incorporated controls for car stereo, GPS or satnav.
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A 2010 study reviewed the incidence of mobile phone use while cycling and its effects on behaviour and safety.[65] In 2013, a national survey in the US reported the number of drivers who reported using their cellphones to access the Internet while driving had risen to nearly one of four.[66] A study conducted by the University of Vienna examined approaches for reducing inappropriate and problematic use of mobile phones, such as using mobile phones while driving.[67]
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Accidents involving a driver being distracted by talking on a mobile phone have begun to be prosecuted as negligence similar to speeding. In the United Kingdom, from 27 February 2007, motorists who are caught using a hand-held mobile phone while driving will have three penalty points added to their license in addition to the fine of £60.[68] This increase was introduced to try to stem the increase in drivers ignoring the law.[69] Japan prohibits all mobile phone use while driving, including use of hands-free devices. New Zealand has banned hand-held cell phone use since 1 November 2009. Many states in the United States have banned texting on cell phones while driving. Illinois became the 17th American state to enforce this law.[70] As of July 2010[update], 30 states had banned texting while driving, with Kentucky becoming the most recent addition on 15 July.[71]
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Public Health Law Research maintains a list of distracted driving laws in the United States. This database of laws provides a comprehensive view of the provisions of laws that restrict the use of mobile communication devices while driving for all 50 states and the District of Columbia between 1992 when first law was passed, through 1 December 2010. The dataset contains information on 22 dichotomous, continuous or categorical variables including, for example, activities regulated (e.g., texting versus talking, hands-free versus handheld), targeted populations, and exemptions.[72]
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In 2010, an estimated 1500 pedestrians were injured in the US while using a cellphone and some jurisdictions have attempted to ban pedestrians from using their cellphones.[73][74]
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The effect of mobile phone radiation on human health is the subject of recent[when?] interest and study, as a result of the enormous increase in mobile phone usage throughout the world. Mobile phones use electromagnetic radiation in the microwave range, which some believe may be harmful to human health. A large body of research exists, both epidemiological and experimental, in non-human animals and in humans. The majority of this research shows no definite causative relationship between exposure to mobile phones and harmful biological effects in humans. This is often paraphrased simply as the balance of evidence showing no harm to humans from mobile phones, although a significant number of individual studies do suggest such a relationship, or are inconclusive. Other digital wireless systems, such as data communication networks, produce similar radiation.[citation needed]
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On 31 May 2011, the World Health Organization stated that mobile phone use may possibly represent a long-term health risk,[75][76] classifying mobile phone radiation as "possibly carcinogenic to humans" after a team of scientists reviewed studies on mobile phone safety.[77] The mobile phone is in category 2B, which ranks it alongside coffee and other possibly carcinogenic substances.[78][79]
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Some recent[when?] studies have found an association between mobile phone use and certain kinds of brain and salivary gland tumors. Lennart Hardell and other authors of a 2009 meta-analysis of 11 studies from peer-reviewed journals concluded that cell phone usage for at least ten years "approximately doubles the risk of being diagnosed with a brain tumor on the same ('ipsilateral') side of the head as that preferred for cell phone use".[80]
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One study of past mobile phone use cited in the report showed a "40% increased risk for gliomas (brain cancer) in the highest category of heavy users (reported average: 30 minutes per day over a 10‐year period)".[81] This is a reversal of the study's prior position that cancer was unlikely to be caused by cellular phones or their base stations and that reviews had found no convincing evidence for other health effects.[76][82] However, a study published 24 March 2012, in the British Medical Journal questioned these estimates because the increase in brain cancers has not paralleled the increase in mobile phone use.[83] Certain countries, including France, have warned against the use of mobile phones by minors in particular, due to health risk uncertainties.[84] Mobile pollution by transmitting electromagnetic waves can be decreased up to 90% by adopting the circuit as designed in mobile phone and mobile exchange.[85]
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In May 2016, preliminary findings of a long-term study by the U.S. government suggested that radio-frequency (RF) radiation, the type emitted by cell phones, can cause cancer.[86][87]
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A study by the London School of Economics found that banning mobile phones in schools could increase pupils' academic performance, providing benefits equal to one extra week of schooling per year.[88]
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Studies have shown that around 40–50% of the environmental impact of mobile phones occurs during the manufacture of their printed wiring boards and integrated circuits.[89]
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The average user replaces their mobile phone every 11 to 18 months,[90] and the discarded phones then contribute to electronic waste. Mobile phone manufacturers within Europe are subject to the WEEE directive, and Australia has introduced a mobile phone recycling scheme.[91]
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Apple Inc. had an advanced robotic disassembler and sorter called Liam specifically for recycling outdated or broken iPhones.[350]
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According to the Federal Communications Commission, one out of three robberies involve the theft of a cellular phone.[citation needed] Police data in San Francisco show that half of all robberies in 2012 were thefts of cellular phones.[citation needed] An online petition on Change.org, called Secure our Smartphones, urged smartphone manufacturers to install kill switches in their devices to make them unusable if stolen. The petition is part of a joint effort by New York Attorney General Eric Schneiderman and San Francisco District Attorney George Gascón and was directed to the CEOs of the major smartphone manufacturers and telecommunication carriers.[92] On 10 June 2013, Apple announced that it would install a "kill switch" on its next iPhone operating system, due to debut in October 2013.[93]
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All mobile phones have a unique identifier called IMEI. Anyone can report their phone as lost or stolen with their Telecom Carrier, and the IMEI would be blacklisted with a central registry.[94] Telecom carriers, depending upon local regulation can or must implement blocking of blacklisted phones in their network. There are, however, a number of ways to circumvent a blacklist. One method is to send the phone to a country where the telecom carriers are not required to implement the blacklisting and sell it there,[95] another involves altering the phone's IMEI number.[96] Even so, mobile phones typically have less value on the second-hand market if the phones original IMEI is blacklisted.
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An unusual example of a phone bill caused by theft (reported on 28 June 2018) was when a biological group in Poland put a GPS tracker on a white stork and released it; during autumn migration over the Blue Nile valley in eastern Sudan someone got hold of the stork's GPS tracker, and found in it a mobile-phone-type sim card, which he put in his mobile phone, and made 20 hours of calls on it, running up a bill of over 10,000 Polish zlotys (US$2,700) for the biological group.[97]
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Demand for metals used in mobile phones and other electronics fuelled the Second Congo War, which claimed almost 5.5 million lives.[98] In a 2012 news story, The Guardian reported: "In unsafe mines deep underground in eastern Congo, children are working to extract minerals essential for the electronics industry. The profits from the minerals finance the bloodiest conflict since the second world war; the war has lasted nearly 20 years and has recently flared up again. ... For the last 15 years, the Democratic Republic of the Congo has been a major source of natural resources for the mobile phone industry."[99] The company Fairphone has worked to develop a mobile phone that does not contain conflict minerals.
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en/5637.html.txt
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A telephone is a telecommunications device that permits two or more users to conduct a conversation when they are too far apart to be heard directly. A telephone converts sound, typically and most efficiently the human voice, into electronic signals that are transmitted via cables and other communication channels to another telephone which reproduces the sound to the receiving user. The term is derived from Greek: τῆλε (tēle, far) and φωνή (phōnē, voice), together meaning distant voice. A common short form of the term is phone, which has been in use since the early 20th century.
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In 1876, Alexander Graham Bell was the first to be granted a United States patent for a device that produced clearly intelligible replication of the human voice. This instrument was further developed by many others, and became rapidly indispensable in business, government, and in households.
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The essential elements of a telephone are a microphone (transmitter) to speak into and an earphone (receiver) which reproduces the voice in a distant location.[1] In addition, most telephones contain a ringer to announce an incoming telephone call, and a dial or keypad to enter a telephone number when initiating a call to another telephone. The receiver and transmitter are usually built into a handset which is held up to the ear and mouth during conversation. The dial may be located either on the handset or on a base unit to which the handset is connected. The transmitter converts the sound waves to electrical signals which are sent through a telephone network to the receiving telephone, which converts the signals into audible sound in the receiver or sometimes a loudspeaker. Telephones are duplex devices, meaning they permit transmission in both directions simultaneously.
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The first telephones were directly connected to each other from one customer's office or residence to another customer's location. Being impractical beyond just a few customers, these systems were quickly replaced by manually operated centrally located switchboards. These exchanges were soon connected together, eventually forming an automated, worldwide public switched telephone network. For greater mobility, various radio systems were developed for transmission between mobile stations on ships and automobiles in the mid-20th century. Hand-held mobile phones were introduced for personal service starting in 1973. In later decades their analog cellular system evolved into digital networks with greater capability and lower cost.
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Convergence has given most modern cell phones capabilities far beyond simple voice conversation. Most are smartphones, integrating all mobile communication and many computing needs.
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A traditional landline telephone system, also known as plain old telephone service (POTS), commonly carries both control and audio signals on the same twisted pair (C in diagram) of insulated wires, the telephone line. The control and signaling equipment consists of three components, the ringer, the hookswitch, and a dial. The ringer, or beeper, light or other device (A7), alerts the user to incoming calls. The hookswitch signals to the central office that the user has picked up the handset to either answer a call or initiate a call. A dial, if present, is used by the subscriber to transmit a telephone number to the central office when initiating a call. Until the 1960s dials used almost exclusively the rotary technology, which was replaced by dual-tone multi-frequency signaling (DTMF) with pushbutton telephones (A4).
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A major expense of wire-line telephone service is the outside wire plant. Telephones transmit both the incoming and outgoing speech signals on a single pair of wires. A twisted pair line rejects electromagnetic interference (EMI) and crosstalk better than a single wire or an untwisted pair. The strong outgoing speech signal from the microphone (transmitter) does not overpower the weaker incoming speaker (receiver) signal with sidetone because a hybrid coil (A3) and other components compensate the imbalance. The junction box (B) arrests lightning (B2) and adjusts the line's resistance (B1) to maximize the signal power for the line length. Telephones have similar adjustments for inside line lengths (A8). The line voltages are negative compared to earth, to reduce galvanic corrosion. Negative voltage attracts positive metal ions toward the wires.
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The landline telephone contains a switchhook (A4) and an alerting device, usually a ringer (A7), that remains connected to the phone line whenever the phone is "on hook" (i.e. the switch (A4) is open), and other components which are connected when the phone is "off hook". The off-hook components include a transmitter (microphone, A2), a receiver (speaker, A1), and other circuits for dialing, filtering (A3), and amplification.
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The place a telephone call, the calling party picks up the telephone's handset, thereby operating a lever which closes the hook switch (A4). This powers the telephone by connecting the transmission hybrid transformer, as well as the transmitter (microphone) and receiver (speaker) to the line. In this off-hook state, the telephone circuitry has a low resistance of typically than 300 ohms, which causes the flow of direct current (DC) in the line (C) from the telephone exchange. The exchange detects this current, attaches a digit receiver circuit to the line, and sends dial tone to indicate its readiness. On a modern push-button telephone, the caller then presses the number keys to send the telephone number of the destination, the called party. The keys control a tone generator circuit (not shown) that sends DTMF tones to the exchange. A rotary-dial telephone uses pulse dialing, sending electrical pulses, that the exchange counts to decode each digit of the telephone number. If the called party's line is available, the terminating exchange applies an intermittent alternating current (AC) ringing signal of 40 to 90 volts to alert the called party of the incoming call. If the called party's line is in use, however, the exchange returns a busy signal to the calling party. If the called party's line is in use but subscribes to call waiting service, the exchange sends an intermittent audible tone to the called party to indicate another call.
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The electromechanical ringer of a telephone (A7) is connected to the line through a capacitor (A6), which blocks direct current and passes the alternating current of the ringing power. The telephone draws no current when it is on hook, while a DC voltage is continually applied to the line. Exchange circuitry (D2) can send an alternating current down the line to activate the ringer and announce an incoming call. In manual service exchange areas, before dial service was installed, telephones had hand-cranked magneto generators to generate a ringing voltage back to the exchange or any other telephone on the same line. When a landline telephone is inactive (on hook), the circuitry at the telephone exchange detects the absence of direct current to indicate that the line is not in use.[2] When a party initiates a call to this line, the exchange sends the ringing signal. When the called party picks up the handset, they actuate a double-circuit switchhook (not shown) which may simultaneously disconnects the alerting device and connects the audio circuitry to the line. This, in turn, draws direct current through the line, confirming that the called phone is now active. The exchange circuitry turns off the ring signal, and both telephones are now active and connected through the exchange. The parties may now converse as long as both phones remain off hook. When a party hangs up, placing the handset back on the cradle or hook, direct current ceases in that line, signaling the exchange to disconnect the call.
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Calls to parties beyond the local exchange are carried over trunk lines which establish connections between exchanges. In modern telephone networks, fiber-optic cable and digital technology are often employed in such connections. Satellite technology may be used for communication over very long distances.
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In most landline telephones, the transmitter and receiver (microphone and speaker) are located in the handset, although in a speakerphone these components may be located in the base or in a separate enclosure. Powered by the line, the microphone (A2) produces a modulated electric current which varies its frequency and amplitude in response to the sound waves arriving at its diaphragm. The resulting current is transmitted along the telephone line to the local exchange then on to the other phone (via the local exchange or via a larger network), where it passes through the coil of the receiver (A3). The varying current in the coil produces a corresponding movement of the receiver's diaphragm, reproducing the original sound waves present at the transmitter.
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Along with the microphone and speaker, additional circuitry is incorporated to prevent the incoming speaker signal and the outgoing microphone signal from interfering with each other. This is accomplished through a hybrid coil (A3). The incoming audio signal passes through a resistor (A8) and the primary winding of the coil (A3) which passes it to the speaker (A1). Since the current path A8 – A3 has a far lower impedance than the microphone (A2), virtually all of the incoming signal passes through it and bypasses the microphone.
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At the same time the DC voltage across the line causes a DC current which is split between the resistor-coil (A8-A3) branch and the microphone-coil (A2-A3) branch. The DC current through the resistor-coil branch has no effect on the incoming audio signal. But the DC current passing through the microphone is turned into AC (in response to voice sounds) which then passes through only the upper branch of the coil's (A3) primary winding, which has far fewer turns than the lower primary winding. This causes a small portion of the microphone output to be fed back to the speaker, while the rest of the AC goes out through the phone line.
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A lineman's handset is a telephone designed for testing the telephone network, and may be attached directly to aerial lines and other infrastructure components.
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Before the development of the electric telephone, the term "telephone" was applied to other inventions, and not all early researchers of the electrical device called it "telephone". Perhaps the earliest use of the word for a communications system was the telephon created by Gottfried Huth in 1796. Huth proposed an alternative to the optical telegraph of Claude Chappe in which the operators in the signalling towers would shout to each other by means of what he called "speaking tubes", but would now be called giant megaphones.[4] A communication device for sailing vessels called a "telephone" was invented by the captain John Taylor in 1844. This instrument used four air horns to communicate with vessels in foggy weather.[5][6]
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Johann Philipp Reis used the term in reference to his invention, commonly known as the Reis telephone, in c. 1860. His device appears to be the first device based on conversion of sound into electrical impulses. The term telephone was adopted into the vocabulary of many languages. It is derived from the Greek: τῆλε, tēle, "far" and φωνή, phōnē, "voice", together meaning "distant voice".
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Credit for the invention of the electric telephone is frequently disputed. As with other influential inventions such as radio, television, the light bulb, and the computer, several inventors pioneered experimental work on voice transmission over a wire and improved on each other's ideas. New controversies over the issue still arise from time to time. Charles Bourseul, Antonio Meucci, Johann Philipp Reis, Alexander Graham Bell, and Elisha Gray, amongst others, have all been credited with the invention of the telephone.[7][2]
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Alexander Graham Bell was the first to be awarded a patent for the electric telephone by the United States Patent and Trademark Office (USPTO) in March 1876.[8] Before Bell's patent, the telephone transmitted sound in a way that was similar to the telegraph. This method used vibrations and circuits to send electrical pulses, but was missing key features. Bell found that this method produced a sound through intermittent currents, but in order for the telephone to work a fluctuating current reproduced sounds the best. The fluctuating currents became the basis for the working telephone, creating Bell's patent.[9] That first patent by Bell was the master patent of the telephone, from which other patents for electric telephone devices and features flowed.[10] The Bell patents were forensically victorious and commercially decisive.
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In 1876, shortly after Bell's patent application, Hungarian engineer Tivadar Puskás proposed the telephone switch, which allowed for the formation of telephone exchanges, and eventually networks.[11]
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In the United Kingdom the blower is used as a slang term for a telephone. The term came from navy slang for a speaking tube.[12]
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Early telephones were technically diverse. Some used a water microphone, some had a metal diaphragm that induced current in an electromagnet wound around a permanent magnet, and some were dynamic – their diaphragm vibrated a coil of wire in the field of a permanent magnet or the coil vibrated the diaphragm. The sound-powered dynamic variants survived in small numbers through the 20th century in military and maritime applications, where its ability to create its own electrical power was crucial. Most, however, used the Edison/Berliner carbon transmitter, which was much louder than the other kinds, even though it required an induction coil which was an impedance matching transformer to make it compatible with the impedance of the line. The Edison patents kept the Bell monopoly viable into the 20th century, by which time the network was more important than the instrument.
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Early telephones were locally powered, using either a dynamic transmitter or by the powering of a transmitter with a local battery. One of the jobs of outside plant personnel was to visit each telephone periodically to inspect the battery. During the 20th century, telephones powered from the telephone exchange over the same wires that carried the voice signals became common.
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Early telephones used a single wire for the subscriber's line, with ground return used to complete the circuit (as used in telegraphs). The earliest dynamic telephones also had only one port opening for sound, with the user alternately listening and speaking (or rather, shouting) into the same hole. Sometimes the instruments were operated in pairs at each end, making conversation more convenient but also more expensive.
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At first, the benefits of a telephone exchange were not exploited. Instead telephones were leased in pairs to a subscriber, who had to arrange for a telegraph contractor to construct a line between them, for example between a home and a shop. Users who wanted the ability to speak to several different locations would need to obtain and set up three or four pairs of telephones. Western Union, already using telegraph exchanges, quickly extended the principle to its telephones in New York City and San Francisco, and Bell was not slow in appreciating the potential.
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Signalling began in an appropriately primitive manner. The user alerted the other end, or the exchange operator, by whistling into the transmitter. Exchange operation soon resulted in telephones being equipped with a bell in a ringer box, first operated over a second wire, and later over the same wire, but with a condenser (capacitor) in series with the bell coil to allow the AC ringer signal through while still blocking DC (keeping the phone "on hook"). Telephones connected to the earliest Strowger switch automatic exchanges had seven wires, one for the knife switch, one for each telegraph key, one for the bell, one for the push-button and two for speaking. Large wall telephones in the early 20th century usually incorporated the bell, and separate bell boxes for desk phones dwindled away in the middle of the century.
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Rural and other telephones that were not on a common battery exchange had a magneto hand-cranked generator to produce a high voltage alternating signal to ring the bells of other telephones on the line and to alert the operator. Some local farming communities that were not connected to the main networks set up barbed wire telephone lines that exploited the existing system of field fences to transmit the signal.
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In the 1890s a new smaller style of telephone was introduced, packaged in three parts. The transmitter stood on a stand, known as a "candlestick" for its shape. When not in use, the receiver hung on a hook with a switch in it, known as a "switchhook". Previous telephones required the user to operate a separate switch to connect either the voice or the bell. With the new kind, the user was less likely to leave the phone "off the hook". In phones connected to magneto exchanges, the bell, induction coil, battery and magneto were in a separate bell box or "ringer box".[13] In phones connected to common battery exchanges, the ringer box was installed under a desk, or other out-of-the-way place, since it did not need a battery or magneto.
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Cradle designs were also used at this time, having a handle with the receiver and transmitter attached, now called a handset, separate from the cradle base that housed the magneto crank and other parts. They were larger than the "candlestick" and more popular.
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Disadvantages of single-wire operation such as crosstalk and hum from nearby AC power wires had already led to the use of twisted pairs and, for long-distance telephones, four-wire circuits. Users at the beginning of the 20th century did not place long-distance calls from their own telephones but made an appointment to use a special soundproofed long-distance telephone booth furnished with the latest technology.
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What turned out to be the most popular and longest-lasting physical style of telephone was introduced in the early 20th century, including Bell's 202-type desk set. A carbon granule transmitter and electromagnetic receiver were united in a single molded plastic handle, which when not in use sat in a cradle in the base unit. The circuit diagram of the model 202 shows the direct connection of the transmitter to the line, while the receiver was induction coupled. In local battery configurations, when the local loop was too long to provide sufficient current from the exchange, the transmitter was powered by a local battery and inductively coupled, while the receiver was included in the local loop.[14] The coupling transformer and the ringer were mounted in a separate enclosure, called the subscriber set. The dial switch in the base interrupted the line current by repeatedly but very briefly disconnecting the line 1 to 10 times for each digit, and the hook switch (in the center of the circuit diagram) disconnected the line and the transmitter battery while the handset was on the cradle.
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In the 1930s, telephone sets were developed that combined the bell and induction coil with the desk set, obviating a separate ringer box. The rotary dial becoming commonplace in the 1930s in many areas enabled customer-dialed service, but some magneto systems remained even into the 1960s. After World War II, the telephone networks saw rapid expansion and more efficient telephone sets, such as the model 500 telephone in the United States, were developed that permitted larger local networks centered around central offices. A breakthrough new technology was the introduction of Touch-Tone signaling using push-button telephones by American Telephone & Telegraph Company (AT&T) in 1963.
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Ericsson DBH 1001 (ca. 1931), the first combined telephone made with a Bakelite housing and handset.
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Telephone used by American soldiers (WWII, Minalin, Pampanga, Philippines)
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Video shows the operation of an Ericofon
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Modern sound-powered emergency telephone
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One type of mobile phone, called a cell phone
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The invention of the transistor in 1947 dramatically changed the technology used in telephone systems and in the long-distance transmission networks, over the next several decades. Along with the development of stored program control for electronic switching systems, and new transmission technologies, such as pulse-code modulation (PCM), telephony gradually evolved towards digital telephony, which improved the capacity, quality, and cost of the network.[15]
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The development of digital data communications methods made it possible to digitize voice and transmit it as real-time data across computer networks and the Internet, giving rise to the field of Internet Protocol (IP) telephony, also known as voice over Internet Protocol (VoIP), a term that reflects the methodology memorably. VoIP has proven to be a disruptive technology that is rapidly replacing traditional telephone network infrastructure.
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By January 2005, up to 10% of telephone subscribers in Japan and South Korea had switched to this digital telephone service. A January 2005 Newsweek article suggested that Internet telephony may be "the next big thing."[16] The technology has spawned a new industry comprising many VoIP companies that offer services to consumers and businesses.
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IP telephony uses high-bandwidth Internet connections and specialized customer premises equipment to transmit telephone calls via the Internet, or any modern private data network. The customer equipment may be an analog telephone adapter (ATA) which translates the signals of a conventional analog telephone to packet-switched IP messages. IP Phones have these function combined in standalone device, and computer softphone applications use microphone and headset devices of a personal computer.
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While traditional analog telephones are typically powered from the central office through the telephone line, digital telephones require a local power supply. Internet-based digital service also requires special provisions to provide the service location to the emergency services when an emergency telephone number is called.
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In 2002, only 10% of the world's population used mobile phones and by 2005 that percentage had risen to 46%.[17] By the end of 2009, there were a total of nearly 6 billion mobile and fixed-line telephone subscribers worldwide. This included 1.26 billion fixed-line subscribers and 4.6 billion mobile subscribers.[18]
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The Unicode system provides various code points for graphic symbols used in designating telephone devices, services, or information, for print, signage, and other media.
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A telescope is an optical instrument using lenses, curved mirrors, or a combination of both to observe distant objects, or various devices used to observe distant objects by their emission, absorption, or reflection of electromagnetic radiation.[1] The first known practical telescopes were refracting telescopes invented in the Netherlands at the beginning of the 17th century, by using glass lenses. They were used for both terrestrial applications and astronomy.
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The reflecting telescope, which uses mirrors to collect and focus light, was invented within a few decades of the first refracting telescope. In the 20th century, many new types of telescopes were invented, including radio telescopes in the 1930s and infrared telescopes in the 1960s. The word telescope now refers to a wide range of instruments capable of detecting different regions of the electromagnetic spectrum, and in some cases other types of detectors.
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The word telescope (from the Ancient Greek τῆλε, tele "far" and σκοπεῖν, skopein "to look or see"; τηλεσκόπος, teleskopos "far-seeing") was coined in 1611 by the Greek mathematician Giovanni Demisiani for one of Galileo Galilei's instruments presented at a banquet at the Accademia dei Lincei.[2][3] In the Starry Messenger, Galileo had used the term perspicillum.
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The earliest existing record of a telescope was a 1608 patent submitted to the government in the Netherlands by Middelburg spectacle maker Hans Lippershey for a refracting telescope.[4] The actual inventor is unknown but word of it spread through Europe. Galileo heard about it and, in 1609, built his own version, and made his telescopic observations of celestial objects.[5][6]
|
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The idea that the objective, or light-gathering element, could be a mirror instead of a lens was being investigated soon after the invention of the refracting telescope.[7] The potential advantages of using parabolic mirrors—reduction of spherical aberration and no chromatic aberration—led to many proposed designs and several attempts to build reflecting telescopes.[8] In 1668, Isaac Newton built the first practical reflecting telescope, of a design which now bears his name, the Newtonian reflector.
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The invention of the achromatic lens in 1733 partially corrected color aberrations present in the simple lens and enabled the construction of shorter, more functional refracting telescopes. Reflecting telescopes, though not limited by the color problems seen in refractors, were hampered by the use of fast tarnishing speculum metal mirrors employed during the 18th and early 19th century—a problem alleviated by the introduction of silver coated glass mirrors in 1857,[9] and aluminized mirrors in 1932.[10] The maximum physical size limit for refracting telescopes is about 1 meter (40 inches), dictating that the vast majority of large optical researching telescopes built since the turn of the 20th century have been reflectors. The largest reflecting telescopes currently have objectives larger than 10 m (33 feet), and work is underway on several 30-40m designs.
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The 20th century also saw the development of telescopes that worked in a wide range of wavelengths from radio to gamma-rays. The first purpose built radio telescope went into operation in 1937. Since then, a large variety of complex astronomical instruments have been developed.
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|
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The name "telescope" covers a wide range of instruments. Most detect electromagnetic radiation, but there are major differences in how astronomers must go about collecting light (electromagnetic radiation) in different frequency bands.
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|
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Telescopes may be classified by the wavelengths of light they detect:
|
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As wavelengths become longer, it becomes easier to use antenna technology to interact with electromagnetic radiation (although it is possible to make very tiny antenna). The near-infrared can be collected much like visible light, however in the far-infrared and submillimetre range, telescopes can operate more like a radio telescope. For example, the James Clerk Maxwell Telescope observes from wavelengths from 3 μm (0.003 mm) to 2000 μm (2 mm), but uses a parabolic aluminum antenna.[11] On the other hand, the Spitzer Space Telescope, observing from about 3 μm (0.003 mm) to 180 μm (0.18 mm) uses a mirror (reflecting optics). Also using reflecting optics, the Hubble Space Telescope with Wide Field Camera 3 can observe in the frequency range from about 0.2 μm (0.0002 mm) to 1.7 μm (0.0017 mm) (from ultra-violet to infrared light).[12]
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With photons of the shorter wavelengths, with the higher frequencies, glancing-incident optics, rather than fully reflecting optics are used. Telescopes such as TRACE and SOHO use special mirrors to reflect Extreme ultraviolet, producing higher resolution and brighter images than are otherwise possible. A larger aperture does not just mean that more light is collected, it also enables a finer angular resolution.
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Telescopes may also be classified by location: ground telescope, space telescope, or flying telescope. They may also be classified by whether they are operated by professional astronomers or amateur astronomers. A vehicle or permanent campus containing one or more telescopes or other instruments is called an observatory.
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An optical telescope gathers and focuses light mainly from the visible part of the electromagnetic spectrum (although some work in the infrared and ultraviolet).[13] Optical telescopes increase the apparent angular size of distant objects as well as their apparent brightness. In order for the image to be observed, photographed, studied, and sent to a computer, telescopes work by employing one or more curved optical elements, usually made from glass lenses and/or mirrors, to gather light and other electromagnetic radiation to bring that light or radiation to a focal point. Optical telescopes are used for astronomy and in many non-astronomical instruments, including: theodolites (including transits), spotting scopes, monoculars, binoculars, camera lenses, and spyglasses. There are three main optical types:
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A Fresnel Imager is a proposed ultra-lightweight design for a space telescope that uses a Fresnel lens to focus light.
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Beyond these basic optical types there are many sub-types of varying optical design classified by the task they perform such as astrographs, comet seekers and solar telescopes.
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Radio telescopes are directional radio antennas that typically employ a large dish to collect radio waves. The dishes are sometimes constructed of a conductive wire mesh whose openings are smaller than the wavelength being observed.
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Unlike an optical telescope, which produces a magnified image of the patch of sky being observed, a traditional radio telescope dish contains a single receiver and records a single time-varying signal characteristic of the observed region; this signal may be sampled at various frequencies. In some newer radio telescope designs, a single dish contains an array of several receivers; this is known as a focal-plane array.
|
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+
|
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+
By collecting and correlating signals simultaneously received by several dishes, high-resolution images can be computed. Such multi-dish arrays are known as astronomical interferometers and the technique is called aperture synthesis. The 'virtual' apertures of these arrays are similar in size to the distance between the telescopes. As of 2005, the record array size is many times the diameter of the Earth — utilizing space-based Very Long Baseline Interferometry (VLBI) telescopes such as the Japanese HALCA (Highly Advanced Laboratory for Communications and Astronomy) VSOP (VLBI Space Observatory Program) satellite.
|
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|
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+
Aperture synthesis is now also being applied to optical telescopes using optical interferometers (arrays of optical telescopes) and aperture masking interferometry at single reflecting telescopes.
|
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+
Radio telescopes are also used to collect microwave radiation, which has the advantage of being able to pass through the atmosphere and interstellar gas and dust clouds.
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Some radio telescopes are used by programs such as SETI and the Arecibo Observatory to search for extraterrestrial life.
|
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X-rays are much harder to collect and focus than electromagnetic radiation of longer wavelengths. X-ray telescopes can use X-ray optics, such as Wolter telescopes composed of ring-shaped 'glancing' mirrors made of heavy metals that are able to reflect the rays just a few degrees. The mirrors are usually a section of a rotated parabola and a hyperbola, or ellipse. In 1952, Hans Wolter outlined 3 ways a telescope could be built using only this kind of mirror.[15][16] Examples of observatories using this type of telescope are the Einstein Observatory, ROSAT, and the Chandra X-Ray Observatory. By 2010, Wolter focusing X-ray telescopes are possible up to photon energies of 79 keV.[14]
|
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|
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Higher energy X-ray and Gamma-ray telescopes refrain from focusing completely and use coded aperture masks: the patterns of the shadow the mask creates can be reconstructed to form an image.
|
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|
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+
X-ray and Gamma-ray telescopes are usually installed on Earth-orbiting satellites or high-flying balloons since the Earth's atmosphere is opaque to this part of the electromagnetic spectrum. An example of this type of telescope is the Fermi Gamma-ray Space Telescope.
|
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|
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The detection of very high energy gamma rays, with shorter wavelength and higher frequency than regular gamma rays, requires further specialization. An example of this type of observatory is VERITAS.
|
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|
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+
A discovery in 2012 may allow focusing gamma-ray telescopes.[17] At photon energies greater than 700 keV, the index of refraction starts to increase again.[17]
|
52 |
+
|
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+
Astronomy is not limited to using electromagnetic radiation. Additional information can be obtained by detecting other signals, with detectors analogous to telescopes. These are:
|
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|
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+
A telescope mount is a mechanical structure which supports a telescope. Telescope mounts are designed to support the mass of the telescope and allow for accurate pointing of the instrument. Many sorts of mounts have been developed over the years, with the majority of effort being put into systems that can track the motion of the stars as the Earth rotates. The two main types of tracking mount are:
|
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|
57 |
+
By the 21 century, although not a structure a type of control system called a GoTo telescope was more popular. In this case a computer software system can in part or whole direct the telescope to a certain coordinate in the sky.
|
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+
|
59 |
+
Since the atmosphere is opaque for most of the electromagnetic spectrum, only a few bands can be observed from the Earth's surface. These bands are visible – near-infrared and a portion of the radio-wave part of the spectrum. For this reason there are no X-ray or far-infrared ground-based telescopes as these have to be observed from orbit. Even if a wavelength is observable from the ground, it might still be advantageous to place a telescope on a satellite due to astronomical seeing.
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Different types of telescope, operating in different wavelength bands, provide different information about the same object. Together they provide a more comprehensive understanding.
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Telescopes that operate in the electromagnetic spectrum:
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*Links to categories.
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1 |
+
A telescope is an optical instrument using lenses, curved mirrors, or a combination of both to observe distant objects, or various devices used to observe distant objects by their emission, absorption, or reflection of electromagnetic radiation.[1] The first known practical telescopes were refracting telescopes invented in the Netherlands at the beginning of the 17th century, by using glass lenses. They were used for both terrestrial applications and astronomy.
|
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|
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The reflecting telescope, which uses mirrors to collect and focus light, was invented within a few decades of the first refracting telescope. In the 20th century, many new types of telescopes were invented, including radio telescopes in the 1930s and infrared telescopes in the 1960s. The word telescope now refers to a wide range of instruments capable of detecting different regions of the electromagnetic spectrum, and in some cases other types of detectors.
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The word telescope (from the Ancient Greek τῆλε, tele "far" and σκοπεῖν, skopein "to look or see"; τηλεσκόπος, teleskopos "far-seeing") was coined in 1611 by the Greek mathematician Giovanni Demisiani for one of Galileo Galilei's instruments presented at a banquet at the Accademia dei Lincei.[2][3] In the Starry Messenger, Galileo had used the term perspicillum.
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|
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The earliest existing record of a telescope was a 1608 patent submitted to the government in the Netherlands by Middelburg spectacle maker Hans Lippershey for a refracting telescope.[4] The actual inventor is unknown but word of it spread through Europe. Galileo heard about it and, in 1609, built his own version, and made his telescopic observations of celestial objects.[5][6]
|
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|
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The idea that the objective, or light-gathering element, could be a mirror instead of a lens was being investigated soon after the invention of the refracting telescope.[7] The potential advantages of using parabolic mirrors—reduction of spherical aberration and no chromatic aberration—led to many proposed designs and several attempts to build reflecting telescopes.[8] In 1668, Isaac Newton built the first practical reflecting telescope, of a design which now bears his name, the Newtonian reflector.
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The invention of the achromatic lens in 1733 partially corrected color aberrations present in the simple lens and enabled the construction of shorter, more functional refracting telescopes. Reflecting telescopes, though not limited by the color problems seen in refractors, were hampered by the use of fast tarnishing speculum metal mirrors employed during the 18th and early 19th century—a problem alleviated by the introduction of silver coated glass mirrors in 1857,[9] and aluminized mirrors in 1932.[10] The maximum physical size limit for refracting telescopes is about 1 meter (40 inches), dictating that the vast majority of large optical researching telescopes built since the turn of the 20th century have been reflectors. The largest reflecting telescopes currently have objectives larger than 10 m (33 feet), and work is underway on several 30-40m designs.
|
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|
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The 20th century also saw the development of telescopes that worked in a wide range of wavelengths from radio to gamma-rays. The first purpose built radio telescope went into operation in 1937. Since then, a large variety of complex astronomical instruments have been developed.
|
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|
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The name "telescope" covers a wide range of instruments. Most detect electromagnetic radiation, but there are major differences in how astronomers must go about collecting light (electromagnetic radiation) in different frequency bands.
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|
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Telescopes may be classified by the wavelengths of light they detect:
|
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|
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As wavelengths become longer, it becomes easier to use antenna technology to interact with electromagnetic radiation (although it is possible to make very tiny antenna). The near-infrared can be collected much like visible light, however in the far-infrared and submillimetre range, telescopes can operate more like a radio telescope. For example, the James Clerk Maxwell Telescope observes from wavelengths from 3 μm (0.003 mm) to 2000 μm (2 mm), but uses a parabolic aluminum antenna.[11] On the other hand, the Spitzer Space Telescope, observing from about 3 μm (0.003 mm) to 180 μm (0.18 mm) uses a mirror (reflecting optics). Also using reflecting optics, the Hubble Space Telescope with Wide Field Camera 3 can observe in the frequency range from about 0.2 μm (0.0002 mm) to 1.7 μm (0.0017 mm) (from ultra-violet to infrared light).[12]
|
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With photons of the shorter wavelengths, with the higher frequencies, glancing-incident optics, rather than fully reflecting optics are used. Telescopes such as TRACE and SOHO use special mirrors to reflect Extreme ultraviolet, producing higher resolution and brighter images than are otherwise possible. A larger aperture does not just mean that more light is collected, it also enables a finer angular resolution.
|
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+
Telescopes may also be classified by location: ground telescope, space telescope, or flying telescope. They may also be classified by whether they are operated by professional astronomers or amateur astronomers. A vehicle or permanent campus containing one or more telescopes or other instruments is called an observatory.
|
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An optical telescope gathers and focuses light mainly from the visible part of the electromagnetic spectrum (although some work in the infrared and ultraviolet).[13] Optical telescopes increase the apparent angular size of distant objects as well as their apparent brightness. In order for the image to be observed, photographed, studied, and sent to a computer, telescopes work by employing one or more curved optical elements, usually made from glass lenses and/or mirrors, to gather light and other electromagnetic radiation to bring that light or radiation to a focal point. Optical telescopes are used for astronomy and in many non-astronomical instruments, including: theodolites (including transits), spotting scopes, monoculars, binoculars, camera lenses, and spyglasses. There are three main optical types:
|
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|
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+
A Fresnel Imager is a proposed ultra-lightweight design for a space telescope that uses a Fresnel lens to focus light.
|
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|
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+
Beyond these basic optical types there are many sub-types of varying optical design classified by the task they perform such as astrographs, comet seekers and solar telescopes.
|
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+
|
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+
Radio telescopes are directional radio antennas that typically employ a large dish to collect radio waves. The dishes are sometimes constructed of a conductive wire mesh whose openings are smaller than the wavelength being observed.
|
32 |
+
|
33 |
+
Unlike an optical telescope, which produces a magnified image of the patch of sky being observed, a traditional radio telescope dish contains a single receiver and records a single time-varying signal characteristic of the observed region; this signal may be sampled at various frequencies. In some newer radio telescope designs, a single dish contains an array of several receivers; this is known as a focal-plane array.
|
34 |
+
|
35 |
+
By collecting and correlating signals simultaneously received by several dishes, high-resolution images can be computed. Such multi-dish arrays are known as astronomical interferometers and the technique is called aperture synthesis. The 'virtual' apertures of these arrays are similar in size to the distance between the telescopes. As of 2005, the record array size is many times the diameter of the Earth — utilizing space-based Very Long Baseline Interferometry (VLBI) telescopes such as the Japanese HALCA (Highly Advanced Laboratory for Communications and Astronomy) VSOP (VLBI Space Observatory Program) satellite.
|
36 |
+
|
37 |
+
Aperture synthesis is now also being applied to optical telescopes using optical interferometers (arrays of optical telescopes) and aperture masking interferometry at single reflecting telescopes.
|
38 |
+
|
39 |
+
Radio telescopes are also used to collect microwave radiation, which has the advantage of being able to pass through the atmosphere and interstellar gas and dust clouds.
|
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+
|
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+
Some radio telescopes are used by programs such as SETI and the Arecibo Observatory to search for extraterrestrial life.
|
42 |
+
|
43 |
+
X-rays are much harder to collect and focus than electromagnetic radiation of longer wavelengths. X-ray telescopes can use X-ray optics, such as Wolter telescopes composed of ring-shaped 'glancing' mirrors made of heavy metals that are able to reflect the rays just a few degrees. The mirrors are usually a section of a rotated parabola and a hyperbola, or ellipse. In 1952, Hans Wolter outlined 3 ways a telescope could be built using only this kind of mirror.[15][16] Examples of observatories using this type of telescope are the Einstein Observatory, ROSAT, and the Chandra X-Ray Observatory. By 2010, Wolter focusing X-ray telescopes are possible up to photon energies of 79 keV.[14]
|
44 |
+
|
45 |
+
Higher energy X-ray and Gamma-ray telescopes refrain from focusing completely and use coded aperture masks: the patterns of the shadow the mask creates can be reconstructed to form an image.
|
46 |
+
|
47 |
+
X-ray and Gamma-ray telescopes are usually installed on Earth-orbiting satellites or high-flying balloons since the Earth's atmosphere is opaque to this part of the electromagnetic spectrum. An example of this type of telescope is the Fermi Gamma-ray Space Telescope.
|
48 |
+
|
49 |
+
The detection of very high energy gamma rays, with shorter wavelength and higher frequency than regular gamma rays, requires further specialization. An example of this type of observatory is VERITAS.
|
50 |
+
|
51 |
+
A discovery in 2012 may allow focusing gamma-ray telescopes.[17] At photon energies greater than 700 keV, the index of refraction starts to increase again.[17]
|
52 |
+
|
53 |
+
Astronomy is not limited to using electromagnetic radiation. Additional information can be obtained by detecting other signals, with detectors analogous to telescopes. These are:
|
54 |
+
|
55 |
+
A telescope mount is a mechanical structure which supports a telescope. Telescope mounts are designed to support the mass of the telescope and allow for accurate pointing of the instrument. Many sorts of mounts have been developed over the years, with the majority of effort being put into systems that can track the motion of the stars as the Earth rotates. The two main types of tracking mount are:
|
56 |
+
|
57 |
+
By the 21 century, although not a structure a type of control system called a GoTo telescope was more popular. In this case a computer software system can in part or whole direct the telescope to a certain coordinate in the sky.
|
58 |
+
|
59 |
+
Since the atmosphere is opaque for most of the electromagnetic spectrum, only a few bands can be observed from the Earth's surface. These bands are visible – near-infrared and a portion of the radio-wave part of the spectrum. For this reason there are no X-ray or far-infrared ground-based telescopes as these have to be observed from orbit. Even if a wavelength is observable from the ground, it might still be advantageous to place a telescope on a satellite due to astronomical seeing.
|
60 |
+
|
61 |
+
Different types of telescope, operating in different wavelength bands, provide different information about the same object. Together they provide a more comprehensive understanding.
|
62 |
+
|
63 |
+
Telescopes that operate in the electromagnetic spectrum:
|
64 |
+
|
65 |
+
*Links to categories.
|
en/564.html.txt
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1 |
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Bartholomew JoJo Simpson is a fictional character in the American animated television series The Simpsons and part of the Simpson family. He is voiced by Nancy Cartwright and first appeared on television in The Tracey Ullman Show short "Good Night" on April 19, 1987. Cartoonist Matt Groening created and designed Bart while waiting in the lobby of James L. Brooks' office. Groening had been called to pitch a series of shorts based on his comic strip, Life in Hell, but instead decided to create a new set of characters. While the rest of the characters were named after Groening's family members, Bart's name is an anagram of the word brat. After appearing on The Tracey Ullman Show for three years, the Simpson family received its own series on Fox, which debuted December 17, 1989. Bart has appeared in every Simpsons episode except "Four Great Women and a Manicure".
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At ten years old, Bart is the eldest child and only son of Homer and Marge, and the brother of Lisa and Maggie. Bart's most prominent and popular character traits are his mischievousness, rebelliousness and disrespect for authority. Hallmarks of the character include his chalkboard gags in the opening sequence; his prank calls to Moe; and his catchphrases "Eat my shorts", "¡Ay, caramba!", "Don't have a cow, man!", and "I'm Bart Simpson. Who the hell are you?". However, the exception of "Ay, caramba!", these hallmarks have been retired or not often used. Bart has appeared in other media relating to The Simpsons – including video games, The Simpsons Movie, The Simpsons Ride, commercials, and comic books – and inspired an entire line of merchandise.
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In casting, Cartwright originally planned to audition for the role of Lisa, while Yeardley Smith tried out for Bart. Smith's voice was too high for a boy, so she was given the role of Lisa. Cartwright found that Lisa was not interesting at the time, so instead auditioned for Bart, which she thought was a better role.[2]
|
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During the first two seasons of The Simpsons, Bart was the show's breakout character and "Bartmania" ensued, spawning Bart Simpson-themed merchandise touting his rebellious attitude and pride at underachieving, which caused many parents and educators to cast him as a bad role model for children. Around the third season, the series started to focus more on the family as a whole, though Bart still remains a prominent character. Time named Bart one of the 100 most important people of the 20th century, and he was named "entertainer of the year" in 1990 by Entertainment Weekly. Cartwright has won several awards for voicing Bart, including a Primetime Emmy Award in 1992 and an Annie Award in 1995. In 2000, Bart, along with the rest of his family, was awarded a star on the Hollywood Walk of Fame.
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The Simpsons uses a floating timeline in which the characters do not age or age very little, and as such, the show is always assumed to be set in the current year. In several episodes, events have been linked to specific times, though sometimes this timeline has been contradicted in subsequent episodes. Bart's year of birth was stated in "I Married Marge" (season three, 1991) as being in the early 1980s.[3] In "Simpsorama" (season 26, 2014) Bart states his birthday as February 23.[4] He lived with his parents in the Lower East Side of Springfield until the Simpsons bought their first house. When Lisa was born, Bart was at first jealous of the attention she received, but he soon warmed to her when he discovered that "Bart" was her first word.[5] Bart's first day of school was in the early 1990s. His initial enthusiasm was crushed by an uncaring teacher and Marge became worried that something was truly wrong with Bart. One day during recess, Bart met Milhouse and started entertaining him and other students with various gestures and rude words. Principal Skinner told him "you've just started school, and the path you choose now may be the one you follow for the rest of your life! Now, what do you say?" In his moment of truth, Bart responded, "eat my shorts".[6] The episode "That '90s Show" (season nineteen, 2008) contradicted much of the backstory's time frame; for example, it was revealed that Homer and Marge were childless in the early 1990s.[7]
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Bart's hobbies include skateboarding, watching television (especially The Krusty the Clown Show which includes The Itchy & Scratchy Show), reading comic books (especially Radioactive Man), playing video games and generally causing mischief.[8] His favorite movies are Jaws and the Star Wars Trilogy. For the duration of the series, Bart has attended Springfield Elementary School and has been in Edna Krabappel's fourth grade class. While he is too young to hold a full-time job, he has had occasional part-time jobs. He works as a bartender at Fat Tony's social club in "Bart the Murderer" (season three, 1991);[9] as Krusty the Clown's assistant in "Bart Gets Famous" (season five, 1994);[10] as a doorman in Springfield's burlesque house, the Maison Derrière, in "Bart After Dark" (season eight, 1996);[11] and briefly owns his own factory in "Homer's Enemy" (season eight, 1997).[12]
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Matt Groening first conceived of Bart and the rest of the Simpson family in 1986, while waiting in the lobby of producer James L. Brooks' office. Groening had been called in to pitch a series of animated shorts for The Tracey Ullman Show, and had intended to present an adaptation of his Life in Hell comic strip. When he realized that animating Life in Hell would require him to rescind publication rights, Groening decided to go in another direction.[13] He hurriedly sketched out his version of a dysfunctional family, naming the characters after members of his own family. For the rebellious son, he substituted "Bart", an anagram of the word brat, for his own name,[13] as he decided it would have been too obvious for him to have named the character 'Matt'.[14] Bart's middle initial J is a "tribute" to animated characters such as Bullwinkle J. Moose and Rocket J. Squirrel from The Rocky and Bullwinkle Show, who received their middle initial from Jay Ward.[15][16] According to the book Bart Simpson's Guide to Life, Bart's full middle name is "JoJo".[17]
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Bart had originally been envisioned as "a much milder, troubled youth given to existential angst who talks to himself", but the character was changed based on Cartwright's voice acting.[18] Groening has credited several different figures with providing inspiration for Bart: Matt Groening's older brother Mark provided much of the motivation for Bart's attitude.[19][20][21] Bart was conceived as an extreme version of the typical misbehaving child character, merging all of the extreme traits of characters such as Tom Sawyer and Huckleberry Finn into one person.[14] Groening describes Bart as "what would happen if Eddie Haskell [from Leave It to Beaver] got his own show".[22] Groening has also said that he found the premise of Dennis the Menace disappointing and was inspired to create a character who was actually a menace.[23]
|
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|
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Bart made his debut with the rest of the Simpson family on April 19, 1987 in The Tracey Ullman Show short "Good Night".[24] In 1989, the shorts were adapted into The Simpsons, a half-hour series airing on the Fox Broadcasting Company. Bart and the Simpson family remained the main characters on this new show.[25]
|
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|
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The entire Simpson family was designed so that they would be recognizable in silhouette.[26] The family was crudely drawn, because Groening had submitted basic sketches to the animators, assuming they would clean them up; instead, they just traced over his drawings.[13] Bart's original design, which appeared in the first shorts, had spikier hair, and the spikes were of different lengths. The number was later limited to nine spikes, all of the same size.[27] At the time Groening was primarily drawing in black and "not thinking that [Bart] would eventually be drawn in color" gave him spikes which appear to be an extension of his head.[28] The features of Bart's character design are generally not used in other characters; for example, no other characters in current episodes have Bart's spiky hairline, although several background characters in the first few seasons shared the trait.[29]
|
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The basic rectangular shape of Bart's head is described by director Mark Kirkland as a coffee can. Homer's head is also rectangular (with a dome on top), while spheres are used for Marge, Lisa, and Maggie.[30] Different animators have different methods of drawing Bart. Former director Jeffrey Lynch starts off with a box, then adds the eyes, then the mouth, then the hair spikes, ear, and then the rest of the body. Matt Groening normally starts with the eyes, then the nose, and the rest of the outline of Bart's head. Many of the animators have trouble drawing Bart's spikes evenly; one trick they use is to draw one on the right, one on the left, one in the middle, then continue to add one in the middle of the blank space until there are nine. Originally, whenever Bart was to be drawn from an angle looking down so the top of his head was seen, Groening wanted there to be spikes along the outline of his head, and in the middle as well. Instead, Wes Archer and David Silverman drew him so that there was an outline of the spikes, then just a smooth patch in the middle because "it worked graphically."[31] In "The Blue and the Gray", Bart (along with Lisa and Maggie) finally questions why his hair has no visible border to separate head from hair.
|
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+
In the season seven (1995) episode "Treehouse of Horror VI", Bart (along with Homer) was computer animated into a three-dimensional character for the first time for the "Homer3" segment of the episode. The computer animation was provided by Pacific Data Images.[32] While designing the 3D model of the character, the animators did not know how they would show Bart's hair. They realized that there were vinyl Bart dolls in production and purchased one to use as a model.[32]
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Bart's voice is provided by Nancy Cartwright, who voices several other child characters on The Simpsons, including Nelson Muntz, Ralph Wiggum, Todd Flanders, and Kearney.[33] While the roles of Homer and Marge were given to Dan Castellaneta and Julie Kavner because they were already a part of The Tracey Ullman Show cast,[34] the producers decided to hold casting for the roles of Bart and Lisa. Yeardley Smith had initially been asked to audition for the role of Bart, but casting director Bonita Pietila believed her voice was too high. Smith later recalled, "I always sounded too much like a girl. I read two lines as Bart and they said, 'Thanks for coming!'"[35] Smith was given the role of Lisa instead.[36] On March 13, 1987, Nancy Cartwright went in to audition for the role of Lisa. After arriving at the audition, she found that Lisa was simply described as the "middle child" and at the time did not have much personality. Cartwright became more interested in the role of Bart, who was described as "devious, underachieving, school-hating, irreverent, [and] clever".[37] Matt Groening let her try out for the part instead, and upon hearing her read, gave her the job on the spot.[38] Cartwright is the only one of the six main Simpsons cast members who had been professionally trained in voice acting prior to working on the show.[39]
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Cartwright's normal speaking voice is said to have "no obvious traces of Bart".[40] The voice came naturally to Cartwright; prior to The Tracey Ullman Show, she had used elements of it in shows such as My Little Pony, Snorks, and Pound Puppies.[40] Cartwright describes Bart's voice as easy to perform, saying, "Some characters take a little bit more effort, upper respiratory control, whatever it is technically. But Bart is easy to do. I can just slip into that without difficulty."[40] She usually does five or six readings of every line in order to give the producers more to work with.[38] In flashforward episodes, Cartwright still provides the voice of Bart. For "Lisa's Wedding", (season six, 1995) Bart's voice was electronically lowered.[41]
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Despite Bart's fame, Cartwright is rarely recognized in public. When she is recognized and asked to perform Bart's voice in front of children, Cartwright refuses as it "freaks [them] out".[40] During the first season of The Simpsons, the Fox Network did not allow Cartwright to give interviews because they did not want to publicize that Bart was voiced by a woman.[42]
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Until 1998, Cartwright was paid $30,000 per episode. During a pay dispute in 1998, Fox threatened to replace the six main voice actors with new actors, going as far as preparing for casting of new voices.[43] The dispute was resolved and Cartwright received $125,000 per episode until 2004, when the voice actors demanded that they be paid $360,000 an episode.[43] The dispute was resolved a month later,[44] and Cartwright's pay rose to $250,000 per episode.[45] After salary renegotiations in 2008, the voice actors receive approximately $400,000 per episode.[46] Three years later, with Fox threatening to cancel the series unless production costs were cut, Cartwright and the other cast members accepted a 30 percent pay cut to just over $300,000 per episode.[47]
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In the opening sequence of many Simpsons episodes, the camera zooms in on Springfield Elementary School, where Bart can be seen writing a message on the chalkboard. This message, which changes from episode to episode, has become known as the "chalkboard gag".[48] Chalkboard messages may involve political humor such as "The First Amendment does not cover burping",[49] pop culture references such as "I can't see dead people",[50] and meta-references such as "I am not a 32-year-old woman" and "Nobody reads these anymore".[48] The animators are able to produce the chalkboard gags quickly and in some cases have changed them to fit current events. For example, the chalkboard gag for "Homer the Heretic" (season four, 1992) read, "I will not defame New Orleans." The gag had been written as an apology to the city for a controversial song in the previous week's episode, "A Streetcar Named Marge", which called the city a "home of pirates, drunks and whores".[51][52] Many episodes do not feature a chalkboard gag because a shorter opening title sequence, where the chalkboard gags are cut, is used to make more room for story and plot development.
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One of Bart's early hallmarks were his prank calls to Moe's Tavern owner Moe Szyslak in which Bart calls Moe and asks for a gag name. Moe tries to find that person in the bar, but rapidly realizes it is a prank call and (despite not knowing who actually made the call) angrily threatens Bart. These calls were based on a series of prank calls known as the Tube Bar recordings. Moe was based partly on Tube Bar owner Louis "Red" Deutsch, whose often profane responses inspired Moe's violent side.[53] The prank calls debuted in "Homer's Odyssey", (season one, 1990) the third episode to air, but were included in "Some Enchanted Evening", the first episode of the series that was produced.[54] As the series progressed, it became more difficult for the writers to come up with a fake name and to write Moe's angry response, so the pranks were dropped as a regular joke during the fourth season[54][55] but they have occasionally resurfaced on the show.[56]
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The catchphrase "Eat My Shorts" was an ad-lib by Cartwright in one of the original table readings, harking back to an incident when she was in high school. Cartwright was in the marching band at Fairmont High School, and one day while performing, the band chanted "Eat my shorts" rather than the usual "Fairmont West! Fairmont West!"[35][57]It could also be an homage to The Breakfast Club, as John Bender says the phrase at Principal Vernon, and John Bender would become the inspiration for another Matt Groening creation, Bender from Futurama.[58] Bart's other catchphrases, "¡Ay, caramba!" came from a Portuguese flamenco dancer[57] and "Don't have a cow!" had been around since the 1950s which derived from the British phrase "Don't have kittens";[57] both were featured on T-shirts manufactured during the production of the early seasons of The Simpsons.[59][60] "Cowabunga" is also commonly associated with Bart, although it was mostly used on the show after it had been used as a slogan on the T-shirts.[61] Reiss also stated the writers took the phrase from Chief Thunderthud on The Howdy Doody Show. The use of catchphrase-based humor was mocked in the episode "Bart Gets Famous" (season five, 1994) in which Bart lands a popular role on Krusty the Clown's show for saying the line "I didn't do it."[62] The writers chose the phrase "I didn't do it" because they wanted a "lousy" phrase "to point out how really crummy things can become really popular".[63]
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Bart commonly appears nude in the show, although in every case only his buttocks are visible.[64] In The Simpsons Movie (2007), Bart appears in a sequence where he is skateboarding while fully nude; several different items cover his genitalia, but for a brief moment his penis can be seen. The scene was one of the first worked on for the film, but the producers were nervous about the segment because they thought it would earn the movie an R rating.[65] Despite this, the film was rated PG-13 by the Motion Picture Association of America for "Irreverent Humor Throughout."[66] The scene was later included by Entertainment Weekly in their list of "30 Unforgettable Nude Scenes."[67]
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—Chris Turner, author[68]
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Bart's character traits of rebelliousness and disrespect for authority have been compared to that of America's founding fathers, and he has been described as an updated version of Tom Sawyer and Huckleberry Finn, rolled into one.[69] In his book Planet Simpson, Chris Turner describes Bart as a nihilist, a philosophical position that argues that existence is without objective meaning, purpose, or intrinsic value.[68]
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Bart's rebellious attitude has made him a disruptive student at Springfield Elementary School, where Bart is an underachiever and proud of it.[70] He is constantly at odds with his teacher Ms. Krabappel, Principal Skinner, and occasionally Groundskeeper Willie.[71] Bart does poorly in school and is well aware of it, having once declared, "I am dumb, okay? Dumb as a post! Think I'm happy about it?"[70] On one occasion, Lisa successfully proves that Bart is dumber than a hamster, although Bart ultimately outsmarts her in the end.[72] In "Separate Vocations" (season three, 1992) Bart becomes hall monitor and his grades go up, suggesting that he struggles mainly because he does not pay attention, not because he is stupid.[73] This idea is reinforced in "Brother's Little Helper", (season eleven, 1999) in which it is revealed that Bart suffers from attention deficit disorder.[74] His lack of smarts can also be attributed to the hereditary "Simpson Gene", which affects the intelligence of most male members of the Simpson family.[75] Although he gets into endless trouble and can be sadistic, shallow and selfish, Bart also exhibits many qualities of high integrity. He has, on a few occasions, helped Principal Skinner and Mrs. Krabappel:[76] In "Sweet Seymour Skinner's Baadasssss Song" (season five, 1994), Bart accidentally got Skinner fired and befriended him outside the school environment. Bart missed having Skinner as an adversary and got him rehired, knowing that this would mean that the two could no longer be friends.[77]
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Due to Bart's mischievousness and Homer's often uncaring and incompetent behavior, the two have a turbulent relationship. Bart regularly addresses Homer by his given name instead of "Dad", while Homer in turn often refers to him as "the boy".[78] Homer has a short temper and when enraged by Bart will strangle him on impulse in a cartoonishly violent manner.[79] One of the original ideas for the show was that Homer would be "very angry" and oppressive toward Bart, but these characteristics were toned down somewhat as their characters were explored.[80] Marge is a much more caring, understanding and nurturing parent than Homer, but she also refers to Bart as "a handful" and is often embarrassed by his antics.[81] In "Marge Be Not Proud", (season seven, 1995) she felt she was mothering Bart too much and began acting more distant towards him after he was caught shoplifting. At the beginning of the episode, Bart protested at her over-mothering but as her attitude changed, he felt bad and made it up to her.[82] Despite his attitude, Bart is sometimes willing to experience humiliation if it means pleasing his mom.[83] Marge has expressed an understanding for her "special little guy" and has defended him on many occasions. She once said "I know Bart can be a handful, but I also know what he's like inside. He's got a spark. It's not a bad thing... Of course, it makes him do bad things."[81]
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—Nancy Cartwright[84]
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Bart shares a sibling rivalry with his younger sister, Lisa, but has a buddy-like relationship with his youngest sister Maggie, due to her infant state. While Bart has often hurt Lisa, and even fought her physically, the two are often very close.[72][78] Bart cares for Lisa deeply and has always apologized for going too far.[85] He also believes Lisa to be his superior when it comes to solving problems and frequently goes to her for advice.[73] Bart is also highly protective of Lisa: When a bully destroys her box of cupcakes in "Bart the General", (season one, 1990), Bart immediately stands up for her.[86]
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Bart is portrayed as a popular cool boy[87] and has many friends at school. Out of all of them his best friend is Milhouse Van Houten, although Bart has at times shown embarrassment about their friendship.[87] Bart is a bad influence on Milhouse, and the two have been involved in a lot of mischief together. Because of this behavior, Milhouse's mother forbids Milhouse from playing with Bart in "Homer Defined" (season three, 1991). While at first he pretended that he did not care, Bart eventually realizes that he needs Milhouse, and Marge manages to convince Mrs. Van Houten to reconsider.[81] Milhouse is a frequent target for local bullies Nelson Muntz and his friends Jimbo, Dolph, and Kearney. At times, Bart also finds himself at the hands of their abuse. Despite being the more socially powerful of the two, Bart's social popularity has temporarily subsided various episodes either due to extreme embarrassment caused by his family or other people (or even himself) or an unfortunate coincidence. Milhouse describes their social standing as "Three and a half. We get beat up, but we get an explanation."[76] While Bart and the bullies have been adversaries at times, with Bart once declaring war on Nelson,[86] the school bullies actually like Bart for his ways and hang out with him at times, especially Nelson who eventually becomes close friends with him.[88]
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Bart is one of the biggest fans of children's television host Krusty the Clown. He once declared, "I've based my whole life on Krusty's teachings," and sleeps in a room filled with Krusty merchandise. He has helped the clown on many occasions, for example, foiling Sideshow Bob's attempt to frame Krusty for armed robbery in "Krusty Gets Busted" (season one, 1990), reuniting Krusty with his estranged father in "Like Father, Like Clown".[89] and helping Krusty return to the air with a comeback special and reignite his career in "Krusty Gets Kancelled".[90] For his part, Krusty has remained largely ignorant of Bart's help and treats Bart with disinterest.[89] One summer, Bart enthusiastically attended Kamp Krusty, which turned out to be a disaster, with Krusty nowhere to be seen. Bart keeps his hopes up by believing that Krusty would show up, but is soon pushed over the edge, and finally decides that he is sick of Krusty's shoddy merchandise and takes over the camp. Krusty immediately visits the camp in hopes of ending the conflict and manages to appease Bart.[91][92] One of the original ideas for the series was that Bart worshiped a television clown but had no respect for his father, although this was never directly explored.[93][94] Because of this original plan, Krusty's design is basically Homer in clown make-up.[35] When Bart foiled Sideshow Bob's plans in "Krusty Gets Busted", it sparked a long-standing feud between the two. The writers decided to have Bob repeatedly return to get revenge on Bart. They took the idea of the Coyote chasing the Road Runner and depicted Bob as an intelligent person obsessed with catching a bratty boy.[95] Bob has appeared in fourteen episodes, generally plotting various evil schemes which often have to do with getting revenge on Bart (and sometimes the entire Simpson family by proxy), but is always foiled in the end.[96]
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}}
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In 1990, Bart quickly became one of the most popular characters on television in what was termed "Bartmania".[97][98][99][100] He became the most prevalent Simpsons character on memorabilia, such as T-shirts. In the early 1990s, millions of T-shirts featuring Bart were sold;[101] as many as one million were sold on some days.[102] Believing Bart to be a bad role model, several American public schools banned T-shirts featuring Bart next to captions such as "I'm Bart Simpson. Who the hell are you?" and "Underachiever ('And proud of it, man!')".[103][104][105] The Simpsons merchandise sold well and generated $2 billion in revenue during the first 14 months of sales.[103] The success of Bart Simpson merchandise inspired an entire line of black market counterfeit items, especially T-shirts. Some featured Bart announcing various slogans, others depicted redesigns of the character, including "Teenage Mutant Ninja Bart, Air Simpson Bart, [and] RastaBart".[106] Matt Groening generally did not object to bootleg merchandise, but took exception to a series of "Nazi Bart" shirts which depicted Bart in Nazi uniform or as a white power skinhead.[107] 20th Century Fox sued the creator of the shirts, who eventually agreed to stop making them.[108]
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Bart became so associated with Fox that, when bidding in 1993 to show pro football, the network had to assure the NFL and reporters that the character would not announce games.[109] Due to the show's success, over the summer of 1990 Fox decided to switch The Simpsons' timeslot so that it would move from 8:00 p.m. ET on Sunday night to the same time on Thursday, where it would compete with The Cosby Show on NBC, the number one show at the time.[61][110] Through the summer, several news outlets published stories about the supposed "Bill vs. Bart" rivalry.[61][102] The August 31, 1990 issue of Entertainment Weekly featured a picture of Bill Cosby wearing a Bart Simpson T-shirt.[111] "Bart Gets an "F"" (season two, 1990) was the first episode to air against The Cosby Show, and it received a lower Nielsen rating, tying for eighth behind The Cosby Show, which had an 18.5 rating. The rating is based on the number of household televisions that were tuned into the show, but Nielsen Media Research estimated that 33.6 million viewers watched the episode, making it the number one show in terms of actual viewers that week. At the time, it was the most watched episode in the history of the Fox Network,[112] and it is still the highest rated episode in the history of The Simpsons.[113] Because of his popularity, Bart was often the most promoted member of the Simpson family in advertisements for the show, even for episodes in which he was not involved in the main plot.[114]
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Bart was described as "television's king of 1990",[115] "television's brightest new star"[116] and an "undiminished smash".[102] Entertainment Weekly named Bart the "entertainer of the year" for 1990, writing that "Bart has proved to be a rebel who's also a good kid, a terror who's easily terrorized, and a flake who astonishes us, and himself, with serious displays of fortitude."[117] In the United States congressional, senatorial and gubernatorial elections of 1990, Bart was one of the most popular write-in candidates, and in many areas was second only to Mickey Mouse amongst fictional characters.[118][119] In the 1990 Macy's Thanksgiving Day Parade, Bart made his debut as one of the giant helium-filled balloons for which the parade is known. The Bart Simpson balloon has appeared at every parade since.[120] This was referenced in The Simpsons in the episode "Bart vs. Thanksgiving", which aired the same day as the parade, where Homer tells Bart, "If you start building a balloon for every flash-in-the-pan cartoon character, you turn the parade into a farce!" Meanwhile, behind and unbeknownst to him, the television briefly shows a Bart Simpson balloon.[85][121]
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The album The Simpsons Sing the Blues was released in September 1990 and was a success, peaking at No. 3 on the Billboard 200[122] and becoming certified 2x platinum by the Recording Industry Association of America.[123] The first single from the album was the pop rap song "Do the Bartman", performed by Nancy Cartwright and released on November 20, 1990. The song was written by Bryan Loren, a friend of Michael Jackson.[124] Jackson was a fan of The Simpsons, especially Bart,[125] and had called the producers one night offering to write Bart a number one single and do a guest spot on the show.[126] Jackson eventually guest starred in the episode "Stark Raving Dad" (season three, 1991) under the pseudonym John Jay Smith.[127] While the song was never officially released as a single in the United States, it was successful in the United Kingdom. In 1991 it was the number one song in the UK for three weeks from February 16 to March 9 and was the seventh best-selling song of the year.[128] It sold half a million copies and was certified gold by the British Phonographic Industry on February 1, 1991.[129]
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—Matt Groening, in a 1998 interview when asked, "How do you respond to critics who consider Bart Simpson a dreadful role model for children?"[130]
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Bart's rebellious nature, which frequently resulted in no punishment for his misbehavior, led some parents and conservatives to characterize him as a poor role model for children.[131][132] Robert Bianco of the Pittsburgh Post-Gazette wrote that "[Bart] outwits his parents and outtalks his teachers; in short, he's the child we wish we'd been, and fear our children will become."[133] In schools, educators claimed that Bart was a "threat to learning" because of his "underachiever and proud of it" attitude and negative attitude regarding his education.[134] Others described him as "egotistical, aggressive and mean-spirited."[135] In response to the criticism, James L. Brooks said, "I'm very wary of television where everybody is supposed to be a role model, you don't run across that many role models in real life. Why should television be full of them?"[136]
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In 1990 William Bennett, who at the time was drug czar of the United States, visited a drug treatment centre in Pittsburgh and upon noticing a poster of Bart remarked, "You guys aren't watching The Simpsons, are you? That's not going to help you any."[137] When a backlash over the comment ensued, Bennett apologized, claiming he "was just kidding"[138] and saying "I'll sit down with the little spike head. We'll straighten this thing out."[139] In a 1991 interview, Bill Cosby described Bart as a bad role model for children, calling him "angry, confused, frustrated." In response, Matt Groening said, "That sums up Bart, all right. Most people are in a struggle to be normal. He thinks normal is very boring, and does things that others just wished they dare do."[140] On January 27, 1992, then-President George H. W. Bush said, "We are going to keep on trying to strengthen the American family, to make American families a lot more like the Waltons and a lot less like the Simpsons."[103] The writers rushed out a tongue-in-cheek reply in the form of a short segment which aired three days later before a rerun of "Stark Raving Dad" in which Bart replied, "Hey, we're just like the Waltons. We're praying for an end to the Depression, too."[141][142]
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Although there were many critics of the character, favorable comments came from several quarters. Columnist Erma Bombeck wrote, "Kids need to know that somewhere in this world is a contemporary who can pull off all the things they can only fantasize about, someone who can stick it to their parents once in a while and still be permitted to live."[134] In 2003, Bart placed first in a poll of parents in the United Kingdom who were asked "which made-up character had the most influence" on children under 12 years old.[143]
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In 1998, Time named Bart one of the 100 most important people of the 20th century. He was the only fictional character to make the list.[60] He had previously appeared on the cover of the edition of December 31, 1990.[144] He was also ranked No. 48 in TV Guide's "50 Greatest TV Stars of All Time" in 1996[145] and both he and Lisa ranked No. 11 in TV Guide's "Top 50 Greatest Cartoon Characters of All Time" in 2002.[146]
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At the 44th Primetime Emmy Awards in 1992, Cartwright won a Primetime Emmy Award for Outstanding Voice-Over Performance for voicing Bart in the season three episode "Separate Vocations". She shared the award with five other voice actors from The Simpsons.[147] Various episodes in which Bart is strongly featured have been nominated for Emmy Awards for Outstanding Animated Program, including "Radio Bart" in 1992, "Future-Drama" in 2005, "The Haw-Hawed Couple" in 2006 and "Homer's Phobia", which won the award in 1997.[148] In 1995, Cartwright won an Annie Award for "Voice Acting in the Field of Animation" for her portrayal of Bart in an episode.[149] In 2000, Bart and the rest of the Simpson family were awarded a star on the Hollywood Walk of Fame located at 7021 Hollywood Boulevard.[150]
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In 2014 Bart Simpson became the second mascot of Russian football club FC Zenit Saint Petersburg, wearing number 87 on his back (referring to The Simpsons debut in 1987; the club's first mascot is a blue-maned lion).[151]
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Alongside T-shirts, Bart has been included in various other The Simpsons-related merchandise, including air fresheners, baseball caps, bumper stickers,[102] cardboard standups, refrigerator magnets, key rings, buttons, dolls, posters, figurines, clocks, soapstone carvings, Chia Pets, bowling balls and boxer shorts.[152][153] The Bart Book, a book about Bart's personality and attributes, was released in 2004.[154][155] Other books include Bart Simpson's Guide to Life. The Simpsons and Philosophy: The D'oh! of Homer, which is not an official publication, includes a chapter analyzing Bart's character and comparing him to the "Nietzschean ideal".[156]
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Bart has appeared in other media relating to The Simpsons. He has appeared in every one of The Simpsons video games, including Bart vs. the World, Bart Simpson's Escape from Camp Deadly, Bart vs. the Space Mutants, Bart's House of Weirdness, Bart vs. The Juggernauts, Bartman Meets Radioactive Man, Bart's Nightmare, Bart & the Beanstalk and The Simpsons Game, released in 2007.[157] Alongside the television series, Bart regularly appears in issues of Simpsons Comics, which were first published on November 29, 1993 and are still issued monthly, and also has his own series called Bart Simpson Comics which have been released since 2000.[158][159] Bart also plays a role in The Simpsons Ride, launched in 2008 at Universal Studios Florida and Hollywood.[160] Bart appears as a playable character in the toys-to-life video game Lego Dimensions, released via a "Fun Pack" packaged with a Gravity Sprinter accessory in November 2015.[161]
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Bart, and other The Simpsons characters, have appeared in numerous television commercials for Nestlé's Butterfinger candy bars from 1990 to 2001, with the slogan "Nobody better lay a finger on my Butterfinger!"[162] Lisa would occasionally advertise it too. Matt Groening would later say that the Butterfinger advertising campaign was a large part of the reason why Fox decided to pick up the half-hour show.[163] The campaign was discontinued in 2001, much to the disappointment of Cartwright.[164] Bart has also appeared in commercials for Burger King and Ramada Inn.[164] In 2001, Kellogg's launched a brand of cereal called "Bart Simpson Peanut Butter Chocolate Crunch", which was available for a limited time.[165][166] Before the half-hour series went on the air, Matt Groening pitched Bart as a spokesperson for Jell-O. He wanted Bart to sing "J-E-L-L-O", then burp the letter O. His belief was that kids would try to do it the next day, but he was rejected.[167]
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On April 9, 2009, the United States Postal Service unveiled a series of five 44-cent stamps featuring Bart and the four other members of the Simpson family. They are the first characters, other than Sesame Street characters, to receive this accolade while the show is still in production. The stamps, designed by Matt Groening, were made available for purchase on May 7, 2009.[168][169]
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Television (TV), sometimes shortened to tele or telly, is a telecommunication medium used for transmitting moving images in monochrome (black and white), or in color, and in two or three dimensions and sound. The term can refer to a television set, a television show, or the medium of television transmission. Television is a mass medium for advertising, entertainment, news, and sports.
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Television became available in crude experimental forms in the late 1920s, but it would still be several years before the new technology would be marketed to consumers. After World War II, an improved form of black-and-white TV broadcasting became popular in the United States and Britain, and television sets became commonplace in homes, businesses, and institutions. During the 1950s, television was the primary medium for influencing public opinion.[1] In the mid-1960s, color broadcasting was introduced in the US and most other developed countries. The availability of multiple types of archival storage media such as Betamax and VHS tapes, high-capacity hard disk drives, DVDs, flash drives, high-definition Blu-ray Discs, and cloud digital video recorders has enabled viewers to watch pre-recorded material—such as movies—at home on their own time schedule. For many reasons, especially the convenience of remote retrieval, the storage of television and video programming now occurs on the cloud. At the end of the first decade of the 2000s, digital television transmissions greatly increased in popularity. Another development was the move from standard-definition television (SDTV) (576i, with 576 interlaced lines of resolution and 480i) to high-definition television (HDTV), which provides a resolution that is substantially higher. HDTV may be transmitted in various formats: 1080p, 1080i and 720p. Since 2010, with the invention of smart television, Internet television has increased the availability of television programs and movies via the Internet through streaming video services such as Netflix, Amazon Video, iPlayer and Hulu.
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In 2013, 79% of the world's households owned a television set.[2] The replacement of early bulky, high-voltage cathode ray tube (CRT) screen displays with compact, energy-efficient, flat-panel alternative technologies such as LCDs (both fluorescent-backlit and LED), OLED displays, and plasma displays was a hardware revolution that began with computer monitors in the late 1990s. Most TV sets sold in the 2000s were flat-panel, mainly LEDs. Major manufacturers announced the discontinuation of CRT, DLP, plasma, and even fluorescent-backlit LCDs by the mid-2010s.[3][4] In the near future, LEDs are expected to be gradually replaced by OLEDs.[5] Also, major manufacturers have announced that they will increasingly produce smart TVs in the mid-2010s.[6][7][8] Smart TVs with integrated Internet and Web 2.0 functions became the dominant form of television by the late 2010s.[9]
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Television signals were initially distributed only as terrestrial television using high-powered radio-frequency transmitters to broadcast the signal to individual television receivers. Alternatively television signals are distributed by coaxial cable or optical fiber, satellite systems and, since the 2000s via the Internet. Until the early 2000s, these were transmitted as analog signals, but a transition to digital television is expected to be completed worldwide by the late 2010s. A standard television set is composed of multiple internal electronic circuits, including a tuner for receiving and decoding broadcast signals. A visual display device which lacks a tuner is correctly called a video monitor rather than a television.
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The word television comes from Ancient Greek τῆλε (tèle), meaning 'far', and Latin visio, meaning 'sight'. The first documented usage of the term dates back to 1900, when the Russian scientist Constantin Perskyi used it in a paper that he presented in French at the 1st International Congress of Electricity, which ran from 18 to 25 August 1900 during the International World Fair in Paris. The Anglicised version of the term is first attested in 1907, when it was still "...a theoretical system to transmit moving images over telegraph or telephone wires".[10] It was "...formed in English or borrowed from French télévision."[10] In the 19th century and early 20th century, other "...proposals for the name of a then-hypothetical technology for sending pictures over distance were telephote (1880) and televista (1904)."[10] The abbreviation "TV" is from 1948. The use of the term to mean "a television set" dates from 1941.[10] The use of the term to mean "television as a medium" dates from 1927.[10] The slang term "telly" is more common in the UK. The slang term "the tube" or the "boob tube" derives from the bulky cathode ray tube used on most TVs until the advent of flat-screen TVs. Another slang term for the TV is "idiot box".[11] Also, in the 1940s and throughout the 1950s, during the early rapid growth of television programming and television-set ownership in the United States, another slang term became widely used in that period and continues to be used today to distinguish productions originally created for broadcast on television from films developed for presentation in movie theaters.[12] The "small screen", as both a compound adjective and noun, became specific references to television, while the "big screen" was used to identify productions made for theatrical release.[12]
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Facsimile transmission systems for still photographs pioneered methods of mechanical scanning of images in the early 19th century. Alexander Bain introduced the facsimile machine between 1843 and 1846. Frederick Bakewell demonstrated a working laboratory version in 1851.[citation needed] Willoughby Smith discovered the photoconductivity of the element selenium in 1873. As a 23-year-old German university student, Paul Julius Gottlieb Nipkow proposed and patented the Nipkow disk in 1884.[13] This was a spinning disk with a spiral pattern of holes in it, so each hole scanned a line of the image. Although he never built a working model of the system, variations of Nipkow's spinning-disk "image rasterizer" became exceedingly common.[14] Constantin Perskyi had coined the word television in a paper read to the International Electricity Congress at the International World Fair in Paris on 24 August 1900. Perskyi's paper reviewed the existing electromechanical technologies, mentioning the work of Nipkow and others.[15] However, it was not until 1907 that developments in amplification tube technology by Lee de Forest and Arthur Korn, among others, made the design practical.[16]
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The first demonstration of the live transmission of images was by Georges Rignoux and A. Fournier in Paris in 1909. A matrix of 64 selenium cells, individually wired to a mechanical commutator, served as an electronic retina. In the receiver, a type of Kerr cell modulated the light and a series of variously angled mirrors attached to the edge of a rotating disc scanned the modulated beam onto the display screen. A separate circuit regulated synchronization. The 8x8 pixel resolution in this proof-of-concept demonstration was just sufficient to clearly transmit individual letters of the alphabet. An updated image was transmitted "several times" each second.[17]
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In 1911, Boris Rosing and his student Vladimir Zworykin created a system that used a mechanical mirror-drum scanner to transmit, in Zworykin's words, "very crude images" over wires to the "Braun tube" (cathode ray tube or "CRT") in the receiver. Moving images were not possible because, in the scanner: "the sensitivity was not enough and the selenium cell was very laggy".[18]
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In 1921, Edouard Belin sent the first image via radio waves with his belinograph.[citation needed]
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By the 1920s, when amplification made television practical, Scottish inventor John Logie Baird employed the Nipkow disk in his prototype video systems. On 25 March 1925, Baird gave the first public demonstration of televised silhouette images in motion, at Selfridge's Department Store in London.[19] Since human faces had inadequate contrast to show up on his primitive system, he televised a ventriloquist's dummy named "Stooky Bill", whose painted face had higher contrast, talking and moving. By 26 January 1926, he demonstrated the transmission of the image of a face in motion by radio. This is widely regarded as the first television demonstration. The subject was Baird's business partner Oliver Hutchinson. Baird's system used the Nipkow disk for both scanning the image and displaying it. A bright light shining through a spinning Nipkow disk set with lenses projected a bright spot of light which swept across the subject. A Selenium photoelectric tube detected the light reflected from the subject and converted it into a proportional electrical signal. This was transmitted by AM radio waves to a receiver unit, where the video signal was applied to a neon light behind a second Nipkow disk rotating synchronized with the first. The brightness of the neon lamp was varied in proportion to the brightness of each spot on the image. As each hole in the disk passed by, one scan line of the image was reproduced. Baird's disk had 30 holes, producing an image with only 30 scan lines, just enough to recognize a human face. In 1927, Baird transmitted a signal over 438 miles (705 km) of telephone line between London and Glasgow.[citation needed]
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In 1928, Baird's company (Baird Television Development Company/Cinema Television) broadcast the first transatlantic television signal, between London and New York, and the first shore-to-ship transmission. In 1929, he became involved in the first experimental mechanical television service in Germany. In November of the same year, Baird and Bernard Natan of Pathé established France's first television company, Télévision-Baird-Natan. In 1931, he made the first outdoor remote broadcast, of The Derby.[20] In 1932, he demonstrated ultra-short wave television. Baird's mechanical system reached a peak of 240-lines of resolution on BBC television broadcasts in 1936, though the mechanical system did not scan the televised scene directly. Instead a 17.5mm film was shot, rapidly developed and then scanned while the film was still wet.[citation needed]
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An American inventor, Charles Francis Jenkins, also pioneered the television. He published an article on "Motion Pictures by Wireless" in 1913, but it was not until December 1923 that he transmitted moving silhouette images for witnesses; and it was on 13 June 1925, that he publicly demonstrated synchronized transmission of silhouette pictures. In 1925 Jenkins used the Nipkow disk and transmitted the silhouette image of a toy windmill in motion, over a distance of 5 miles (8 km), from a naval radio station in Maryland to his laboratory in Washington, D.C., using a lensed disk scanner with a 48-line resolution.[21][22] He was granted U.S. Patent No. 1,544,156 (Transmitting Pictures over Wireless) on 30 June 1925 (filed 13 March 1922).[23]
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Herbert E. Ives and Frank Gray of Bell Telephone Laboratories gave a dramatic demonstration of mechanical television on 7 April 1927. Their reflected-light television system included both small and large viewing screens. The small receiver had a 2-inch-wide by 2.5-inch-high screen (5 by 6 cm). The large receiver had a screen 24 inches wide by 30 inches high (60 by 75 cm). Both sets were capable of reproducing reasonably accurate, monochromatic, moving images. Along with the pictures, the sets received synchronized sound. The system transmitted images over two paths: first, a copper wire link from Washington to New York City, then a radio link from Whippany, New Jersey. Comparing the two transmission methods, viewers noted no difference in quality. Subjects of the telecast included Secretary of Commerce Herbert Hoover. A flying-spot scanner beam illuminated these subjects. The scanner that produced the beam had a 50-aperture disk. The disc revolved at a rate of 18 frames per second, capturing one frame about every 56 milliseconds. (Today's systems typically transmit 30 or 60 frames per second, or one frame every 33.3 or 16.7 milliseconds respectively.) Television historian Albert Abramson underscored the significance of the Bell Labs demonstration: "It was in fact the best demonstration of a mechanical television system ever made to this time. It would be several years before any other system could even begin to compare with it in picture quality."[24]
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In 1928, WRGB, then W2XB, was started as the world's first television station. It broadcast from the General Electric facility in Schenectady, NY. It was popularly known as "WGY Television". Meanwhile, in the Soviet Union, Léon Theremin had been developing a mirror drum-based television, starting with 16 lines resolution in 1925, then 32 lines and eventually 64 using interlacing in 1926. As part of his thesis, on 7 May 1926, he electrically transmitted, and then projected, near-simultaneous moving images on a 5-square-foot (0.46 m2) screen.[22]
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By 1927, Theremin had achieved an image of 100 lines, a resolution that was not surpassed until May 1932 by RCA, with 120 lines.[25]
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On 25 December 1926, Kenjiro Takayanagi demonstrated a television system with a 40-line resolution that employed a Nipkow disk scanner and CRT display at Hamamatsu Industrial High School in Japan. This prototype is still on display at the Takayanagi Memorial Museum in Shizuoka University, Hamamatsu Campus. His research in creating a production model was halted by the SCAP after World War II.[26]
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Because only a limited number of holes could be made in the disks, and disks beyond a certain diameter became impractical, image resolution on mechanical television broadcasts was relatively low, ranging from about 30 lines up to 120 or so. Nevertheless, the image quality of 30-line transmissions steadily improved with technical advances, and by 1933 the UK broadcasts using the Baird system were remarkably clear.[27] A few systems ranging into the 200-line region also went on the air. Two of these were the 180-line system that Compagnie des Compteurs (CDC) installed in Paris in 1935, and the 180-line system that Peck Television Corp. started in 1935 at station VE9AK in Montreal.[28][29] The advancement of all-electronic television (including image dissectors and other camera tubes and cathode ray tubes for the reproducer) marked the beginning of the end for mechanical systems as the dominant form of television. Mechanical television, despite its inferior image quality and generally smaller picture, would remain the primary television technology until the 1930s. The last mechanical television broadcasts ended in 1939 at stations run by a handful of public universities in the United States.[citation needed]
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In 1897, English physicist J. J. Thomson was able, in his three famous experiments, to deflect cathode rays, a fundamental function of the modern cathode ray tube (CRT). The earliest version of the CRT was invented by the German physicist Ferdinand Braun in 1897 and is also known as the "Braun" tube.[30][31] It was a cold-cathode diode, a modification of the Crookes tube, with a phosphor-coated screen. In 1906 the Germans Max Dieckmann and Gustav Glage produced raster images for the first time in a CRT.[32] In 1907, Russian scientist Boris Rosing used a CRT in the receiving end of an experimental video signal to form a picture. He managed to display simple geometric shapes onto the screen.[33]
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In 1908 Alan Archibald Campbell-Swinton, fellow of the Royal Society (UK), published a letter in the scientific journal Nature in which he described how "distant electric vision" could be achieved by using a cathode ray tube, or Braun tube, as both a transmitting and receiving device,[34][35] He expanded on his vision in a speech given in London in 1911 and reported in The Times[36] and the Journal of the Röntgen Society.[37][38] In a letter to Nature published in October 1926, Campbell-Swinton also announced the results of some "not very successful experiments" he had conducted with G. M. Minchin and J. C. M. Stanton. They had attempted to generate an electrical signal by projecting an image onto a selenium-coated metal plate that was simultaneously scanned by a cathode ray beam.[39][40] These experiments were conducted before March 1914, when Minchin died,[41] but they were later repeated by two different teams in 1937, by H. Miller and J. W. Strange from EMI,[42] and by H. Iams and A. Rose from RCA.[43] Both teams succeeded in transmitting "very faint" images with the original Campbell-Swinton's selenium-coated plate. Although others had experimented with using a cathode ray tube as a receiver, the concept of using one as a transmitter was novel.[44] The first cathode ray tube to use a hot cathode was developed by John B. Johnson (who gave his name to the term Johnson noise) and Harry Weiner Weinhart of Western Electric, and became a commercial product in 1922.[citation needed]
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In 1926, Hungarian engineer Kálmán Tihanyi designed a television system utilizing fully electronic scanning and display elements and employing the principle of "charge storage" within the scanning (or "camera") tube.[45][46][47][48] The problem of low sensitivity to light resulting in low electrical output from transmitting or "camera" tubes would be solved with the introduction of charge-storage technology by Kálmán Tihanyi beginning in 1924.[49] His solution was a camera tube that accumulated and stored electrical charges ("photoelectrons") within the tube throughout each scanning cycle. The device was first described in a patent application he filed in Hungary in March 1926 for a television system he dubbed "Radioskop".[50] After further refinements included in a 1928 patent application,[49] Tihanyi's patent was declared void in Great Britain in 1930,[51] so he applied for patents in the United States. Although his breakthrough would be incorporated into the design of RCA's "iconoscope" in 1931, the U.S. patent for Tihanyi's transmitting tube would not be granted until May 1939. The patent for his receiving tube had been granted the previous October. Both patents had been purchased by RCA prior to their approval.[52][53] Charge storage remains a basic principle in the design of imaging devices for television to the present day.[50] On 25 December 1926, at Hamamatsu Industrial High School in Japan, Japanese inventor Kenjiro Takayanagi demonstrated a TV system with a 40-line resolution that employed a CRT display.[26] This was the first working example of a fully electronic television receiver. Takayanagi did not apply for a patent.[54]
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On 7 September 1927, American inventor Philo Farnsworth's image dissector camera tube transmitted its first image, a simple straight line, at his laboratory at 202 Green Street in San Francisco.[55][56] By 3 September 1928, Farnsworth had developed the system sufficiently to hold a demonstration for the press. This is widely regarded as the first electronic television demonstration.[56] In 1929, the system was improved further by the elimination of a motor generator, so that his television system now had no mechanical parts.[57] That year, Farnsworth transmitted the first live human images with his system, including a three and a half-inch image of his wife Elma ("Pem") with her eyes closed (possibly due to the bright lighting required).[58]
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Meanwhile, Vladimir Zworykin was also experimenting with the cathode ray tube to create and show images. While working for Westinghouse Electric in 1923, he began to develop an electronic camera tube. But in a 1925 demonstration, the image was dim, had low contrast, and poor definition, and was stationary.[59] Zworykin's imaging tube never got beyond the laboratory stage. But RCA, which acquired the Westinghouse patent, asserted that the patent for Farnsworth's 1927 image dissector was written so broadly that it would exclude any other electronic imaging device. Thus RCA, on the basis of Zworykin's 1923 patent application, filed a patent interference suit against Farnsworth. The U.S. Patent Office examiner disagreed in a 1935 decision, finding priority of invention for Farnsworth against Zworykin. Farnsworth claimed that Zworykin's 1923 system would be unable to produce an electrical image of the type to challenge his patent. Zworykin received a patent in 1928 for a color transmission version of his 1923 patent application;[60] he also divided his original application in 1931.[61] Zworykin was unable or unwilling to introduce evidence of a working model of his tube that was based on his 1923 patent application. In September 1939, after losing an appeal in the courts, and determined to go forward with the commercial manufacturing of television equipment, RCA agreed to pay Farnsworth US$1 million over a ten-year period, in addition to license payments, to use his patents.[62][63]
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In 1933, RCA introduced an improved camera tube that relied on Tihanyi's charge storage principle.[64] Dubbed the "Iconoscope" by Zworykin, the new tube had a light sensitivity of about 75,000 lux, and thus was claimed to be much more sensitive than Farnsworth's image dissector.[citation needed] However, Farnsworth had overcome his power problems with his Image Dissector through the invention of a completely unique "multipactor" device that he began work on in 1930, and demonstrated in 1931.[65][66] This small tube could amplify a signal reportedly to the 60th power or better[67] and showed great promise in all fields of electronics. Unfortunately, a problem with the multipactor was that it wore out at an unsatisfactory rate.[68]
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At the Berlin Radio Show in August 1931, Manfred von Ardenne gave a public demonstration of a television system using a CRT for both transmission and reception. However, Ardenne had not developed a camera tube, using the CRT instead as a flying-spot scanner to scan slides and film.[69] Philo Farnsworth gave the world's first public demonstration of an all-electronic television system, using a live camera, at the Franklin Institute of Philadelphia on 25 August 1934, and for ten days afterwards.[70][71] Mexican inventor Guillermo González Camarena also played an important role in early TV. His experiments with TV (known as telectroescopía at first) began in 1931 and led to a patent for the "trichromatic field sequential system" color television in 1940.[72] In Britain, the EMI engineering team led by Isaac Shoenberg applied in 1932 for a patent for a new device they dubbed "the Emitron",[73][74] which formed the heart of the cameras they designed for the BBC. On 2 November 1936, a 405-line broadcasting service employing the Emitron began at studios in Alexandra Palace, and transmitted from a specially built mast atop one of the Victorian building's towers. It alternated for a short time with Baird's mechanical system in adjoining studios, but was more reliable and visibly superior. This was the world's first regular "high-definition" television service.[75]
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The original American iconoscope was noisy, had a high ratio of interference to signal, and ultimately gave disappointing results, especially when compared to the high definition mechanical scanning systems then becoming available.[76][77] The EMI team, under the supervision of Isaac Shoenberg, analyzed how the iconoscope (or Emitron) produces an electronic signal and concluded that its real efficiency was only about 5% of the theoretical maximum.[78][79] They solved this problem by developing, and patenting in 1934, two new camera tubes dubbed super-Emitron and CPS Emitron.[80][81][82] The super-Emitron was between ten and fifteen times more sensitive than the original Emitron and iconoscope tubes and, in some cases, this ratio was considerably greater.[78] It was used for outside broadcasting by the BBC, for the first time, on Armistice Day 1937, when the general public could watch on a television set as the King laid a wreath at the Cenotaph.[83] This was the first time that anyone had broadcast a live street scene from cameras installed on the roof of neighboring buildings, because neither Farnsworth nor RCA would do the same until the 1939 New York World's Fair.
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On the other hand, in 1934, Zworykin shared some patent rights with the German licensee company Telefunken.[84] The "image iconoscope" ("Superikonoskop" in Germany) was produced as a result of the collaboration. This tube is essentially identical to the super-Emitron.[citation needed] The production and commercialization of the super-Emitron and image iconoscope in Europe were not affected by the patent war between Zworykin and Farnsworth, because Dieckmann and Hell had priority in Germany for the invention of the image dissector, having submitted a patent application for their Lichtelektrische Bildzerlegerröhre für Fernseher (Photoelectric Image Dissector Tube for Television) in Germany in 1925,[85] two years before Farnsworth did the same in the United States.[86] The image iconoscope (Superikonoskop) became the industrial standard for public broadcasting in Europe from 1936 until 1960, when it was replaced by the vidicon and plumbicon tubes. Indeed, it was the representative of the European tradition in electronic tubes competing against the American tradition represented by the image orthicon.[87][88] The German company Heimann produced the Superikonoskop for the 1936 Berlin Olympic Games,[89][90] later Heimann also produced and commercialized it from 1940 to 1955;[91] finally the Dutch company Philips produced and commercialized the image iconoscope and multicon from 1952 to 1958.[88][92]
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American television broadcasting, at the time, consisted of a variety of markets in a wide range of sizes, each competing for programming and dominance with separate technology, until deals were made and standards agreed upon in 1941.[93] RCA, for example, used only Iconoscopes in the New York area, but Farnsworth Image Dissectors in Philadelphia and San Francisco.[94] In September 1939, RCA agreed to pay the Farnsworth Television and Radio Corporation royalties over the next ten years for access to Farnsworth's patents.[95] With this historic agreement in place, RCA integrated much of what was best about the Farnsworth Technology into their systems.[94] In 1941, the United States implemented 525-line television.[96][97] Electrical engineer Benjamin Adler played a prominent role in the development of television.[98][99]
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The world's first 625-line television standard was designed in the Soviet Union in 1944 and became a national standard in 1946.[100] The first broadcast in 625-line standard occurred in Moscow in 1948.[101] The concept of 625 lines per frame was subsequently implemented in the European CCIR standard.[102] In 1936, Kálmán Tihanyi described the principle of plasma display, the first flat panel display system.[103][104]
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Early electronic television sets were large and bulky, with analog circuits made of vacuum tubes. Following the invention of the first working transistor at Bell Labs, Sony founder Masaru Ibuka predicted in 1952 that the transition to electronic circuits made of transistors would lead to smaller and more portable television sets.[105] The first fully transistorized, portable solid-state television set was the 8-inch Sony TV8-301, developed in 1959 and released in 1960.[106][107] This began the transformation of television viewership from a communal viewing experience to a solitary viewing experience.[108] By 1960, Sony had sold over 4 million portable television sets worldwide.[109]
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The basic idea of using three monochrome images to produce a color image had been experimented with almost as soon as black-and-white televisions had first been built. Although he gave no practical details, among the earliest published proposals for television was one by Maurice Le Blanc, in 1880, for a color system, including the first mentions in television literature of line and frame scanning.[110] Polish inventor Jan Szczepanik patented a color television system in 1897, using a selenium photoelectric cell at the transmitter and an electromagnet controlling an oscillating mirror and a moving prism at the receiver. But his system contained no means of analyzing the spectrum of colors at the transmitting end, and could not have worked as he described it.[111] Another inventor, Hovannes Adamian, also experimented with color television as early as 1907. The first color television project is claimed by him,[112] and was patented in Germany on 31 March 1908, patent No. 197183, then in Britain, on 1 April 1908, patent No. 7219,[113] in France (patent No. 390326) and in Russia in 1910 (patent No. 17912).[114]
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Scottish inventor John Logie Baird demonstrated the world's first color transmission on 3 July 1928, using scanning discs at the transmitting and receiving ends with three spirals of apertures, each spiral with filters of a different primary color; and three light sources at the receiving end, with a commutator to alternate their illumination.[115] Baird also made the world's first color broadcast on 4 February 1938, sending a mechanically scanned 120-line image from Baird's Crystal Palace studios to a projection screen at London's Dominion Theatre.[116] Mechanically scanned color television was also demonstrated by Bell Laboratories in June 1929 using three complete systems of photoelectric cells, amplifiers, glow-tubes, and color filters, with a series of mirrors to superimpose the red, green, and blue images into one full color image.
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The first practical hybrid system was again pioneered by John Logie Baird. In 1940 he publicly demonstrated a color television combining a traditional black-and-white display with a rotating colored disk. This device was very "deep", but was later improved with a mirror folding the light path into an entirely practical device resembling a large conventional console.[117] However, Baird was not happy with the design, and, as early as 1944, had commented to a British government committee that a fully electronic device would be better.
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In 1939, Hungarian engineer Peter Carl Goldmark introduced an electro-mechanical system while at CBS, which contained an Iconoscope sensor. The CBS field-sequential color system was partly mechanical, with a disc made of red, blue, and green filters spinning inside the television camera at 1,200 rpm, and a similar disc spinning in synchronization in front of the cathode ray tube inside the receiver set.[118] The system was first demonstrated to the Federal Communications Commission (FCC) on 29 August 1940, and shown to the press on 4 September.[119][120][121][122]
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CBS began experimental color field tests using film as early as 28 August 1940, and live cameras by 12 November.[120][123] NBC (owned by RCA) made its first field test of color television on 20 February 1941. CBS began daily color field tests on 1 June 1941.[124] These color systems were not compatible with existing black-and-white television sets, and, as no color television sets were available to the public at this time, viewing of the color field tests was restricted to RCA and CBS engineers and the invited press. The War Production Board halted the manufacture of television and radio equipment for civilian use from 22 April 1942 to 20 August 1945, limiting any opportunity to introduce color television to the general public.[125][126]
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As early as 1940, Baird had started work on a fully electronic system he called Telechrome. Early Telechrome devices used two electron guns aimed at either side of a phosphor plate. The phosphor was patterned so the electrons from the guns only fell on one side of the patterning or the other. Using cyan and magenta phosphors, a reasonable limited-color image could be obtained. He also demonstrated the same system using monochrome signals to produce a 3D image (called "stereoscopic" at the time). A demonstration on 16 August 1944 was the first example of a practical color television system. Work on the Telechrome continued and plans were made to introduce a three-gun version for full color. However, Baird's untimely death in 1946 ended development of the Telechrome system.[127][128]
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Similar concepts were common through the 1940s and 1950s, differing primarily in the way they re-combined the colors generated by the three guns. The Geer tube was similar to Baird's concept, but used small pyramids with the phosphors deposited on their outside faces, instead of Baird's 3D patterning on a flat surface. The Penetron used three layers of phosphor on top of each other and increased the power of the beam to reach the upper layers when drawing those colors. The Chromatron used a set of focusing wires to select the colored phosphors arranged in vertical stripes on the tube.
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One of the great technical challenges of introducing color broadcast television was the desire to conserve bandwidth, potentially three times that of the existing black-and-white standards, and not use an excessive amount of radio spectrum. In the United States, after considerable research, the National Television Systems Committee[129] approved an all-electronic system developed by RCA, which encoded the color information separately from the brightness information and greatly reduced the resolution of the color information in order to conserve bandwidth. As black-and-white TVs could receive the same transmission and display it in black-and-white, the color system adopted is [backwards] "compatible". ("Compatible Color", featured in RCA advertisements of the period, is mentioned in the song "America", of West Side Story, 1957.) The brightness image remained compatible with existing black-and-white television sets at slightly reduced resolution, while color televisions could decode the extra information in the signal and produce a limited-resolution color display. The higher resolution black-and-white and lower resolution color images combine in the brain to produce a seemingly high-resolution color image. The NTSC standard represented a major technical achievement.
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The first color broadcast (the first episode of the live program The Marriage (TV series)) occurred on 8 July 1954, but during the following ten years most network broadcasts, and nearly all local programming, continued to be in black-and-white. It was not until the mid-1960s that color sets started selling in large numbers, due in part to the color transition of 1965 in which it was announced that over half of all network prime-time programming would be broadcast in color that fall. The first all-color prime-time season came just one year later. In 1972, the last holdout among daytime network programs converted to color, resulting in the first completely all-color network season.
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Early color sets were either floor-standing console models or tabletop versions nearly as bulky and heavy, so in practice they remained firmly anchored in one place. GE's relatively compact and lightweight Porta-Color set was introduced in the spring of 1966. It used a transistor-based UHF tuner.[130] The first fully transistorized color television in the United States was the Quasar television introduced in 1967.[131] These developments made watching color television a more flexible and convenient proposition.
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The MOSFET (metal-oxide-semiconductor field-effect transistor, or MOS transistor) was invented by Mohamed M. Atalla and Dawon Kahng at Bell Labs in 1959,[132] and presented in 1960.[133] By the mid-1960s, RCA were using MOSFETs in their consumer television products.[134] RCA Laboratories researchers W.M. Austin, J.A. Dean, D.M. Griswold and O.P. Hart in 1966 described the use of the MOSFET in television circuits, including RF amplifier, low-level video, chroma and AGC circuits.[135] The power MOSFET was later widely adopted for television receiver circuits.[136]
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In 1972, sales of color sets finally surpassed sales of black-and-white sets. Color broadcasting in Europe was not standardized on the PAL format until the 1960s, and broadcasts did not start until 1967. By this point many of the technical problems in the early sets had been worked out, and the spread of color sets in Europe was fairly rapid. By the mid-1970s, the only stations broadcasting in black-and-white were a few high-numbered UHF stations in small markets, and a handful of low-power repeater stations in even smaller markets such as vacation spots. By 1979, even the last of these had converted to color and, by the early 1980s, B&W sets had been pushed into niche markets, notably low-power uses, small portable sets, or for use as video monitor screens in lower-cost consumer equipment. By the late 1980s even these areas switched to color sets.
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Digital television (DTV) is the transmission of audio and video by digitally processed and multiplexed signals, in contrast to the totally analog and channel separated signals used by analog television. Due to data compression, digital TV can support more than one program in the same channel bandwidth.[137] It is an innovative service that represents the most significant evolution in television broadcast technology since color television emerged in the 1950s.[138] Digital TV's roots have been tied very closely to the availability of inexpensive, high performance computers. It was not until the 1990s that digital TV became feasible.[139] Digital television was previously not practically feasible due to the impractically high bandwidth requirements of uncompressed digital video,[140][141] requiring around 200 Mbit/s bit-rate for a standard-definition television (SDTV) signal,[140] and over 1 Gbit/s for high-definition television (HDTV).[141]
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Digital TV became practically feasible in the early 1990s due to a major technological development, discrete cosine transform (DCT) video compression.[140][141] DCT coding is a lossy compression technique that was first proposed for image compression by Nasir Ahmed in 1972,[142] and was later adapted into a motion-compensated DCT video coding algorithm, for video coding standards such as the H.26x formats from 1988 onwards and the MPEG formats from 1991 onwards.[143][144] Motion-compensated DCT video compression significantly reduced the amount of bandwidth required for a digital TV signal.[140][141] DCT coding compressed down the bandwidth requirements of digital television signals to about 34 Mpps bit-rate for SDTV and around 70–140 Mbit/s for HDTV while maintaining near-studio-quality transmission, making digital television a practical reality in the 1990s.[141]
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A digital TV service was proposed in 1986 by Nippon Telegraph and Telephone (NTT) and the Ministry of Posts and Telecommunication (MPT) in Japan, where there were plans to develop an "Integrated Network System" service. However, it was not possible to practically implement such a digital TV service until the adoption of DCT video compression technology made it possible in the early 1990s.[140]
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In the mid-1980s, as Japanese consumer electronics firms forged ahead with the development of HDTV technology, the MUSE analog format proposed by NHK, a Japanese company, was seen as a pacesetter that threatened to eclipse U.S. electronics companies' technologies. Until June 1990, the Japanese MUSE standard, based on an analog system, was the front-runner among the more than 23 different technical concepts under consideration. Then, an American company, General Instrument, demonstrated the feasibility of a digital television signal. This breakthrough was of such significance that the FCC was persuaded to delay its decision on an ATV standard until a digitally based standard could be developed.
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In March 1990, when it became clear that a digital standard was feasible, the FCC made a number of critical decisions. First, the Commission declared that the new ATV standard must be more than an enhanced analog signal, but be able to provide a genuine HDTV signal with at least twice the resolution of existing television images.(7) Then, to ensure that viewers who did not wish to buy a new digital television set could continue to receive conventional television broadcasts, it dictated that the new ATV standard must be capable of being "simulcast" on different channels.(8)The new ATV standard also allowed the new DTV signal to be based on entirely new design principles. Although incompatible with the existing NTSC standard, the new DTV standard would be able to incorporate many improvements.
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The final standards adopted by the FCC did not require a single standard for scanning formats, aspect ratios, or lines of resolution. This compromise resulted from a dispute between the consumer electronics industry (joined by some broadcasters) and the computer industry (joined by the film industry and some public interest groups) over which of the two scanning processes—interlaced or progressive—would be best suited for the newer digital HDTV compatible display devices.[145] Interlaced scanning, which had been specifically designed for older analogue CRT display technologies, scans even-numbered lines first, then odd-numbered ones. In fact, interlaced scanning can be looked at as the first video compression model as it was partly designed in the 1940s to double the image resolution to exceed the limitations of the television broadcast bandwidth. Another reason for its adoption was to limit the flickering on early CRT screens whose phosphor coated screens could only retain the image from the electron scanning gun for a relatively short duration.[146] However interlaced scanning does not work as efficiently on newer display devices such as Liquid-crystal (LCD), for example, which are better suited to a more frequent progressive refresh rate.[145]
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Progressive scanning, the format that the computer industry had long adopted for computer display monitors, scans every line in sequence, from top to bottom. Progressive scanning in effect doubles the amount of data generated for every full screen displayed in comparison to interlaced scanning by painting the screen in one pass in 1/60-second, instead of two passes in 1/30-second. The computer industry argued that progressive scanning is superior because it does not "flicker" on the new standard of display devices in the manner of interlaced scanning. It also argued that progressive scanning enables easier connections with the Internet, and is more cheaply converted to interlaced formats than vice versa. The film industry also supported progressive scanning because it offered a more efficient means of converting filmed programming into digital formats. For their part, the consumer electronics industry and broadcasters argued that interlaced scanning was the only technology that could transmit the highest quality pictures then (and currently) feasible, i.e., 1,080 lines per picture and 1,920 pixels per line. Broadcasters also favored interlaced scanning because their vast archive of interlaced programming is not readily compatible with a progressive format. William F. Schreiber, who was director of the Advanced Television Research Program at the Massachusetts Institute of Technology from 1983 until his retirement in 1990, thought that the continued advocacy of interlaced equipment originated from consumer electronics companies that were trying to get back the substantial investments they made in the interlaced technology.[147]
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Digital television transition started in late 2000s. All governments across the world set the deadline for analog shutdown by 2010s. Initially the adoption rate was low, as the first digital tuner-equipped TVs were costly. But soon, as the price of digital-capable TVs dropped, more and more households were converting to digital televisions. The transition is expected to be completed worldwide by mid to late 2010s.
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The advent of digital television allowed innovations like smart TVs. A smart television, sometimes referred to as connected TV or hybrid TV, is a television set or set-top box with integrated Internet and Web 2.0 features, and is an example of technological convergence between computers, television sets and set-top boxes. Besides the traditional functions of television sets and set-top boxes provided through traditional Broadcasting media, these devices can also provide Internet TV, online interactive media, over-the-top content, as well as on-demand streaming media, and home networking access. These TVs come pre-loaded with an operating system.[9][148][149][150]
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Smart TV should not to be confused with Internet TV, Internet Protocol television (IPTV) or with Web TV. Internet television refers to the receiving of television content over the Internet instead of by traditional systems—terrestrial, cable and satellite (although internet itself is received by these methods). IPTV is one of the emerging Internet television technology standards for use by television broadcasters. Web television (WebTV) is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV. A first patent was filed in 1994[151] (and extended the following year)[152] for an "intelligent" television system, linked with data processing systems, by means of a digital or analog network. Apart from being linked to data networks, one key point is its ability to automatically download necessary software routines, according to a user's demand, and process their needs. Major TV manufacturers have announced production of smart TVs only, for middle-end and high-end TVs in 2015.[6][7][8] Smart TVs have gotten more affordable compared to when they were first introduced, with 46 million of U.S. households having at least one as of 2019.[153]
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3D television conveys depth perception to the viewer by employing techniques such as stereoscopic display, multi-view display, 2D-plus-depth, or any other form of 3D display. Most modern 3D television sets use an active shutter 3D system or a polarized 3D system, and some are autostereoscopic without the need of glasses. Stereoscopic 3D television was demonstrated for the first time on 10 August 1928, by John Logie Baird in his company's premises at 133 Long Acre, London.[154] Baird pioneered a variety of 3D television systems using electromechanical and cathode-ray tube techniques. The first 3D TV was produced in 1935. The advent of digital television in the 2000s greatly improved 3D TVs. Although 3D TV sets are quite popular for watching 3D home media such as on Blu-ray discs, 3D programming has largely failed to make inroads with the public. Many 3D television channels which started in the early 2010s were shut down by the mid-2010s. According to DisplaySearch 3D televisions shipments totaled 41.45 million units in 2012, compared with 24.14 in 2011 and 2.26 in 2010.[155] As of late 2013, the number of 3D TV viewers started to decline.[156][157][158][159][160]
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Programming is broadcast by television stations, sometimes called "channels", as stations are licensed by their governments to broadcast only over assigned channels in the television band. At first, terrestrial broadcasting was the only way television could be widely distributed, and because bandwidth was limited, i.e., there were only a small number of channels available, government regulation was the norm. In the U.S., the Federal Communications Commission (FCC) allowed stations to broadcast advertisements beginning in July 1941, but required public service programming commitments as a requirement for a license. By contrast, the United Kingdom chose a different route, imposing a television license fee on owners of television reception equipment to fund the British Broadcasting Corporation (BBC), which had public service as part of its Royal Charter.
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WRGB claims to be the world's oldest television station, tracing its roots to an experimental station founded on 13 January 1928, broadcasting from the General Electric factory in Schenectady, NY, under the call letters W2XB.[161] It was popularly known as "WGY Television" after its sister radio station. Later in 1928, General Electric started a second facility, this one in New York City, which had the call letters W2XBS and which today is known as WNBC. The two stations were experimental in nature and had no regular programming, as receivers were operated by engineers within the company. The image of a Felix the Cat doll rotating on a turntable was broadcast for 2 hours every day for several years as new technology was being tested by the engineers. On 2 November 1936, the BBC began transmitting the world's first public regular high-definition service from the Victorian Alexandra Palace in north London.[162] It therefore claims to be the birthplace of TV broadcasting as we know it today.
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With the widespread adoption of cable across the United States in the 1970s and 80s, terrestrial television broadcasts have been in decline; in 2013 it was estimated that about 7% of US households used an antenna.[163][164] A slight increase in use began around 2010 due to switchover to digital terrestrial television broadcasts, which offered pristine image quality over very large areas, and offered an alternate to cable television (CATV) for cord cutters. All other countries around the world are also in the process of either shutting down analog terrestrial television or switching over to digital terrestrial television.
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Cable television is a system of broadcasting television programming to paying subscribers via radio frequency (RF) signals transmitted through coaxial cables or light pulses through fiber-optic cables. This contrasts with traditional terrestrial television, in which the television signal is transmitted over the air by radio waves and received by a television antenna attached to the television. In the 2000s, FM radio programming, high-speed Internet, telephone service, and similar non-television services may also be provided through these cables. The abbreviation CATV is often used for cable television. It originally stood for Community Access Television or Community Antenna Television, from cable television's origins in 1948: in areas where over-the-air reception was limited by distance from transmitters or mountainous terrain, large "community antennas" were constructed, and cable was run from them to individual homes.[165] The origins of cable broadcasting are even older as radio programming was distributed by cable in some European cities as far back as 1924. Earlier cable television was analog, but since the 2000s, all cable operators have switched to, or are in the process of switching to, digital cable television.
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Satellite television is a system of supplying television programming using broadcast signals relayed from communication satellites. The signals are received via an outdoor parabolic reflector antenna usually referred to as a satellite dish and a low-noise block downconverter (LNB). A satellite receiver then decodes the desired television program for viewing on a television set. Receivers can be external set-top boxes, or a built-in television tuner. Satellite television provides a wide range of channels and services, especially to geographic areas without terrestrial television or cable television.
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The most common method of reception is direct-broadcast satellite television (DBSTV), also known as "direct to home" (DTH).[166] In DBSTV systems, signals are relayed from a direct broadcast satellite on the Ku wavelength and are completely digital.[167] Satellite TV systems formerly used systems known as television receive-only. These systems received analog signals transmitted in the C-band spectrum from FSS type satellites, and required the use of large dishes. Consequently, these systems were nicknamed "big dish" systems, and were more expensive and less popular.[168]
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The direct-broadcast satellite television signals were earlier analog signals and later digital signals, both of which require a compatible receiver. Digital signals may include high-definition television (HDTV). Some transmissions and channels are free-to-air or free-to-view, while many other channels are pay television requiring a subscription.[169]
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In 1945, British science fiction writer Arthur C. Clarke proposed a worldwide communications system which would function by means of three satellites equally spaced apart in earth orbit.[170][171] This was published in the October 1945 issue of the Wireless World magazine and won him the Franklin Institute's Stuart Ballantine Medal in 1963.[172][173]
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The first satellite television signals from Europe to North America were relayed via the Telstar satellite over the Atlantic ocean on 23 July 1962.[174] The signals were received and broadcast in North American and European countries and watched by over 100 million.[174] Launched in 1962, the Relay 1 satellite was the first satellite to transmit television signals from the US to Japan.[175] The first geosynchronous communication satellite, Syncom 2, was launched on 26 July 1963.[176]
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The world's first commercial communications satellite, called Intelsat I and nicknamed "Early Bird", was launched into geosynchronous orbit on 6 April 1965.[177] The first national network of television satellites, called Orbita, was created by the Soviet Union in October 1967, and was based on the principle of using the highly elliptical Molniya satellite for rebroadcasting and delivering of television signals to ground downlink stations.[178] The first commercial North American satellite to carry television transmissions was Canada's geostationary Anik 1, which was launched on 9 November 1972.[179] ATS-6, the world's first experimental educational and Direct Broadcast Satellite (DBS), was launched on 30 May 1974.[180] It transmitted at 860 MHz using wideband FM modulation and had two sound channels. The transmissions were focused on the Indian subcontinent but experimenters were able to receive the signal in Western Europe using home constructed equipment that drew on UHF television design techniques already in use.[181]
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The first in a series of Soviet geostationary satellites to carry Direct-To-Home television, Ekran 1, was launched on 26 October 1976.[182] It used a 714 MHz UHF downlink frequency so that the transmissions could be received with existing UHF television technology rather than microwave technology.[183]
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Internet television (Internet TV) (or online television) is the digital distribution of television content via the Internet as opposed to traditional systems like terrestrial, cable, and satellite, although the Internet itself is received by terrestrial, cable, or satellite methods. Internet television is a general term that covers the delivery of television shows, and other video content, over the Internet by video streaming technology, typically by major traditional television broadcasters. Internet television should not be confused with Smart TV, IPTV or with Web TV. Smart television refers to the TV set which has a built-in operating system. Internet Protocol television (IPTV) is one of the emerging Internet television technology standards for use by television broadcasters. Web television is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV.
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A television set, also called a television receiver, television, TV set, TV, or "telly", is a device that combines a tuner, display, an amplifier, and speakers for the purpose of viewing television and hearing its audio components. Introduced in the late 1920s in mechanical form, television sets became a popular consumer product after World War II in electronic form, using cathode ray tubes. The addition of color to broadcast television after 1953 further increased the popularity of television sets and an outdoor antenna became a common feature of suburban homes. The ubiquitous television set became the display device for recorded media in the 1970s, such as Betamax and VHS, which enabled viewers to record TV shows and watch prerecorded movies. In the subsequent decades, TVs were used to watch DVDs and Blu-ray Discs of movies and other content. Major TV manufacturers announced the discontinuation of CRT, DLP, plasma and fluorescent-backlit LCDs by the mid-2010s. Televisions since 2010s mostly use LEDs.[3][4][184][185] LEDs are expected to be gradually replaced by OLEDs in the near future.[5]
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The earliest systems employed a spinning disk to create and reproduce images.[186] These usually had a low resolution and screen size and never became popular with the public.
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The cathode ray tube (CRT) is a vacuum tube containing one or more electron guns (a source of electrons or electron emitter) and a fluorescent screen used to view images.[33] It has a means to accelerate and deflect the electron beam(s) onto the screen to create the images. The images may represent electrical waveforms (oscilloscope), pictures (television, computer monitor), radar targets or others. The CRT uses an evacuated glass envelope which is large, deep (i.e. long from front screen face to rear end), fairly heavy, and relatively fragile. As a matter of safety, the face is typically made of thick lead glass so as to be highly shatter-resistant and to block most X-ray emissions, particularly if the CRT is used in a consumer product.
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In television sets and computer monitors, the entire front area of the tube is scanned repetitively and systematically in a fixed pattern called a raster. An image is produced by controlling the intensity of each of the three electron beams, one for each additive primary color (red, green, and blue) with a video signal as a reference.[187] In all modern CRT monitors and televisions, the beams are bent by magnetic deflection, a varying magnetic field generated by coils and driven by electronic circuits around the neck of the tube, although electrostatic deflection is commonly used in oscilloscopes, a type of diagnostic instrument.[187]
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Digital Light Processing (DLP) is a type of video projector technology that uses a digital micromirror device. Some DLPs have a TV tuner, which makes them a type of TV display. It was originally developed in 1987 by Dr. Larry Hornbeck of Texas Instruments. While the DLP imaging device was invented by Texas Instruments, the first DLP based projector was introduced by Digital Projection Ltd in 1997. Digital Projection and Texas Instruments were both awarded Emmy Awards in 1998 for invention of the DLP projector technology. DLP is used in a variety of display applications from traditional static displays to interactive displays and also non-traditional embedded applications including medical, security, and industrial uses. DLP technology is used in DLP front projectors (standalone projection units for classrooms and business primarily), but also in private homes; in these cases, the image is projected onto a projection screen. DLP is also used in DLP rear projection television sets and digital signs. It is also used in about 85% of digital cinema projection.[188]
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A plasma display panel (PDP) is a type of flat panel display common to large TV displays 30 inches (76 cm) or larger. They are called "plasma" displays because the technology utilizes small cells containing electrically charged ionized gases, or what are in essence chambers more commonly known as fluorescent lamps.
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Liquid-crystal-display televisions (LCD TV) are television sets that use LCD display technology to produce images. LCD televisions are much thinner and lighter than cathode ray tube (CRTs) of similar display size, and are available in much larger sizes (e.g., 90-inch diagonal). When manufacturing costs fell, this combination of features made LCDs practical for television receivers. LCDs come in two types: those using cold cathode fluorescent lamps, simply called LCDs and those using LED as backlight called as LEDs.
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In 2007, LCD televisions surpassed sales of CRT-based televisions worldwide for the first time, and their sales figures relative to other technologies accelerated. LCD TVs have quickly displaced the only major competitors in the large-screen market, the Plasma display panel and rear-projection television.[189] In mid 2010s LCDs especially LEDs became, by far, the most widely produced and sold television display type.[184][185] LCDs also have disadvantages. Other technologies address these weaknesses, including OLEDs, FED and SED, but as of 2014[update] none of these have entered widespread production.
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An OLED (organic light-emitting diode) is a light-emitting diode (LED) in which the emissive electroluminescent layer is a film of organic compound which emits light in response to an electric current. This layer of organic semiconductor is situated between two electrodes. Generally, at least one of these electrodes is transparent. OLEDs are used to create digital displays in devices such as television screens. It is also used for computer monitors, portable systems such as mobile phones, handheld game consoles and PDAs.
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There are two main families of OLED: those based on small molecules and those employing polymers. Adding mobile ions to an OLED creates a light-emitting electrochemical cell or LEC, which has a slightly different mode of operation. OLED displays can use either passive-matrix (PMOLED) or active-matrix (AMOLED) addressing schemes. Active-matrix OLEDs require a thin-film transistor backplane to switch each individual pixel on or off, but allow for higher resolution and larger display sizes.
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An OLED display works without a backlight. Thus, it can display deep black levels and can be thinner and lighter than a liquid crystal display (LCD). In low ambient light conditions such as a dark room an OLED screen can achieve a higher contrast ratio than an LCD, whether the LCD uses cold cathode fluorescent lamps or LED backlight. OLEDs are expected to replace other forms of display in near future.[5]
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Low-definition television or LDTV refers to television systems that have a lower screen resolution than standard-definition television systems such 240p (320*240). It is used in handheld television. The most common source of LDTV programming is the Internet, where mass distribution of higher-resolution video files could overwhelm computer servers and take too long to download. Many mobile phones and portable devices such as Apple's iPod Nano, or Sony's PlayStation Portable use LDTV video, as higher-resolution files would be excessive to the needs of their small screens (320×240 and 480×272 pixels respectively). The current generation of iPod Nanos have LDTV screens, as do the first three generations of iPod Touch and iPhone (480×320). For the first years of its existence, YouTube offered only one, low-definition resolution of 320x240p at 30fps or less. A standard, consumer grade VHS videotape can be considered SDTV due to its resolution (approximately 360 × 480i/576i).
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Standard-definition television or SDTV refers to two different resolutions: 576i, with 576 interlaced lines of resolution, derived from the European-developed PAL and SECAM systems; and 480i based on the American National Television System Committee NTSC system. SDTV is a television system that uses a resolution that is not considered to be either high-definition television (720p, 1080i, 1080p, 1440p, 4K UHDTV, and 8K UHD) or enhanced-definition television (EDTV 480p). In North America, digital SDTV is broadcast in the same 4:3 aspect ratio as NTSC signals with widescreen content being center cut.[190] However, in other parts of the world that used the PAL or SECAM color systems, standard-definition television is now usually shown with a 16:9 aspect ratio, with the transition occurring between the mid-1990s and mid-2000s. Older programs with a 4:3 aspect ratio are shown in the US as 4:3 with non-ATSC countries preferring to reduce the horizontal resolution by anamorphically scaling a pillarboxed image.
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High-definition television (HDTV) provides a resolution that is substantially higher than that of standard-definition television.
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HDTV may be transmitted in various formats:
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Ultra-high-definition television (also known as Super Hi-Vision, Ultra HD television, UltraHD, UHDTV, or UHD) includes 4K UHD (2160p) and 8K UHD (4320p), which are two digital video formats proposed by NHK Science & Technology Research Laboratories and defined and approved by the International Telecommunication Union (ITU). The Consumer Electronics Association announced on 17 October 2012, that "Ultra High Definition", or "Ultra HD", would be used for displays that have an aspect ratio of at least 16:9 and at least one digital input capable of carrying and presenting native video at a minimum resolution of 3840×2160 pixels.[191][192]
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North American consumers purchase a new television set on average every seven years, and the average household owns 2.8 televisions. As of 2011[update], 48 million are sold each year at an average price of $460 and size of 38 in (97 cm).[193]
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Getting TV programming shown to the public can happen in many different ways. After production, the next step is to market and deliver the product to whichever markets are open to using it. This typically happens on two levels:
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First-run programming is increasing on subscription services outside the US, but few domestically produced programs are syndicated on domestic free-to-air (FTA) elsewhere. This practice is increasing, however, generally on digital-only FTA channels or with subscriber-only, first-run material appearing on FTA. Unlike the US, repeat FTA screenings of an FTA network program usually only occur on that network. Also, affiliates rarely buy or produce non-network programming that is not centered on local programming.
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Television genres include a broad range of programming types that entertain, inform, and educate viewers. The most expensive entertainment genres to produce are usually dramas and dramatic miniseries. However, other genres, such as historical Western genres, may also have high production costs.
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Popular culture entertainment genres include action-oriented shows such as police, crime, detective dramas, horror, or thriller shows. As well, there are also other variants of the drama genre, such as medical dramas and daytime soap operas. Science fiction shows can fall into either the drama or action category, depending on whether they emphasize philosophical questions or high adventure. Comedy is a popular genre which includes situation comedy (sitcom) and animated shows for the adult demographic such as South Park.
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The least expensive forms of entertainment programming genres are game shows, talk shows, variety shows, and reality television. Game shows feature contestants answering questions and solving puzzles to win prizes. Talk shows contain interviews with film, television, music and sports celebrities and public figures. Variety shows feature a range of musical performers and other entertainers, such as comedians and magicians, introduced by a host or Master of Ceremonies. There is some crossover between some talk shows and variety shows because leading talk shows often feature performances by bands, singers, comedians, and other performers in between the interview segments. Reality TV shows "regular" people (i.e., not actors) facing unusual challenges or experiences ranging from arrest by police officers (COPS) to significant weight loss (The Biggest Loser). A variant version of reality shows depicts celebrities doing mundane activities such as going about their everyday life (The Osbournes, Snoop Dogg's Father Hood) or doing regular jobs (The Simple Life).
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Fictional television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include series such as Twin Peaks and The Sopranos. Kristin Thompson argues that some of these television series exhibit traits also found in art films, such as psychological realism, narrative complexity, and ambiguous plotlines. Nonfiction television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include a range of serious, noncommercial, programming aimed at a niche audience, such as documentaries and public affairs shows.
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Around the globe, broadcast TV is financed by government, advertising, licensing (a form of tax), subscription, or any combination of these. To protect revenues, subscription TV channels are usually encrypted to ensure that only subscribers receive the decryption codes to see the signal. Unencrypted channels are known as free to air or FTA. In 2009, the global TV market represented 1,217.2 million TV households with at least one TV and total revenues of 268.9 billion EUR (declining 1.2% compared to 2008).[195] North America had the biggest TV revenue market share with 39% followed by Europe (31%), Asia-Pacific (21%), Latin America (8%), and Africa and the Middle East (2%).[196] Globally, the different TV revenue sources divide into 45–50% TV advertising revenues, 40–45% subscription fees and 10% public funding.[197][198]
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TV's broad reach makes it a powerful and attractive medium for advertisers. Many TV networks and stations sell blocks of broadcast time to advertisers ("sponsors") to fund their programming.[199] Television advertisements (variously called a television commercial, commercial or ad in American English, and known in British English as an advert) is a span of television programming produced and paid for by an organization, which conveys a message, typically to market a product or service. Advertising revenue provides a significant portion of the funding for most privately owned television networks. The vast majority of television advertisements today consist of brief advertising spots, ranging in length from a few seconds to several minutes (as well as program-length infomercials). Advertisements of this sort have been used to promote a wide variety of goods, services and ideas since the beginning of television.
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The effects of television advertising upon the viewing public (and the effects of mass media in general) have been the subject of philosophical discourse by such luminaries as Marshall McLuhan. The viewership of television programming, as measured by companies such as Nielsen Media Research, is often used as a metric for television advertisement placement, and consequently, for the rates charged to advertisers to air within a given network, television program, or time of day (called a "daypart"). In many countries, including the United States, television campaign advertisements are considered indispensable for a political campaign. In other countries, such as France, political advertising on television is heavily restricted,[200] while some countries, such as Norway, completely ban political advertisements.
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The first official, paid television advertisement was broadcast in the United States on 1 July 1941 over New York station WNBT (now WNBC) before a baseball game between the Brooklyn Dodgers and Philadelphia Phillies. The announcement for Bulova watches, for which the company paid anywhere from $4.00 to $9.00 (reports vary), displayed a WNBT test pattern modified to look like a clock with the hands showing the time. The Bulova logo, with the phrase "Bulova Watch Time", was shown in the lower right-hand quadrant of the test pattern while the second hand swept around the dial for one minute.[201][202] The first TV ad broadcast in the UK was on ITV on 22 September 1955, advertising Gibbs SR toothpaste. The first TV ad broadcast in Asia was on Nippon Television in Tokyo on 28 August 1953, advertising Seikosha (now Seiko), which also displayed a clock with the current time.[203]
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Since inception in the US in 1941,[204] television commercials have become one of the most effective, persuasive, and popular methods of selling products of many sorts, especially consumer goods. During the 1940s and into the 1950s, programs were hosted by single advertisers. This, in turn, gave great creative license to the advertisers over the content of the show. Perhaps due to the quiz show scandals in the 1950s,[205] networks shifted to the magazine concept, introducing advertising breaks with multiple advertisers.
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US advertising rates are determined primarily by Nielsen ratings. The time of the day and popularity of the channel determine how much a TV commercial can cost. For example, it can cost approximately $750,000 for a 30-second block of commercial time during the highly popular American Idol, while the same amount of time for the Super Bowl can cost several million dollars. Conversely, lesser-viewed time slots, such as early mornings and weekday afternoons, are often sold in bulk to producers of infomercials at far lower rates. In recent years, the paid program or infomercial has become common, usually in lengths of 30 minutes or one hour. Some drug companies and other businesses have even created "news" items for broadcast, known in the industry as video news releases, paying program directors to use them.[206]
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Some TV programs also deliberately place products into their shows as advertisements, a practice started in feature films[207] and known as product placement. For example, a character could be drinking a certain kind of soda, going to a particular chain restaurant, or driving a certain make of car. (This is sometimes very subtle, with shows having vehicles provided by manufacturers for low cost in exchange as a product placement). Sometimes, a specific brand or trade mark, or music from a certain artist or group, is used. (This excludes guest appearances by artists who perform on the show.)
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The TV regulator oversees TV advertising in the United Kingdom. Its restrictions have applied since the early days of commercially funded TV. Despite this, an early TV mogul, Roy Thomson, likened the broadcasting licence as being a "licence to print money".[208] Restrictions mean that the big three national commercial TV channels: ITV, Channel 4, and Channel 5 can show an average of only seven minutes of advertising per hour (eight minutes in the peak period). Other broadcasters must average no more than nine minutes (twelve in the peak). This means that many imported TV shows from the US have unnatural pauses where the UK company does not utilize the narrative breaks intended for more frequent US advertising. Advertisements must not be inserted in the course of certain specific proscribed types of programs which last less than half an hour in scheduled duration; this list includes any news or current affairs programs, documentaries, and programs for children; additionally, advertisements may not be carried in a program designed and broadcast for reception in schools or in any religious broadcasting service or other devotional program or during a formal Royal ceremony or occasion. There also must be clear demarcations in time between the programs and the advertisements. The BBC, being strictly non-commercial, is not allowed to show advertisements on television in the UK, although it has many advertising-funded channels abroad. The majority of its budget comes from television license fees (see below) and broadcast syndication, the sale of content to other broadcasters.
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Broadcast advertising is regulated by the Broadcasting Authority of Ireland.[209]
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Some TV channels are partly funded from subscriptions; therefore, the signals are encrypted during broadcast to ensure that only the paying subscribers have access to the decryption codes to watch pay television or specialty channels. Most subscription services are also funded by advertising.
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Television services in some countries may be funded by a television licence or a form of taxation, which means that advertising plays a lesser role or no role at all. For example, some channels may carry no advertising at all and some very little, including:
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The BBC carries no television advertising on its UK channels and is funded by an annual television licence paid by premises receiving live TV broadcasts. Currently, it is estimated that approximately 26.8 million UK private domestic households own televisions, with approximately 25 million TV licences in all premises in force as of 2010.[210] This television license fee is set by the government, but the BBC is not answerable to or controlled by the government.
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The two main BBC TV channels are watched by almost 90% of the population each week and overall have 27% share of total viewing,[211] despite the fact that 85% of homes are multichannel, with 42% of these having access to 200 free to air channels via satellite and another 43% having access to 30 or more channels via Freeview.[212] The licence that funds the seven advertising-free BBC TV channels costs £147 a year (about US$200) as of 2018 regardless of the number of TV sets owned; the price is reduced by two-thirds if only black and white television is received.[213] When the same sporting event has been presented on both BBC and commercial channels, the BBC always attracts the lion's share of the audience, indicating that viewers prefer to watch TV uninterrupted by advertising.
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Other than internal promotional material, the Australian Broadcasting Corporation (ABC) carries no advertising; it is banned under the ABC Act 1983. The ABC receives its funding from the Australian government every three years. In the 2014/15 federal budget, the ABC received A$1.11 billion.[214] The funds provide for the ABC's television, radio, online, and international outputs. The ABC also receives funds from its many ABC shops across Australia. Although funded by the Australian government, the editorial independence of the ABC is ensured through law.
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In France, government-funded channels carry advertisements, yet those who own television sets have to pay an annual tax ("la redevance audiovisuelle").[215]
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In Japan, NHK is paid for by license fees (known in Japanese as reception fee (受信料, Jushinryō)). The broadcast law that governs NHK's funding stipulates that any television equipped to receive NHK is required to pay. The fee is standardized, with discounts for office workers and students who commute, as well a general discount for residents of Okinawa prefecture.
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Broadcast programming, or TV listings in the United Kingdom, is the practice of organizing television programs in a schedule, with broadcast automation used to regularly change the scheduling of TV programs to build an audience for a new show, retain that audience, or compete with other broadcasters' programs.
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Television has played a pivotal role in the socialization of the 20th and 21st centuries. There are many aspects of television that can be addressed, including negative issues such as media violence. Current research is discovering that individuals suffering from social isolation can employ television to create what is termed a parasocial or faux relationship with characters from their favorite television shows and movies as a way of deflecting feelings of loneliness and social deprivation.[216] Several studies have found that educational television has many advantages. The article "The Good Things about Television"[217] argues that television can be a very powerful and effective learning tool for children if used wisely.
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Methodist denominations in the conservative holiness movement, such as the Allegheny Wesleyan Methodist Connection and the Evangelical Wesleyan Church, eschew the use of the television.[218]
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Children, especially those aged 5 or younger, are at risk of injury from falling televisions.[219] A CRT-style television that falls on a child will, because of its weight, hit with the equivalent force of falling multiple stories from a building.[220] Newer flat-screen televisions are "top-heavy and have narrow bases", which means that a small child can easily pull one over.[221] As of 2015[update], TV tip-overs were responsible for more than 10,000 injuries per year to children, at a cost of more than $8 million per year in emergency care.[219][221]
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A 2017 study in The Journal of Human Resources found that exposure to cable television reduced cognitive ability and high school graduation rates for boys. This effect was stronger for boys from more educated families. The article suggests a mechanism where light television entertainment crowds out more cognitively stimulating activities.[222]
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With high lead content in CRTs and the rapid diffusion of new flat-panel display technologies, some of which (LCDs) use lamps which contain mercury, there is growing concern about electronic waste from discarded televisions. Related occupational health concerns exist, as well, for disassemblers removing copper wiring and other materials from CRTs. Further environmental concerns related to television design and use relate to the devices' increasing electrical energy requirements.[223]
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Television (TV), sometimes shortened to tele or telly, is a telecommunication medium used for transmitting moving images in monochrome (black and white), or in color, and in two or three dimensions and sound. The term can refer to a television set, a television show, or the medium of television transmission. Television is a mass medium for advertising, entertainment, news, and sports.
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Television became available in crude experimental forms in the late 1920s, but it would still be several years before the new technology would be marketed to consumers. After World War II, an improved form of black-and-white TV broadcasting became popular in the United States and Britain, and television sets became commonplace in homes, businesses, and institutions. During the 1950s, television was the primary medium for influencing public opinion.[1] In the mid-1960s, color broadcasting was introduced in the US and most other developed countries. The availability of multiple types of archival storage media such as Betamax and VHS tapes, high-capacity hard disk drives, DVDs, flash drives, high-definition Blu-ray Discs, and cloud digital video recorders has enabled viewers to watch pre-recorded material—such as movies—at home on their own time schedule. For many reasons, especially the convenience of remote retrieval, the storage of television and video programming now occurs on the cloud. At the end of the first decade of the 2000s, digital television transmissions greatly increased in popularity. Another development was the move from standard-definition television (SDTV) (576i, with 576 interlaced lines of resolution and 480i) to high-definition television (HDTV), which provides a resolution that is substantially higher. HDTV may be transmitted in various formats: 1080p, 1080i and 720p. Since 2010, with the invention of smart television, Internet television has increased the availability of television programs and movies via the Internet through streaming video services such as Netflix, Amazon Video, iPlayer and Hulu.
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In 2013, 79% of the world's households owned a television set.[2] The replacement of early bulky, high-voltage cathode ray tube (CRT) screen displays with compact, energy-efficient, flat-panel alternative technologies such as LCDs (both fluorescent-backlit and LED), OLED displays, and plasma displays was a hardware revolution that began with computer monitors in the late 1990s. Most TV sets sold in the 2000s were flat-panel, mainly LEDs. Major manufacturers announced the discontinuation of CRT, DLP, plasma, and even fluorescent-backlit LCDs by the mid-2010s.[3][4] In the near future, LEDs are expected to be gradually replaced by OLEDs.[5] Also, major manufacturers have announced that they will increasingly produce smart TVs in the mid-2010s.[6][7][8] Smart TVs with integrated Internet and Web 2.0 functions became the dominant form of television by the late 2010s.[9]
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Television signals were initially distributed only as terrestrial television using high-powered radio-frequency transmitters to broadcast the signal to individual television receivers. Alternatively television signals are distributed by coaxial cable or optical fiber, satellite systems and, since the 2000s via the Internet. Until the early 2000s, these were transmitted as analog signals, but a transition to digital television is expected to be completed worldwide by the late 2010s. A standard television set is composed of multiple internal electronic circuits, including a tuner for receiving and decoding broadcast signals. A visual display device which lacks a tuner is correctly called a video monitor rather than a television.
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The word television comes from Ancient Greek τῆλε (tèle), meaning 'far', and Latin visio, meaning 'sight'. The first documented usage of the term dates back to 1900, when the Russian scientist Constantin Perskyi used it in a paper that he presented in French at the 1st International Congress of Electricity, which ran from 18 to 25 August 1900 during the International World Fair in Paris. The Anglicised version of the term is first attested in 1907, when it was still "...a theoretical system to transmit moving images over telegraph or telephone wires".[10] It was "...formed in English or borrowed from French télévision."[10] In the 19th century and early 20th century, other "...proposals for the name of a then-hypothetical technology for sending pictures over distance were telephote (1880) and televista (1904)."[10] The abbreviation "TV" is from 1948. The use of the term to mean "a television set" dates from 1941.[10] The use of the term to mean "television as a medium" dates from 1927.[10] The slang term "telly" is more common in the UK. The slang term "the tube" or the "boob tube" derives from the bulky cathode ray tube used on most TVs until the advent of flat-screen TVs. Another slang term for the TV is "idiot box".[11] Also, in the 1940s and throughout the 1950s, during the early rapid growth of television programming and television-set ownership in the United States, another slang term became widely used in that period and continues to be used today to distinguish productions originally created for broadcast on television from films developed for presentation in movie theaters.[12] The "small screen", as both a compound adjective and noun, became specific references to television, while the "big screen" was used to identify productions made for theatrical release.[12]
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Facsimile transmission systems for still photographs pioneered methods of mechanical scanning of images in the early 19th century. Alexander Bain introduced the facsimile machine between 1843 and 1846. Frederick Bakewell demonstrated a working laboratory version in 1851.[citation needed] Willoughby Smith discovered the photoconductivity of the element selenium in 1873. As a 23-year-old German university student, Paul Julius Gottlieb Nipkow proposed and patented the Nipkow disk in 1884.[13] This was a spinning disk with a spiral pattern of holes in it, so each hole scanned a line of the image. Although he never built a working model of the system, variations of Nipkow's spinning-disk "image rasterizer" became exceedingly common.[14] Constantin Perskyi had coined the word television in a paper read to the International Electricity Congress at the International World Fair in Paris on 24 August 1900. Perskyi's paper reviewed the existing electromechanical technologies, mentioning the work of Nipkow and others.[15] However, it was not until 1907 that developments in amplification tube technology by Lee de Forest and Arthur Korn, among others, made the design practical.[16]
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The first demonstration of the live transmission of images was by Georges Rignoux and A. Fournier in Paris in 1909. A matrix of 64 selenium cells, individually wired to a mechanical commutator, served as an electronic retina. In the receiver, a type of Kerr cell modulated the light and a series of variously angled mirrors attached to the edge of a rotating disc scanned the modulated beam onto the display screen. A separate circuit regulated synchronization. The 8x8 pixel resolution in this proof-of-concept demonstration was just sufficient to clearly transmit individual letters of the alphabet. An updated image was transmitted "several times" each second.[17]
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In 1911, Boris Rosing and his student Vladimir Zworykin created a system that used a mechanical mirror-drum scanner to transmit, in Zworykin's words, "very crude images" over wires to the "Braun tube" (cathode ray tube or "CRT") in the receiver. Moving images were not possible because, in the scanner: "the sensitivity was not enough and the selenium cell was very laggy".[18]
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In 1921, Edouard Belin sent the first image via radio waves with his belinograph.[citation needed]
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By the 1920s, when amplification made television practical, Scottish inventor John Logie Baird employed the Nipkow disk in his prototype video systems. On 25 March 1925, Baird gave the first public demonstration of televised silhouette images in motion, at Selfridge's Department Store in London.[19] Since human faces had inadequate contrast to show up on his primitive system, he televised a ventriloquist's dummy named "Stooky Bill", whose painted face had higher contrast, talking and moving. By 26 January 1926, he demonstrated the transmission of the image of a face in motion by radio. This is widely regarded as the first television demonstration. The subject was Baird's business partner Oliver Hutchinson. Baird's system used the Nipkow disk for both scanning the image and displaying it. A bright light shining through a spinning Nipkow disk set with lenses projected a bright spot of light which swept across the subject. A Selenium photoelectric tube detected the light reflected from the subject and converted it into a proportional electrical signal. This was transmitted by AM radio waves to a receiver unit, where the video signal was applied to a neon light behind a second Nipkow disk rotating synchronized with the first. The brightness of the neon lamp was varied in proportion to the brightness of each spot on the image. As each hole in the disk passed by, one scan line of the image was reproduced. Baird's disk had 30 holes, producing an image with only 30 scan lines, just enough to recognize a human face. In 1927, Baird transmitted a signal over 438 miles (705 km) of telephone line between London and Glasgow.[citation needed]
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In 1928, Baird's company (Baird Television Development Company/Cinema Television) broadcast the first transatlantic television signal, between London and New York, and the first shore-to-ship transmission. In 1929, he became involved in the first experimental mechanical television service in Germany. In November of the same year, Baird and Bernard Natan of Pathé established France's first television company, Télévision-Baird-Natan. In 1931, he made the first outdoor remote broadcast, of The Derby.[20] In 1932, he demonstrated ultra-short wave television. Baird's mechanical system reached a peak of 240-lines of resolution on BBC television broadcasts in 1936, though the mechanical system did not scan the televised scene directly. Instead a 17.5mm film was shot, rapidly developed and then scanned while the film was still wet.[citation needed]
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An American inventor, Charles Francis Jenkins, also pioneered the television. He published an article on "Motion Pictures by Wireless" in 1913, but it was not until December 1923 that he transmitted moving silhouette images for witnesses; and it was on 13 June 1925, that he publicly demonstrated synchronized transmission of silhouette pictures. In 1925 Jenkins used the Nipkow disk and transmitted the silhouette image of a toy windmill in motion, over a distance of 5 miles (8 km), from a naval radio station in Maryland to his laboratory in Washington, D.C., using a lensed disk scanner with a 48-line resolution.[21][22] He was granted U.S. Patent No. 1,544,156 (Transmitting Pictures over Wireless) on 30 June 1925 (filed 13 March 1922).[23]
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Herbert E. Ives and Frank Gray of Bell Telephone Laboratories gave a dramatic demonstration of mechanical television on 7 April 1927. Their reflected-light television system included both small and large viewing screens. The small receiver had a 2-inch-wide by 2.5-inch-high screen (5 by 6 cm). The large receiver had a screen 24 inches wide by 30 inches high (60 by 75 cm). Both sets were capable of reproducing reasonably accurate, monochromatic, moving images. Along with the pictures, the sets received synchronized sound. The system transmitted images over two paths: first, a copper wire link from Washington to New York City, then a radio link from Whippany, New Jersey. Comparing the two transmission methods, viewers noted no difference in quality. Subjects of the telecast included Secretary of Commerce Herbert Hoover. A flying-spot scanner beam illuminated these subjects. The scanner that produced the beam had a 50-aperture disk. The disc revolved at a rate of 18 frames per second, capturing one frame about every 56 milliseconds. (Today's systems typically transmit 30 or 60 frames per second, or one frame every 33.3 or 16.7 milliseconds respectively.) Television historian Albert Abramson underscored the significance of the Bell Labs demonstration: "It was in fact the best demonstration of a mechanical television system ever made to this time. It would be several years before any other system could even begin to compare with it in picture quality."[24]
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In 1928, WRGB, then W2XB, was started as the world's first television station. It broadcast from the General Electric facility in Schenectady, NY. It was popularly known as "WGY Television". Meanwhile, in the Soviet Union, Léon Theremin had been developing a mirror drum-based television, starting with 16 lines resolution in 1925, then 32 lines and eventually 64 using interlacing in 1926. As part of his thesis, on 7 May 1926, he electrically transmitted, and then projected, near-simultaneous moving images on a 5-square-foot (0.46 m2) screen.[22]
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By 1927, Theremin had achieved an image of 100 lines, a resolution that was not surpassed until May 1932 by RCA, with 120 lines.[25]
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On 25 December 1926, Kenjiro Takayanagi demonstrated a television system with a 40-line resolution that employed a Nipkow disk scanner and CRT display at Hamamatsu Industrial High School in Japan. This prototype is still on display at the Takayanagi Memorial Museum in Shizuoka University, Hamamatsu Campus. His research in creating a production model was halted by the SCAP after World War II.[26]
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Because only a limited number of holes could be made in the disks, and disks beyond a certain diameter became impractical, image resolution on mechanical television broadcasts was relatively low, ranging from about 30 lines up to 120 or so. Nevertheless, the image quality of 30-line transmissions steadily improved with technical advances, and by 1933 the UK broadcasts using the Baird system were remarkably clear.[27] A few systems ranging into the 200-line region also went on the air. Two of these were the 180-line system that Compagnie des Compteurs (CDC) installed in Paris in 1935, and the 180-line system that Peck Television Corp. started in 1935 at station VE9AK in Montreal.[28][29] The advancement of all-electronic television (including image dissectors and other camera tubes and cathode ray tubes for the reproducer) marked the beginning of the end for mechanical systems as the dominant form of television. Mechanical television, despite its inferior image quality and generally smaller picture, would remain the primary television technology until the 1930s. The last mechanical television broadcasts ended in 1939 at stations run by a handful of public universities in the United States.[citation needed]
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In 1897, English physicist J. J. Thomson was able, in his three famous experiments, to deflect cathode rays, a fundamental function of the modern cathode ray tube (CRT). The earliest version of the CRT was invented by the German physicist Ferdinand Braun in 1897 and is also known as the "Braun" tube.[30][31] It was a cold-cathode diode, a modification of the Crookes tube, with a phosphor-coated screen. In 1906 the Germans Max Dieckmann and Gustav Glage produced raster images for the first time in a CRT.[32] In 1907, Russian scientist Boris Rosing used a CRT in the receiving end of an experimental video signal to form a picture. He managed to display simple geometric shapes onto the screen.[33]
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In 1908 Alan Archibald Campbell-Swinton, fellow of the Royal Society (UK), published a letter in the scientific journal Nature in which he described how "distant electric vision" could be achieved by using a cathode ray tube, or Braun tube, as both a transmitting and receiving device,[34][35] He expanded on his vision in a speech given in London in 1911 and reported in The Times[36] and the Journal of the Röntgen Society.[37][38] In a letter to Nature published in October 1926, Campbell-Swinton also announced the results of some "not very successful experiments" he had conducted with G. M. Minchin and J. C. M. Stanton. They had attempted to generate an electrical signal by projecting an image onto a selenium-coated metal plate that was simultaneously scanned by a cathode ray beam.[39][40] These experiments were conducted before March 1914, when Minchin died,[41] but they were later repeated by two different teams in 1937, by H. Miller and J. W. Strange from EMI,[42] and by H. Iams and A. Rose from RCA.[43] Both teams succeeded in transmitting "very faint" images with the original Campbell-Swinton's selenium-coated plate. Although others had experimented with using a cathode ray tube as a receiver, the concept of using one as a transmitter was novel.[44] The first cathode ray tube to use a hot cathode was developed by John B. Johnson (who gave his name to the term Johnson noise) and Harry Weiner Weinhart of Western Electric, and became a commercial product in 1922.[citation needed]
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In 1926, Hungarian engineer Kálmán Tihanyi designed a television system utilizing fully electronic scanning and display elements and employing the principle of "charge storage" within the scanning (or "camera") tube.[45][46][47][48] The problem of low sensitivity to light resulting in low electrical output from transmitting or "camera" tubes would be solved with the introduction of charge-storage technology by Kálmán Tihanyi beginning in 1924.[49] His solution was a camera tube that accumulated and stored electrical charges ("photoelectrons") within the tube throughout each scanning cycle. The device was first described in a patent application he filed in Hungary in March 1926 for a television system he dubbed "Radioskop".[50] After further refinements included in a 1928 patent application,[49] Tihanyi's patent was declared void in Great Britain in 1930,[51] so he applied for patents in the United States. Although his breakthrough would be incorporated into the design of RCA's "iconoscope" in 1931, the U.S. patent for Tihanyi's transmitting tube would not be granted until May 1939. The patent for his receiving tube had been granted the previous October. Both patents had been purchased by RCA prior to their approval.[52][53] Charge storage remains a basic principle in the design of imaging devices for television to the present day.[50] On 25 December 1926, at Hamamatsu Industrial High School in Japan, Japanese inventor Kenjiro Takayanagi demonstrated a TV system with a 40-line resolution that employed a CRT display.[26] This was the first working example of a fully electronic television receiver. Takayanagi did not apply for a patent.[54]
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On 7 September 1927, American inventor Philo Farnsworth's image dissector camera tube transmitted its first image, a simple straight line, at his laboratory at 202 Green Street in San Francisco.[55][56] By 3 September 1928, Farnsworth had developed the system sufficiently to hold a demonstration for the press. This is widely regarded as the first electronic television demonstration.[56] In 1929, the system was improved further by the elimination of a motor generator, so that his television system now had no mechanical parts.[57] That year, Farnsworth transmitted the first live human images with his system, including a three and a half-inch image of his wife Elma ("Pem") with her eyes closed (possibly due to the bright lighting required).[58]
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Meanwhile, Vladimir Zworykin was also experimenting with the cathode ray tube to create and show images. While working for Westinghouse Electric in 1923, he began to develop an electronic camera tube. But in a 1925 demonstration, the image was dim, had low contrast, and poor definition, and was stationary.[59] Zworykin's imaging tube never got beyond the laboratory stage. But RCA, which acquired the Westinghouse patent, asserted that the patent for Farnsworth's 1927 image dissector was written so broadly that it would exclude any other electronic imaging device. Thus RCA, on the basis of Zworykin's 1923 patent application, filed a patent interference suit against Farnsworth. The U.S. Patent Office examiner disagreed in a 1935 decision, finding priority of invention for Farnsworth against Zworykin. Farnsworth claimed that Zworykin's 1923 system would be unable to produce an electrical image of the type to challenge his patent. Zworykin received a patent in 1928 for a color transmission version of his 1923 patent application;[60] he also divided his original application in 1931.[61] Zworykin was unable or unwilling to introduce evidence of a working model of his tube that was based on his 1923 patent application. In September 1939, after losing an appeal in the courts, and determined to go forward with the commercial manufacturing of television equipment, RCA agreed to pay Farnsworth US$1 million over a ten-year period, in addition to license payments, to use his patents.[62][63]
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In 1933, RCA introduced an improved camera tube that relied on Tihanyi's charge storage principle.[64] Dubbed the "Iconoscope" by Zworykin, the new tube had a light sensitivity of about 75,000 lux, and thus was claimed to be much more sensitive than Farnsworth's image dissector.[citation needed] However, Farnsworth had overcome his power problems with his Image Dissector through the invention of a completely unique "multipactor" device that he began work on in 1930, and demonstrated in 1931.[65][66] This small tube could amplify a signal reportedly to the 60th power or better[67] and showed great promise in all fields of electronics. Unfortunately, a problem with the multipactor was that it wore out at an unsatisfactory rate.[68]
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At the Berlin Radio Show in August 1931, Manfred von Ardenne gave a public demonstration of a television system using a CRT for both transmission and reception. However, Ardenne had not developed a camera tube, using the CRT instead as a flying-spot scanner to scan slides and film.[69] Philo Farnsworth gave the world's first public demonstration of an all-electronic television system, using a live camera, at the Franklin Institute of Philadelphia on 25 August 1934, and for ten days afterwards.[70][71] Mexican inventor Guillermo González Camarena also played an important role in early TV. His experiments with TV (known as telectroescopía at first) began in 1931 and led to a patent for the "trichromatic field sequential system" color television in 1940.[72] In Britain, the EMI engineering team led by Isaac Shoenberg applied in 1932 for a patent for a new device they dubbed "the Emitron",[73][74] which formed the heart of the cameras they designed for the BBC. On 2 November 1936, a 405-line broadcasting service employing the Emitron began at studios in Alexandra Palace, and transmitted from a specially built mast atop one of the Victorian building's towers. It alternated for a short time with Baird's mechanical system in adjoining studios, but was more reliable and visibly superior. This was the world's first regular "high-definition" television service.[75]
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The original American iconoscope was noisy, had a high ratio of interference to signal, and ultimately gave disappointing results, especially when compared to the high definition mechanical scanning systems then becoming available.[76][77] The EMI team, under the supervision of Isaac Shoenberg, analyzed how the iconoscope (or Emitron) produces an electronic signal and concluded that its real efficiency was only about 5% of the theoretical maximum.[78][79] They solved this problem by developing, and patenting in 1934, two new camera tubes dubbed super-Emitron and CPS Emitron.[80][81][82] The super-Emitron was between ten and fifteen times more sensitive than the original Emitron and iconoscope tubes and, in some cases, this ratio was considerably greater.[78] It was used for outside broadcasting by the BBC, for the first time, on Armistice Day 1937, when the general public could watch on a television set as the King laid a wreath at the Cenotaph.[83] This was the first time that anyone had broadcast a live street scene from cameras installed on the roof of neighboring buildings, because neither Farnsworth nor RCA would do the same until the 1939 New York World's Fair.
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On the other hand, in 1934, Zworykin shared some patent rights with the German licensee company Telefunken.[84] The "image iconoscope" ("Superikonoskop" in Germany) was produced as a result of the collaboration. This tube is essentially identical to the super-Emitron.[citation needed] The production and commercialization of the super-Emitron and image iconoscope in Europe were not affected by the patent war between Zworykin and Farnsworth, because Dieckmann and Hell had priority in Germany for the invention of the image dissector, having submitted a patent application for their Lichtelektrische Bildzerlegerröhre für Fernseher (Photoelectric Image Dissector Tube for Television) in Germany in 1925,[85] two years before Farnsworth did the same in the United States.[86] The image iconoscope (Superikonoskop) became the industrial standard for public broadcasting in Europe from 1936 until 1960, when it was replaced by the vidicon and plumbicon tubes. Indeed, it was the representative of the European tradition in electronic tubes competing against the American tradition represented by the image orthicon.[87][88] The German company Heimann produced the Superikonoskop for the 1936 Berlin Olympic Games,[89][90] later Heimann also produced and commercialized it from 1940 to 1955;[91] finally the Dutch company Philips produced and commercialized the image iconoscope and multicon from 1952 to 1958.[88][92]
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American television broadcasting, at the time, consisted of a variety of markets in a wide range of sizes, each competing for programming and dominance with separate technology, until deals were made and standards agreed upon in 1941.[93] RCA, for example, used only Iconoscopes in the New York area, but Farnsworth Image Dissectors in Philadelphia and San Francisco.[94] In September 1939, RCA agreed to pay the Farnsworth Television and Radio Corporation royalties over the next ten years for access to Farnsworth's patents.[95] With this historic agreement in place, RCA integrated much of what was best about the Farnsworth Technology into their systems.[94] In 1941, the United States implemented 525-line television.[96][97] Electrical engineer Benjamin Adler played a prominent role in the development of television.[98][99]
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The world's first 625-line television standard was designed in the Soviet Union in 1944 and became a national standard in 1946.[100] The first broadcast in 625-line standard occurred in Moscow in 1948.[101] The concept of 625 lines per frame was subsequently implemented in the European CCIR standard.[102] In 1936, Kálmán Tihanyi described the principle of plasma display, the first flat panel display system.[103][104]
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Early electronic television sets were large and bulky, with analog circuits made of vacuum tubes. Following the invention of the first working transistor at Bell Labs, Sony founder Masaru Ibuka predicted in 1952 that the transition to electronic circuits made of transistors would lead to smaller and more portable television sets.[105] The first fully transistorized, portable solid-state television set was the 8-inch Sony TV8-301, developed in 1959 and released in 1960.[106][107] This began the transformation of television viewership from a communal viewing experience to a solitary viewing experience.[108] By 1960, Sony had sold over 4 million portable television sets worldwide.[109]
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The basic idea of using three monochrome images to produce a color image had been experimented with almost as soon as black-and-white televisions had first been built. Although he gave no practical details, among the earliest published proposals for television was one by Maurice Le Blanc, in 1880, for a color system, including the first mentions in television literature of line and frame scanning.[110] Polish inventor Jan Szczepanik patented a color television system in 1897, using a selenium photoelectric cell at the transmitter and an electromagnet controlling an oscillating mirror and a moving prism at the receiver. But his system contained no means of analyzing the spectrum of colors at the transmitting end, and could not have worked as he described it.[111] Another inventor, Hovannes Adamian, also experimented with color television as early as 1907. The first color television project is claimed by him,[112] and was patented in Germany on 31 March 1908, patent No. 197183, then in Britain, on 1 April 1908, patent No. 7219,[113] in France (patent No. 390326) and in Russia in 1910 (patent No. 17912).[114]
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Scottish inventor John Logie Baird demonstrated the world's first color transmission on 3 July 1928, using scanning discs at the transmitting and receiving ends with three spirals of apertures, each spiral with filters of a different primary color; and three light sources at the receiving end, with a commutator to alternate their illumination.[115] Baird also made the world's first color broadcast on 4 February 1938, sending a mechanically scanned 120-line image from Baird's Crystal Palace studios to a projection screen at London's Dominion Theatre.[116] Mechanically scanned color television was also demonstrated by Bell Laboratories in June 1929 using three complete systems of photoelectric cells, amplifiers, glow-tubes, and color filters, with a series of mirrors to superimpose the red, green, and blue images into one full color image.
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The first practical hybrid system was again pioneered by John Logie Baird. In 1940 he publicly demonstrated a color television combining a traditional black-and-white display with a rotating colored disk. This device was very "deep", but was later improved with a mirror folding the light path into an entirely practical device resembling a large conventional console.[117] However, Baird was not happy with the design, and, as early as 1944, had commented to a British government committee that a fully electronic device would be better.
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In 1939, Hungarian engineer Peter Carl Goldmark introduced an electro-mechanical system while at CBS, which contained an Iconoscope sensor. The CBS field-sequential color system was partly mechanical, with a disc made of red, blue, and green filters spinning inside the television camera at 1,200 rpm, and a similar disc spinning in synchronization in front of the cathode ray tube inside the receiver set.[118] The system was first demonstrated to the Federal Communications Commission (FCC) on 29 August 1940, and shown to the press on 4 September.[119][120][121][122]
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CBS began experimental color field tests using film as early as 28 August 1940, and live cameras by 12 November.[120][123] NBC (owned by RCA) made its first field test of color television on 20 February 1941. CBS began daily color field tests on 1 June 1941.[124] These color systems were not compatible with existing black-and-white television sets, and, as no color television sets were available to the public at this time, viewing of the color field tests was restricted to RCA and CBS engineers and the invited press. The War Production Board halted the manufacture of television and radio equipment for civilian use from 22 April 1942 to 20 August 1945, limiting any opportunity to introduce color television to the general public.[125][126]
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As early as 1940, Baird had started work on a fully electronic system he called Telechrome. Early Telechrome devices used two electron guns aimed at either side of a phosphor plate. The phosphor was patterned so the electrons from the guns only fell on one side of the patterning or the other. Using cyan and magenta phosphors, a reasonable limited-color image could be obtained. He also demonstrated the same system using monochrome signals to produce a 3D image (called "stereoscopic" at the time). A demonstration on 16 August 1944 was the first example of a practical color television system. Work on the Telechrome continued and plans were made to introduce a three-gun version for full color. However, Baird's untimely death in 1946 ended development of the Telechrome system.[127][128]
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Similar concepts were common through the 1940s and 1950s, differing primarily in the way they re-combined the colors generated by the three guns. The Geer tube was similar to Baird's concept, but used small pyramids with the phosphors deposited on their outside faces, instead of Baird's 3D patterning on a flat surface. The Penetron used three layers of phosphor on top of each other and increased the power of the beam to reach the upper layers when drawing those colors. The Chromatron used a set of focusing wires to select the colored phosphors arranged in vertical stripes on the tube.
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One of the great technical challenges of introducing color broadcast television was the desire to conserve bandwidth, potentially three times that of the existing black-and-white standards, and not use an excessive amount of radio spectrum. In the United States, after considerable research, the National Television Systems Committee[129] approved an all-electronic system developed by RCA, which encoded the color information separately from the brightness information and greatly reduced the resolution of the color information in order to conserve bandwidth. As black-and-white TVs could receive the same transmission and display it in black-and-white, the color system adopted is [backwards] "compatible". ("Compatible Color", featured in RCA advertisements of the period, is mentioned in the song "America", of West Side Story, 1957.) The brightness image remained compatible with existing black-and-white television sets at slightly reduced resolution, while color televisions could decode the extra information in the signal and produce a limited-resolution color display. The higher resolution black-and-white and lower resolution color images combine in the brain to produce a seemingly high-resolution color image. The NTSC standard represented a major technical achievement.
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The first color broadcast (the first episode of the live program The Marriage (TV series)) occurred on 8 July 1954, but during the following ten years most network broadcasts, and nearly all local programming, continued to be in black-and-white. It was not until the mid-1960s that color sets started selling in large numbers, due in part to the color transition of 1965 in which it was announced that over half of all network prime-time programming would be broadcast in color that fall. The first all-color prime-time season came just one year later. In 1972, the last holdout among daytime network programs converted to color, resulting in the first completely all-color network season.
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Early color sets were either floor-standing console models or tabletop versions nearly as bulky and heavy, so in practice they remained firmly anchored in one place. GE's relatively compact and lightweight Porta-Color set was introduced in the spring of 1966. It used a transistor-based UHF tuner.[130] The first fully transistorized color television in the United States was the Quasar television introduced in 1967.[131] These developments made watching color television a more flexible and convenient proposition.
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The MOSFET (metal-oxide-semiconductor field-effect transistor, or MOS transistor) was invented by Mohamed M. Atalla and Dawon Kahng at Bell Labs in 1959,[132] and presented in 1960.[133] By the mid-1960s, RCA were using MOSFETs in their consumer television products.[134] RCA Laboratories researchers W.M. Austin, J.A. Dean, D.M. Griswold and O.P. Hart in 1966 described the use of the MOSFET in television circuits, including RF amplifier, low-level video, chroma and AGC circuits.[135] The power MOSFET was later widely adopted for television receiver circuits.[136]
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In 1972, sales of color sets finally surpassed sales of black-and-white sets. Color broadcasting in Europe was not standardized on the PAL format until the 1960s, and broadcasts did not start until 1967. By this point many of the technical problems in the early sets had been worked out, and the spread of color sets in Europe was fairly rapid. By the mid-1970s, the only stations broadcasting in black-and-white were a few high-numbered UHF stations in small markets, and a handful of low-power repeater stations in even smaller markets such as vacation spots. By 1979, even the last of these had converted to color and, by the early 1980s, B&W sets had been pushed into niche markets, notably low-power uses, small portable sets, or for use as video monitor screens in lower-cost consumer equipment. By the late 1980s even these areas switched to color sets.
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Digital television (DTV) is the transmission of audio and video by digitally processed and multiplexed signals, in contrast to the totally analog and channel separated signals used by analog television. Due to data compression, digital TV can support more than one program in the same channel bandwidth.[137] It is an innovative service that represents the most significant evolution in television broadcast technology since color television emerged in the 1950s.[138] Digital TV's roots have been tied very closely to the availability of inexpensive, high performance computers. It was not until the 1990s that digital TV became feasible.[139] Digital television was previously not practically feasible due to the impractically high bandwidth requirements of uncompressed digital video,[140][141] requiring around 200 Mbit/s bit-rate for a standard-definition television (SDTV) signal,[140] and over 1 Gbit/s for high-definition television (HDTV).[141]
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Digital TV became practically feasible in the early 1990s due to a major technological development, discrete cosine transform (DCT) video compression.[140][141] DCT coding is a lossy compression technique that was first proposed for image compression by Nasir Ahmed in 1972,[142] and was later adapted into a motion-compensated DCT video coding algorithm, for video coding standards such as the H.26x formats from 1988 onwards and the MPEG formats from 1991 onwards.[143][144] Motion-compensated DCT video compression significantly reduced the amount of bandwidth required for a digital TV signal.[140][141] DCT coding compressed down the bandwidth requirements of digital television signals to about 34 Mpps bit-rate for SDTV and around 70–140 Mbit/s for HDTV while maintaining near-studio-quality transmission, making digital television a practical reality in the 1990s.[141]
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A digital TV service was proposed in 1986 by Nippon Telegraph and Telephone (NTT) and the Ministry of Posts and Telecommunication (MPT) in Japan, where there were plans to develop an "Integrated Network System" service. However, it was not possible to practically implement such a digital TV service until the adoption of DCT video compression technology made it possible in the early 1990s.[140]
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In the mid-1980s, as Japanese consumer electronics firms forged ahead with the development of HDTV technology, the MUSE analog format proposed by NHK, a Japanese company, was seen as a pacesetter that threatened to eclipse U.S. electronics companies' technologies. Until June 1990, the Japanese MUSE standard, based on an analog system, was the front-runner among the more than 23 different technical concepts under consideration. Then, an American company, General Instrument, demonstrated the feasibility of a digital television signal. This breakthrough was of such significance that the FCC was persuaded to delay its decision on an ATV standard until a digitally based standard could be developed.
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In March 1990, when it became clear that a digital standard was feasible, the FCC made a number of critical decisions. First, the Commission declared that the new ATV standard must be more than an enhanced analog signal, but be able to provide a genuine HDTV signal with at least twice the resolution of existing television images.(7) Then, to ensure that viewers who did not wish to buy a new digital television set could continue to receive conventional television broadcasts, it dictated that the new ATV standard must be capable of being "simulcast" on different channels.(8)The new ATV standard also allowed the new DTV signal to be based on entirely new design principles. Although incompatible with the existing NTSC standard, the new DTV standard would be able to incorporate many improvements.
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The final standards adopted by the FCC did not require a single standard for scanning formats, aspect ratios, or lines of resolution. This compromise resulted from a dispute between the consumer electronics industry (joined by some broadcasters) and the computer industry (joined by the film industry and some public interest groups) over which of the two scanning processes—interlaced or progressive—would be best suited for the newer digital HDTV compatible display devices.[145] Interlaced scanning, which had been specifically designed for older analogue CRT display technologies, scans even-numbered lines first, then odd-numbered ones. In fact, interlaced scanning can be looked at as the first video compression model as it was partly designed in the 1940s to double the image resolution to exceed the limitations of the television broadcast bandwidth. Another reason for its adoption was to limit the flickering on early CRT screens whose phosphor coated screens could only retain the image from the electron scanning gun for a relatively short duration.[146] However interlaced scanning does not work as efficiently on newer display devices such as Liquid-crystal (LCD), for example, which are better suited to a more frequent progressive refresh rate.[145]
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Progressive scanning, the format that the computer industry had long adopted for computer display monitors, scans every line in sequence, from top to bottom. Progressive scanning in effect doubles the amount of data generated for every full screen displayed in comparison to interlaced scanning by painting the screen in one pass in 1/60-second, instead of two passes in 1/30-second. The computer industry argued that progressive scanning is superior because it does not "flicker" on the new standard of display devices in the manner of interlaced scanning. It also argued that progressive scanning enables easier connections with the Internet, and is more cheaply converted to interlaced formats than vice versa. The film industry also supported progressive scanning because it offered a more efficient means of converting filmed programming into digital formats. For their part, the consumer electronics industry and broadcasters argued that interlaced scanning was the only technology that could transmit the highest quality pictures then (and currently) feasible, i.e., 1,080 lines per picture and 1,920 pixels per line. Broadcasters also favored interlaced scanning because their vast archive of interlaced programming is not readily compatible with a progressive format. William F. Schreiber, who was director of the Advanced Television Research Program at the Massachusetts Institute of Technology from 1983 until his retirement in 1990, thought that the continued advocacy of interlaced equipment originated from consumer electronics companies that were trying to get back the substantial investments they made in the interlaced technology.[147]
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Digital television transition started in late 2000s. All governments across the world set the deadline for analog shutdown by 2010s. Initially the adoption rate was low, as the first digital tuner-equipped TVs were costly. But soon, as the price of digital-capable TVs dropped, more and more households were converting to digital televisions. The transition is expected to be completed worldwide by mid to late 2010s.
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The advent of digital television allowed innovations like smart TVs. A smart television, sometimes referred to as connected TV or hybrid TV, is a television set or set-top box with integrated Internet and Web 2.0 features, and is an example of technological convergence between computers, television sets and set-top boxes. Besides the traditional functions of television sets and set-top boxes provided through traditional Broadcasting media, these devices can also provide Internet TV, online interactive media, over-the-top content, as well as on-demand streaming media, and home networking access. These TVs come pre-loaded with an operating system.[9][148][149][150]
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Smart TV should not to be confused with Internet TV, Internet Protocol television (IPTV) or with Web TV. Internet television refers to the receiving of television content over the Internet instead of by traditional systems—terrestrial, cable and satellite (although internet itself is received by these methods). IPTV is one of the emerging Internet television technology standards for use by television broadcasters. Web television (WebTV) is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV. A first patent was filed in 1994[151] (and extended the following year)[152] for an "intelligent" television system, linked with data processing systems, by means of a digital or analog network. Apart from being linked to data networks, one key point is its ability to automatically download necessary software routines, according to a user's demand, and process their needs. Major TV manufacturers have announced production of smart TVs only, for middle-end and high-end TVs in 2015.[6][7][8] Smart TVs have gotten more affordable compared to when they were first introduced, with 46 million of U.S. households having at least one as of 2019.[153]
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3D television conveys depth perception to the viewer by employing techniques such as stereoscopic display, multi-view display, 2D-plus-depth, or any other form of 3D display. Most modern 3D television sets use an active shutter 3D system or a polarized 3D system, and some are autostereoscopic without the need of glasses. Stereoscopic 3D television was demonstrated for the first time on 10 August 1928, by John Logie Baird in his company's premises at 133 Long Acre, London.[154] Baird pioneered a variety of 3D television systems using electromechanical and cathode-ray tube techniques. The first 3D TV was produced in 1935. The advent of digital television in the 2000s greatly improved 3D TVs. Although 3D TV sets are quite popular for watching 3D home media such as on Blu-ray discs, 3D programming has largely failed to make inroads with the public. Many 3D television channels which started in the early 2010s were shut down by the mid-2010s. According to DisplaySearch 3D televisions shipments totaled 41.45 million units in 2012, compared with 24.14 in 2011 and 2.26 in 2010.[155] As of late 2013, the number of 3D TV viewers started to decline.[156][157][158][159][160]
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Programming is broadcast by television stations, sometimes called "channels", as stations are licensed by their governments to broadcast only over assigned channels in the television band. At first, terrestrial broadcasting was the only way television could be widely distributed, and because bandwidth was limited, i.e., there were only a small number of channels available, government regulation was the norm. In the U.S., the Federal Communications Commission (FCC) allowed stations to broadcast advertisements beginning in July 1941, but required public service programming commitments as a requirement for a license. By contrast, the United Kingdom chose a different route, imposing a television license fee on owners of television reception equipment to fund the British Broadcasting Corporation (BBC), which had public service as part of its Royal Charter.
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WRGB claims to be the world's oldest television station, tracing its roots to an experimental station founded on 13 January 1928, broadcasting from the General Electric factory in Schenectady, NY, under the call letters W2XB.[161] It was popularly known as "WGY Television" after its sister radio station. Later in 1928, General Electric started a second facility, this one in New York City, which had the call letters W2XBS and which today is known as WNBC. The two stations were experimental in nature and had no regular programming, as receivers were operated by engineers within the company. The image of a Felix the Cat doll rotating on a turntable was broadcast for 2 hours every day for several years as new technology was being tested by the engineers. On 2 November 1936, the BBC began transmitting the world's first public regular high-definition service from the Victorian Alexandra Palace in north London.[162] It therefore claims to be the birthplace of TV broadcasting as we know it today.
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With the widespread adoption of cable across the United States in the 1970s and 80s, terrestrial television broadcasts have been in decline; in 2013 it was estimated that about 7% of US households used an antenna.[163][164] A slight increase in use began around 2010 due to switchover to digital terrestrial television broadcasts, which offered pristine image quality over very large areas, and offered an alternate to cable television (CATV) for cord cutters. All other countries around the world are also in the process of either shutting down analog terrestrial television or switching over to digital terrestrial television.
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Cable television is a system of broadcasting television programming to paying subscribers via radio frequency (RF) signals transmitted through coaxial cables or light pulses through fiber-optic cables. This contrasts with traditional terrestrial television, in which the television signal is transmitted over the air by radio waves and received by a television antenna attached to the television. In the 2000s, FM radio programming, high-speed Internet, telephone service, and similar non-television services may also be provided through these cables. The abbreviation CATV is often used for cable television. It originally stood for Community Access Television or Community Antenna Television, from cable television's origins in 1948: in areas where over-the-air reception was limited by distance from transmitters or mountainous terrain, large "community antennas" were constructed, and cable was run from them to individual homes.[165] The origins of cable broadcasting are even older as radio programming was distributed by cable in some European cities as far back as 1924. Earlier cable television was analog, but since the 2000s, all cable operators have switched to, or are in the process of switching to, digital cable television.
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Satellite television is a system of supplying television programming using broadcast signals relayed from communication satellites. The signals are received via an outdoor parabolic reflector antenna usually referred to as a satellite dish and a low-noise block downconverter (LNB). A satellite receiver then decodes the desired television program for viewing on a television set. Receivers can be external set-top boxes, or a built-in television tuner. Satellite television provides a wide range of channels and services, especially to geographic areas without terrestrial television or cable television.
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The most common method of reception is direct-broadcast satellite television (DBSTV), also known as "direct to home" (DTH).[166] In DBSTV systems, signals are relayed from a direct broadcast satellite on the Ku wavelength and are completely digital.[167] Satellite TV systems formerly used systems known as television receive-only. These systems received analog signals transmitted in the C-band spectrum from FSS type satellites, and required the use of large dishes. Consequently, these systems were nicknamed "big dish" systems, and were more expensive and less popular.[168]
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The direct-broadcast satellite television signals were earlier analog signals and later digital signals, both of which require a compatible receiver. Digital signals may include high-definition television (HDTV). Some transmissions and channels are free-to-air or free-to-view, while many other channels are pay television requiring a subscription.[169]
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In 1945, British science fiction writer Arthur C. Clarke proposed a worldwide communications system which would function by means of three satellites equally spaced apart in earth orbit.[170][171] This was published in the October 1945 issue of the Wireless World magazine and won him the Franklin Institute's Stuart Ballantine Medal in 1963.[172][173]
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The first satellite television signals from Europe to North America were relayed via the Telstar satellite over the Atlantic ocean on 23 July 1962.[174] The signals were received and broadcast in North American and European countries and watched by over 100 million.[174] Launched in 1962, the Relay 1 satellite was the first satellite to transmit television signals from the US to Japan.[175] The first geosynchronous communication satellite, Syncom 2, was launched on 26 July 1963.[176]
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The world's first commercial communications satellite, called Intelsat I and nicknamed "Early Bird", was launched into geosynchronous orbit on 6 April 1965.[177] The first national network of television satellites, called Orbita, was created by the Soviet Union in October 1967, and was based on the principle of using the highly elliptical Molniya satellite for rebroadcasting and delivering of television signals to ground downlink stations.[178] The first commercial North American satellite to carry television transmissions was Canada's geostationary Anik 1, which was launched on 9 November 1972.[179] ATS-6, the world's first experimental educational and Direct Broadcast Satellite (DBS), was launched on 30 May 1974.[180] It transmitted at 860 MHz using wideband FM modulation and had two sound channels. The transmissions were focused on the Indian subcontinent but experimenters were able to receive the signal in Western Europe using home constructed equipment that drew on UHF television design techniques already in use.[181]
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The first in a series of Soviet geostationary satellites to carry Direct-To-Home television, Ekran 1, was launched on 26 October 1976.[182] It used a 714 MHz UHF downlink frequency so that the transmissions could be received with existing UHF television technology rather than microwave technology.[183]
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Internet television (Internet TV) (or online television) is the digital distribution of television content via the Internet as opposed to traditional systems like terrestrial, cable, and satellite, although the Internet itself is received by terrestrial, cable, or satellite methods. Internet television is a general term that covers the delivery of television shows, and other video content, over the Internet by video streaming technology, typically by major traditional television broadcasters. Internet television should not be confused with Smart TV, IPTV or with Web TV. Smart television refers to the TV set which has a built-in operating system. Internet Protocol television (IPTV) is one of the emerging Internet television technology standards for use by television broadcasters. Web television is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV.
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A television set, also called a television receiver, television, TV set, TV, or "telly", is a device that combines a tuner, display, an amplifier, and speakers for the purpose of viewing television and hearing its audio components. Introduced in the late 1920s in mechanical form, television sets became a popular consumer product after World War II in electronic form, using cathode ray tubes. The addition of color to broadcast television after 1953 further increased the popularity of television sets and an outdoor antenna became a common feature of suburban homes. The ubiquitous television set became the display device for recorded media in the 1970s, such as Betamax and VHS, which enabled viewers to record TV shows and watch prerecorded movies. In the subsequent decades, TVs were used to watch DVDs and Blu-ray Discs of movies and other content. Major TV manufacturers announced the discontinuation of CRT, DLP, plasma and fluorescent-backlit LCDs by the mid-2010s. Televisions since 2010s mostly use LEDs.[3][4][184][185] LEDs are expected to be gradually replaced by OLEDs in the near future.[5]
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The earliest systems employed a spinning disk to create and reproduce images.[186] These usually had a low resolution and screen size and never became popular with the public.
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The cathode ray tube (CRT) is a vacuum tube containing one or more electron guns (a source of electrons or electron emitter) and a fluorescent screen used to view images.[33] It has a means to accelerate and deflect the electron beam(s) onto the screen to create the images. The images may represent electrical waveforms (oscilloscope), pictures (television, computer monitor), radar targets or others. The CRT uses an evacuated glass envelope which is large, deep (i.e. long from front screen face to rear end), fairly heavy, and relatively fragile. As a matter of safety, the face is typically made of thick lead glass so as to be highly shatter-resistant and to block most X-ray emissions, particularly if the CRT is used in a consumer product.
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In television sets and computer monitors, the entire front area of the tube is scanned repetitively and systematically in a fixed pattern called a raster. An image is produced by controlling the intensity of each of the three electron beams, one for each additive primary color (red, green, and blue) with a video signal as a reference.[187] In all modern CRT monitors and televisions, the beams are bent by magnetic deflection, a varying magnetic field generated by coils and driven by electronic circuits around the neck of the tube, although electrostatic deflection is commonly used in oscilloscopes, a type of diagnostic instrument.[187]
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Digital Light Processing (DLP) is a type of video projector technology that uses a digital micromirror device. Some DLPs have a TV tuner, which makes them a type of TV display. It was originally developed in 1987 by Dr. Larry Hornbeck of Texas Instruments. While the DLP imaging device was invented by Texas Instruments, the first DLP based projector was introduced by Digital Projection Ltd in 1997. Digital Projection and Texas Instruments were both awarded Emmy Awards in 1998 for invention of the DLP projector technology. DLP is used in a variety of display applications from traditional static displays to interactive displays and also non-traditional embedded applications including medical, security, and industrial uses. DLP technology is used in DLP front projectors (standalone projection units for classrooms and business primarily), but also in private homes; in these cases, the image is projected onto a projection screen. DLP is also used in DLP rear projection television sets and digital signs. It is also used in about 85% of digital cinema projection.[188]
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A plasma display panel (PDP) is a type of flat panel display common to large TV displays 30 inches (76 cm) or larger. They are called "plasma" displays because the technology utilizes small cells containing electrically charged ionized gases, or what are in essence chambers more commonly known as fluorescent lamps.
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Liquid-crystal-display televisions (LCD TV) are television sets that use LCD display technology to produce images. LCD televisions are much thinner and lighter than cathode ray tube (CRTs) of similar display size, and are available in much larger sizes (e.g., 90-inch diagonal). When manufacturing costs fell, this combination of features made LCDs practical for television receivers. LCDs come in two types: those using cold cathode fluorescent lamps, simply called LCDs and those using LED as backlight called as LEDs.
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In 2007, LCD televisions surpassed sales of CRT-based televisions worldwide for the first time, and their sales figures relative to other technologies accelerated. LCD TVs have quickly displaced the only major competitors in the large-screen market, the Plasma display panel and rear-projection television.[189] In mid 2010s LCDs especially LEDs became, by far, the most widely produced and sold television display type.[184][185] LCDs also have disadvantages. Other technologies address these weaknesses, including OLEDs, FED and SED, but as of 2014[update] none of these have entered widespread production.
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An OLED (organic light-emitting diode) is a light-emitting diode (LED) in which the emissive electroluminescent layer is a film of organic compound which emits light in response to an electric current. This layer of organic semiconductor is situated between two electrodes. Generally, at least one of these electrodes is transparent. OLEDs are used to create digital displays in devices such as television screens. It is also used for computer monitors, portable systems such as mobile phones, handheld game consoles and PDAs.
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There are two main families of OLED: those based on small molecules and those employing polymers. Adding mobile ions to an OLED creates a light-emitting electrochemical cell or LEC, which has a slightly different mode of operation. OLED displays can use either passive-matrix (PMOLED) or active-matrix (AMOLED) addressing schemes. Active-matrix OLEDs require a thin-film transistor backplane to switch each individual pixel on or off, but allow for higher resolution and larger display sizes.
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An OLED display works without a backlight. Thus, it can display deep black levels and can be thinner and lighter than a liquid crystal display (LCD). In low ambient light conditions such as a dark room an OLED screen can achieve a higher contrast ratio than an LCD, whether the LCD uses cold cathode fluorescent lamps or LED backlight. OLEDs are expected to replace other forms of display in near future.[5]
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Low-definition television or LDTV refers to television systems that have a lower screen resolution than standard-definition television systems such 240p (320*240). It is used in handheld television. The most common source of LDTV programming is the Internet, where mass distribution of higher-resolution video files could overwhelm computer servers and take too long to download. Many mobile phones and portable devices such as Apple's iPod Nano, or Sony's PlayStation Portable use LDTV video, as higher-resolution files would be excessive to the needs of their small screens (320×240 and 480×272 pixels respectively). The current generation of iPod Nanos have LDTV screens, as do the first three generations of iPod Touch and iPhone (480×320). For the first years of its existence, YouTube offered only one, low-definition resolution of 320x240p at 30fps or less. A standard, consumer grade VHS videotape can be considered SDTV due to its resolution (approximately 360 × 480i/576i).
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Standard-definition television or SDTV refers to two different resolutions: 576i, with 576 interlaced lines of resolution, derived from the European-developed PAL and SECAM systems; and 480i based on the American National Television System Committee NTSC system. SDTV is a television system that uses a resolution that is not considered to be either high-definition television (720p, 1080i, 1080p, 1440p, 4K UHDTV, and 8K UHD) or enhanced-definition television (EDTV 480p). In North America, digital SDTV is broadcast in the same 4:3 aspect ratio as NTSC signals with widescreen content being center cut.[190] However, in other parts of the world that used the PAL or SECAM color systems, standard-definition television is now usually shown with a 16:9 aspect ratio, with the transition occurring between the mid-1990s and mid-2000s. Older programs with a 4:3 aspect ratio are shown in the US as 4:3 with non-ATSC countries preferring to reduce the horizontal resolution by anamorphically scaling a pillarboxed image.
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High-definition television (HDTV) provides a resolution that is substantially higher than that of standard-definition television.
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HDTV may be transmitted in various formats:
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Ultra-high-definition television (also known as Super Hi-Vision, Ultra HD television, UltraHD, UHDTV, or UHD) includes 4K UHD (2160p) and 8K UHD (4320p), which are two digital video formats proposed by NHK Science & Technology Research Laboratories and defined and approved by the International Telecommunication Union (ITU). The Consumer Electronics Association announced on 17 October 2012, that "Ultra High Definition", or "Ultra HD", would be used for displays that have an aspect ratio of at least 16:9 and at least one digital input capable of carrying and presenting native video at a minimum resolution of 3840×2160 pixels.[191][192]
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North American consumers purchase a new television set on average every seven years, and the average household owns 2.8 televisions. As of 2011[update], 48 million are sold each year at an average price of $460 and size of 38 in (97 cm).[193]
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Getting TV programming shown to the public can happen in many different ways. After production, the next step is to market and deliver the product to whichever markets are open to using it. This typically happens on two levels:
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First-run programming is increasing on subscription services outside the US, but few domestically produced programs are syndicated on domestic free-to-air (FTA) elsewhere. This practice is increasing, however, generally on digital-only FTA channels or with subscriber-only, first-run material appearing on FTA. Unlike the US, repeat FTA screenings of an FTA network program usually only occur on that network. Also, affiliates rarely buy or produce non-network programming that is not centered on local programming.
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Television genres include a broad range of programming types that entertain, inform, and educate viewers. The most expensive entertainment genres to produce are usually dramas and dramatic miniseries. However, other genres, such as historical Western genres, may also have high production costs.
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Popular culture entertainment genres include action-oriented shows such as police, crime, detective dramas, horror, or thriller shows. As well, there are also other variants of the drama genre, such as medical dramas and daytime soap operas. Science fiction shows can fall into either the drama or action category, depending on whether they emphasize philosophical questions or high adventure. Comedy is a popular genre which includes situation comedy (sitcom) and animated shows for the adult demographic such as South Park.
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The least expensive forms of entertainment programming genres are game shows, talk shows, variety shows, and reality television. Game shows feature contestants answering questions and solving puzzles to win prizes. Talk shows contain interviews with film, television, music and sports celebrities and public figures. Variety shows feature a range of musical performers and other entertainers, such as comedians and magicians, introduced by a host or Master of Ceremonies. There is some crossover between some talk shows and variety shows because leading talk shows often feature performances by bands, singers, comedians, and other performers in between the interview segments. Reality TV shows "regular" people (i.e., not actors) facing unusual challenges or experiences ranging from arrest by police officers (COPS) to significant weight loss (The Biggest Loser). A variant version of reality shows depicts celebrities doing mundane activities such as going about their everyday life (The Osbournes, Snoop Dogg's Father Hood) or doing regular jobs (The Simple Life).
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Fictional television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include series such as Twin Peaks and The Sopranos. Kristin Thompson argues that some of these television series exhibit traits also found in art films, such as psychological realism, narrative complexity, and ambiguous plotlines. Nonfiction television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include a range of serious, noncommercial, programming aimed at a niche audience, such as documentaries and public affairs shows.
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Around the globe, broadcast TV is financed by government, advertising, licensing (a form of tax), subscription, or any combination of these. To protect revenues, subscription TV channels are usually encrypted to ensure that only subscribers receive the decryption codes to see the signal. Unencrypted channels are known as free to air or FTA. In 2009, the global TV market represented 1,217.2 million TV households with at least one TV and total revenues of 268.9 billion EUR (declining 1.2% compared to 2008).[195] North America had the biggest TV revenue market share with 39% followed by Europe (31%), Asia-Pacific (21%), Latin America (8%), and Africa and the Middle East (2%).[196] Globally, the different TV revenue sources divide into 45–50% TV advertising revenues, 40–45% subscription fees and 10% public funding.[197][198]
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TV's broad reach makes it a powerful and attractive medium for advertisers. Many TV networks and stations sell blocks of broadcast time to advertisers ("sponsors") to fund their programming.[199] Television advertisements (variously called a television commercial, commercial or ad in American English, and known in British English as an advert) is a span of television programming produced and paid for by an organization, which conveys a message, typically to market a product or service. Advertising revenue provides a significant portion of the funding for most privately owned television networks. The vast majority of television advertisements today consist of brief advertising spots, ranging in length from a few seconds to several minutes (as well as program-length infomercials). Advertisements of this sort have been used to promote a wide variety of goods, services and ideas since the beginning of television.
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The effects of television advertising upon the viewing public (and the effects of mass media in general) have been the subject of philosophical discourse by such luminaries as Marshall McLuhan. The viewership of television programming, as measured by companies such as Nielsen Media Research, is often used as a metric for television advertisement placement, and consequently, for the rates charged to advertisers to air within a given network, television program, or time of day (called a "daypart"). In many countries, including the United States, television campaign advertisements are considered indispensable for a political campaign. In other countries, such as France, political advertising on television is heavily restricted,[200] while some countries, such as Norway, completely ban political advertisements.
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The first official, paid television advertisement was broadcast in the United States on 1 July 1941 over New York station WNBT (now WNBC) before a baseball game between the Brooklyn Dodgers and Philadelphia Phillies. The announcement for Bulova watches, for which the company paid anywhere from $4.00 to $9.00 (reports vary), displayed a WNBT test pattern modified to look like a clock with the hands showing the time. The Bulova logo, with the phrase "Bulova Watch Time", was shown in the lower right-hand quadrant of the test pattern while the second hand swept around the dial for one minute.[201][202] The first TV ad broadcast in the UK was on ITV on 22 September 1955, advertising Gibbs SR toothpaste. The first TV ad broadcast in Asia was on Nippon Television in Tokyo on 28 August 1953, advertising Seikosha (now Seiko), which also displayed a clock with the current time.[203]
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Since inception in the US in 1941,[204] television commercials have become one of the most effective, persuasive, and popular methods of selling products of many sorts, especially consumer goods. During the 1940s and into the 1950s, programs were hosted by single advertisers. This, in turn, gave great creative license to the advertisers over the content of the show. Perhaps due to the quiz show scandals in the 1950s,[205] networks shifted to the magazine concept, introducing advertising breaks with multiple advertisers.
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US advertising rates are determined primarily by Nielsen ratings. The time of the day and popularity of the channel determine how much a TV commercial can cost. For example, it can cost approximately $750,000 for a 30-second block of commercial time during the highly popular American Idol, while the same amount of time for the Super Bowl can cost several million dollars. Conversely, lesser-viewed time slots, such as early mornings and weekday afternoons, are often sold in bulk to producers of infomercials at far lower rates. In recent years, the paid program or infomercial has become common, usually in lengths of 30 minutes or one hour. Some drug companies and other businesses have even created "news" items for broadcast, known in the industry as video news releases, paying program directors to use them.[206]
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Some TV programs also deliberately place products into their shows as advertisements, a practice started in feature films[207] and known as product placement. For example, a character could be drinking a certain kind of soda, going to a particular chain restaurant, or driving a certain make of car. (This is sometimes very subtle, with shows having vehicles provided by manufacturers for low cost in exchange as a product placement). Sometimes, a specific brand or trade mark, or music from a certain artist or group, is used. (This excludes guest appearances by artists who perform on the show.)
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The TV regulator oversees TV advertising in the United Kingdom. Its restrictions have applied since the early days of commercially funded TV. Despite this, an early TV mogul, Roy Thomson, likened the broadcasting licence as being a "licence to print money".[208] Restrictions mean that the big three national commercial TV channels: ITV, Channel 4, and Channel 5 can show an average of only seven minutes of advertising per hour (eight minutes in the peak period). Other broadcasters must average no more than nine minutes (twelve in the peak). This means that many imported TV shows from the US have unnatural pauses where the UK company does not utilize the narrative breaks intended for more frequent US advertising. Advertisements must not be inserted in the course of certain specific proscribed types of programs which last less than half an hour in scheduled duration; this list includes any news or current affairs programs, documentaries, and programs for children; additionally, advertisements may not be carried in a program designed and broadcast for reception in schools or in any religious broadcasting service or other devotional program or during a formal Royal ceremony or occasion. There also must be clear demarcations in time between the programs and the advertisements. The BBC, being strictly non-commercial, is not allowed to show advertisements on television in the UK, although it has many advertising-funded channels abroad. The majority of its budget comes from television license fees (see below) and broadcast syndication, the sale of content to other broadcasters.
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Broadcast advertising is regulated by the Broadcasting Authority of Ireland.[209]
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Some TV channels are partly funded from subscriptions; therefore, the signals are encrypted during broadcast to ensure that only the paying subscribers have access to the decryption codes to watch pay television or specialty channels. Most subscription services are also funded by advertising.
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Television services in some countries may be funded by a television licence or a form of taxation, which means that advertising plays a lesser role or no role at all. For example, some channels may carry no advertising at all and some very little, including:
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The BBC carries no television advertising on its UK channels and is funded by an annual television licence paid by premises receiving live TV broadcasts. Currently, it is estimated that approximately 26.8 million UK private domestic households own televisions, with approximately 25 million TV licences in all premises in force as of 2010.[210] This television license fee is set by the government, but the BBC is not answerable to or controlled by the government.
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The two main BBC TV channels are watched by almost 90% of the population each week and overall have 27% share of total viewing,[211] despite the fact that 85% of homes are multichannel, with 42% of these having access to 200 free to air channels via satellite and another 43% having access to 30 or more channels via Freeview.[212] The licence that funds the seven advertising-free BBC TV channels costs £147 a year (about US$200) as of 2018 regardless of the number of TV sets owned; the price is reduced by two-thirds if only black and white television is received.[213] When the same sporting event has been presented on both BBC and commercial channels, the BBC always attracts the lion's share of the audience, indicating that viewers prefer to watch TV uninterrupted by advertising.
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Other than internal promotional material, the Australian Broadcasting Corporation (ABC) carries no advertising; it is banned under the ABC Act 1983. The ABC receives its funding from the Australian government every three years. In the 2014/15 federal budget, the ABC received A$1.11 billion.[214] The funds provide for the ABC's television, radio, online, and international outputs. The ABC also receives funds from its many ABC shops across Australia. Although funded by the Australian government, the editorial independence of the ABC is ensured through law.
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In France, government-funded channels carry advertisements, yet those who own television sets have to pay an annual tax ("la redevance audiovisuelle").[215]
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In Japan, NHK is paid for by license fees (known in Japanese as reception fee (受信料, Jushinryō)). The broadcast law that governs NHK's funding stipulates that any television equipped to receive NHK is required to pay. The fee is standardized, with discounts for office workers and students who commute, as well a general discount for residents of Okinawa prefecture.
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Broadcast programming, or TV listings in the United Kingdom, is the practice of organizing television programs in a schedule, with broadcast automation used to regularly change the scheduling of TV programs to build an audience for a new show, retain that audience, or compete with other broadcasters' programs.
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Television has played a pivotal role in the socialization of the 20th and 21st centuries. There are many aspects of television that can be addressed, including negative issues such as media violence. Current research is discovering that individuals suffering from social isolation can employ television to create what is termed a parasocial or faux relationship with characters from their favorite television shows and movies as a way of deflecting feelings of loneliness and social deprivation.[216] Several studies have found that educational television has many advantages. The article "The Good Things about Television"[217] argues that television can be a very powerful and effective learning tool for children if used wisely.
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Methodist denominations in the conservative holiness movement, such as the Allegheny Wesleyan Methodist Connection and the Evangelical Wesleyan Church, eschew the use of the television.[218]
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Children, especially those aged 5 or younger, are at risk of injury from falling televisions.[219] A CRT-style television that falls on a child will, because of its weight, hit with the equivalent force of falling multiple stories from a building.[220] Newer flat-screen televisions are "top-heavy and have narrow bases", which means that a small child can easily pull one over.[221] As of 2015[update], TV tip-overs were responsible for more than 10,000 injuries per year to children, at a cost of more than $8 million per year in emergency care.[219][221]
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A 2017 study in The Journal of Human Resources found that exposure to cable television reduced cognitive ability and high school graduation rates for boys. This effect was stronger for boys from more educated families. The article suggests a mechanism where light television entertainment crowds out more cognitively stimulating activities.[222]
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With high lead content in CRTs and the rapid diffusion of new flat-panel display technologies, some of which (LCDs) use lamps which contain mercury, there is growing concern about electronic waste from discarded televisions. Related occupational health concerns exist, as well, for disassemblers removing copper wiring and other materials from CRTs. Further environmental concerns related to television design and use relate to the devices' increasing electrical energy requirements.[223]
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Television (TV), sometimes shortened to tele or telly, is a telecommunication medium used for transmitting moving images in monochrome (black and white), or in color, and in two or three dimensions and sound. The term can refer to a television set, a television show, or the medium of television transmission. Television is a mass medium for advertising, entertainment, news, and sports.
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Television became available in crude experimental forms in the late 1920s, but it would still be several years before the new technology would be marketed to consumers. After World War II, an improved form of black-and-white TV broadcasting became popular in the United States and Britain, and television sets became commonplace in homes, businesses, and institutions. During the 1950s, television was the primary medium for influencing public opinion.[1] In the mid-1960s, color broadcasting was introduced in the US and most other developed countries. The availability of multiple types of archival storage media such as Betamax and VHS tapes, high-capacity hard disk drives, DVDs, flash drives, high-definition Blu-ray Discs, and cloud digital video recorders has enabled viewers to watch pre-recorded material—such as movies—at home on their own time schedule. For many reasons, especially the convenience of remote retrieval, the storage of television and video programming now occurs on the cloud. At the end of the first decade of the 2000s, digital television transmissions greatly increased in popularity. Another development was the move from standard-definition television (SDTV) (576i, with 576 interlaced lines of resolution and 480i) to high-definition television (HDTV), which provides a resolution that is substantially higher. HDTV may be transmitted in various formats: 1080p, 1080i and 720p. Since 2010, with the invention of smart television, Internet television has increased the availability of television programs and movies via the Internet through streaming video services such as Netflix, Amazon Video, iPlayer and Hulu.
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In 2013, 79% of the world's households owned a television set.[2] The replacement of early bulky, high-voltage cathode ray tube (CRT) screen displays with compact, energy-efficient, flat-panel alternative technologies such as LCDs (both fluorescent-backlit and LED), OLED displays, and plasma displays was a hardware revolution that began with computer monitors in the late 1990s. Most TV sets sold in the 2000s were flat-panel, mainly LEDs. Major manufacturers announced the discontinuation of CRT, DLP, plasma, and even fluorescent-backlit LCDs by the mid-2010s.[3][4] In the near future, LEDs are expected to be gradually replaced by OLEDs.[5] Also, major manufacturers have announced that they will increasingly produce smart TVs in the mid-2010s.[6][7][8] Smart TVs with integrated Internet and Web 2.0 functions became the dominant form of television by the late 2010s.[9]
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Television signals were initially distributed only as terrestrial television using high-powered radio-frequency transmitters to broadcast the signal to individual television receivers. Alternatively television signals are distributed by coaxial cable or optical fiber, satellite systems and, since the 2000s via the Internet. Until the early 2000s, these were transmitted as analog signals, but a transition to digital television is expected to be completed worldwide by the late 2010s. A standard television set is composed of multiple internal electronic circuits, including a tuner for receiving and decoding broadcast signals. A visual display device which lacks a tuner is correctly called a video monitor rather than a television.
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The word television comes from Ancient Greek τῆλε (tèle), meaning 'far', and Latin visio, meaning 'sight'. The first documented usage of the term dates back to 1900, when the Russian scientist Constantin Perskyi used it in a paper that he presented in French at the 1st International Congress of Electricity, which ran from 18 to 25 August 1900 during the International World Fair in Paris. The Anglicised version of the term is first attested in 1907, when it was still "...a theoretical system to transmit moving images over telegraph or telephone wires".[10] It was "...formed in English or borrowed from French télévision."[10] In the 19th century and early 20th century, other "...proposals for the name of a then-hypothetical technology for sending pictures over distance were telephote (1880) and televista (1904)."[10] The abbreviation "TV" is from 1948. The use of the term to mean "a television set" dates from 1941.[10] The use of the term to mean "television as a medium" dates from 1927.[10] The slang term "telly" is more common in the UK. The slang term "the tube" or the "boob tube" derives from the bulky cathode ray tube used on most TVs until the advent of flat-screen TVs. Another slang term for the TV is "idiot box".[11] Also, in the 1940s and throughout the 1950s, during the early rapid growth of television programming and television-set ownership in the United States, another slang term became widely used in that period and continues to be used today to distinguish productions originally created for broadcast on television from films developed for presentation in movie theaters.[12] The "small screen", as both a compound adjective and noun, became specific references to television, while the "big screen" was used to identify productions made for theatrical release.[12]
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Facsimile transmission systems for still photographs pioneered methods of mechanical scanning of images in the early 19th century. Alexander Bain introduced the facsimile machine between 1843 and 1846. Frederick Bakewell demonstrated a working laboratory version in 1851.[citation needed] Willoughby Smith discovered the photoconductivity of the element selenium in 1873. As a 23-year-old German university student, Paul Julius Gottlieb Nipkow proposed and patented the Nipkow disk in 1884.[13] This was a spinning disk with a spiral pattern of holes in it, so each hole scanned a line of the image. Although he never built a working model of the system, variations of Nipkow's spinning-disk "image rasterizer" became exceedingly common.[14] Constantin Perskyi had coined the word television in a paper read to the International Electricity Congress at the International World Fair in Paris on 24 August 1900. Perskyi's paper reviewed the existing electromechanical technologies, mentioning the work of Nipkow and others.[15] However, it was not until 1907 that developments in amplification tube technology by Lee de Forest and Arthur Korn, among others, made the design practical.[16]
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The first demonstration of the live transmission of images was by Georges Rignoux and A. Fournier in Paris in 1909. A matrix of 64 selenium cells, individually wired to a mechanical commutator, served as an electronic retina. In the receiver, a type of Kerr cell modulated the light and a series of variously angled mirrors attached to the edge of a rotating disc scanned the modulated beam onto the display screen. A separate circuit regulated synchronization. The 8x8 pixel resolution in this proof-of-concept demonstration was just sufficient to clearly transmit individual letters of the alphabet. An updated image was transmitted "several times" each second.[17]
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In 1911, Boris Rosing and his student Vladimir Zworykin created a system that used a mechanical mirror-drum scanner to transmit, in Zworykin's words, "very crude images" over wires to the "Braun tube" (cathode ray tube or "CRT") in the receiver. Moving images were not possible because, in the scanner: "the sensitivity was not enough and the selenium cell was very laggy".[18]
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In 1921, Edouard Belin sent the first image via radio waves with his belinograph.[citation needed]
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By the 1920s, when amplification made television practical, Scottish inventor John Logie Baird employed the Nipkow disk in his prototype video systems. On 25 March 1925, Baird gave the first public demonstration of televised silhouette images in motion, at Selfridge's Department Store in London.[19] Since human faces had inadequate contrast to show up on his primitive system, he televised a ventriloquist's dummy named "Stooky Bill", whose painted face had higher contrast, talking and moving. By 26 January 1926, he demonstrated the transmission of the image of a face in motion by radio. This is widely regarded as the first television demonstration. The subject was Baird's business partner Oliver Hutchinson. Baird's system used the Nipkow disk for both scanning the image and displaying it. A bright light shining through a spinning Nipkow disk set with lenses projected a bright spot of light which swept across the subject. A Selenium photoelectric tube detected the light reflected from the subject and converted it into a proportional electrical signal. This was transmitted by AM radio waves to a receiver unit, where the video signal was applied to a neon light behind a second Nipkow disk rotating synchronized with the first. The brightness of the neon lamp was varied in proportion to the brightness of each spot on the image. As each hole in the disk passed by, one scan line of the image was reproduced. Baird's disk had 30 holes, producing an image with only 30 scan lines, just enough to recognize a human face. In 1927, Baird transmitted a signal over 438 miles (705 km) of telephone line between London and Glasgow.[citation needed]
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In 1928, Baird's company (Baird Television Development Company/Cinema Television) broadcast the first transatlantic television signal, between London and New York, and the first shore-to-ship transmission. In 1929, he became involved in the first experimental mechanical television service in Germany. In November of the same year, Baird and Bernard Natan of Pathé established France's first television company, Télévision-Baird-Natan. In 1931, he made the first outdoor remote broadcast, of The Derby.[20] In 1932, he demonstrated ultra-short wave television. Baird's mechanical system reached a peak of 240-lines of resolution on BBC television broadcasts in 1936, though the mechanical system did not scan the televised scene directly. Instead a 17.5mm film was shot, rapidly developed and then scanned while the film was still wet.[citation needed]
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An American inventor, Charles Francis Jenkins, also pioneered the television. He published an article on "Motion Pictures by Wireless" in 1913, but it was not until December 1923 that he transmitted moving silhouette images for witnesses; and it was on 13 June 1925, that he publicly demonstrated synchronized transmission of silhouette pictures. In 1925 Jenkins used the Nipkow disk and transmitted the silhouette image of a toy windmill in motion, over a distance of 5 miles (8 km), from a naval radio station in Maryland to his laboratory in Washington, D.C., using a lensed disk scanner with a 48-line resolution.[21][22] He was granted U.S. Patent No. 1,544,156 (Transmitting Pictures over Wireless) on 30 June 1925 (filed 13 March 1922).[23]
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Herbert E. Ives and Frank Gray of Bell Telephone Laboratories gave a dramatic demonstration of mechanical television on 7 April 1927. Their reflected-light television system included both small and large viewing screens. The small receiver had a 2-inch-wide by 2.5-inch-high screen (5 by 6 cm). The large receiver had a screen 24 inches wide by 30 inches high (60 by 75 cm). Both sets were capable of reproducing reasonably accurate, monochromatic, moving images. Along with the pictures, the sets received synchronized sound. The system transmitted images over two paths: first, a copper wire link from Washington to New York City, then a radio link from Whippany, New Jersey. Comparing the two transmission methods, viewers noted no difference in quality. Subjects of the telecast included Secretary of Commerce Herbert Hoover. A flying-spot scanner beam illuminated these subjects. The scanner that produced the beam had a 50-aperture disk. The disc revolved at a rate of 18 frames per second, capturing one frame about every 56 milliseconds. (Today's systems typically transmit 30 or 60 frames per second, or one frame every 33.3 or 16.7 milliseconds respectively.) Television historian Albert Abramson underscored the significance of the Bell Labs demonstration: "It was in fact the best demonstration of a mechanical television system ever made to this time. It would be several years before any other system could even begin to compare with it in picture quality."[24]
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In 1928, WRGB, then W2XB, was started as the world's first television station. It broadcast from the General Electric facility in Schenectady, NY. It was popularly known as "WGY Television". Meanwhile, in the Soviet Union, Léon Theremin had been developing a mirror drum-based television, starting with 16 lines resolution in 1925, then 32 lines and eventually 64 using interlacing in 1926. As part of his thesis, on 7 May 1926, he electrically transmitted, and then projected, near-simultaneous moving images on a 5-square-foot (0.46 m2) screen.[22]
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By 1927, Theremin had achieved an image of 100 lines, a resolution that was not surpassed until May 1932 by RCA, with 120 lines.[25]
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On 25 December 1926, Kenjiro Takayanagi demonstrated a television system with a 40-line resolution that employed a Nipkow disk scanner and CRT display at Hamamatsu Industrial High School in Japan. This prototype is still on display at the Takayanagi Memorial Museum in Shizuoka University, Hamamatsu Campus. His research in creating a production model was halted by the SCAP after World War II.[26]
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Because only a limited number of holes could be made in the disks, and disks beyond a certain diameter became impractical, image resolution on mechanical television broadcasts was relatively low, ranging from about 30 lines up to 120 or so. Nevertheless, the image quality of 30-line transmissions steadily improved with technical advances, and by 1933 the UK broadcasts using the Baird system were remarkably clear.[27] A few systems ranging into the 200-line region also went on the air. Two of these were the 180-line system that Compagnie des Compteurs (CDC) installed in Paris in 1935, and the 180-line system that Peck Television Corp. started in 1935 at station VE9AK in Montreal.[28][29] The advancement of all-electronic television (including image dissectors and other camera tubes and cathode ray tubes for the reproducer) marked the beginning of the end for mechanical systems as the dominant form of television. Mechanical television, despite its inferior image quality and generally smaller picture, would remain the primary television technology until the 1930s. The last mechanical television broadcasts ended in 1939 at stations run by a handful of public universities in the United States.[citation needed]
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In 1897, English physicist J. J. Thomson was able, in his three famous experiments, to deflect cathode rays, a fundamental function of the modern cathode ray tube (CRT). The earliest version of the CRT was invented by the German physicist Ferdinand Braun in 1897 and is also known as the "Braun" tube.[30][31] It was a cold-cathode diode, a modification of the Crookes tube, with a phosphor-coated screen. In 1906 the Germans Max Dieckmann and Gustav Glage produced raster images for the first time in a CRT.[32] In 1907, Russian scientist Boris Rosing used a CRT in the receiving end of an experimental video signal to form a picture. He managed to display simple geometric shapes onto the screen.[33]
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In 1908 Alan Archibald Campbell-Swinton, fellow of the Royal Society (UK), published a letter in the scientific journal Nature in which he described how "distant electric vision" could be achieved by using a cathode ray tube, or Braun tube, as both a transmitting and receiving device,[34][35] He expanded on his vision in a speech given in London in 1911 and reported in The Times[36] and the Journal of the Röntgen Society.[37][38] In a letter to Nature published in October 1926, Campbell-Swinton also announced the results of some "not very successful experiments" he had conducted with G. M. Minchin and J. C. M. Stanton. They had attempted to generate an electrical signal by projecting an image onto a selenium-coated metal plate that was simultaneously scanned by a cathode ray beam.[39][40] These experiments were conducted before March 1914, when Minchin died,[41] but they were later repeated by two different teams in 1937, by H. Miller and J. W. Strange from EMI,[42] and by H. Iams and A. Rose from RCA.[43] Both teams succeeded in transmitting "very faint" images with the original Campbell-Swinton's selenium-coated plate. Although others had experimented with using a cathode ray tube as a receiver, the concept of using one as a transmitter was novel.[44] The first cathode ray tube to use a hot cathode was developed by John B. Johnson (who gave his name to the term Johnson noise) and Harry Weiner Weinhart of Western Electric, and became a commercial product in 1922.[citation needed]
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In 1926, Hungarian engineer Kálmán Tihanyi designed a television system utilizing fully electronic scanning and display elements and employing the principle of "charge storage" within the scanning (or "camera") tube.[45][46][47][48] The problem of low sensitivity to light resulting in low electrical output from transmitting or "camera" tubes would be solved with the introduction of charge-storage technology by Kálmán Tihanyi beginning in 1924.[49] His solution was a camera tube that accumulated and stored electrical charges ("photoelectrons") within the tube throughout each scanning cycle. The device was first described in a patent application he filed in Hungary in March 1926 for a television system he dubbed "Radioskop".[50] After further refinements included in a 1928 patent application,[49] Tihanyi's patent was declared void in Great Britain in 1930,[51] so he applied for patents in the United States. Although his breakthrough would be incorporated into the design of RCA's "iconoscope" in 1931, the U.S. patent for Tihanyi's transmitting tube would not be granted until May 1939. The patent for his receiving tube had been granted the previous October. Both patents had been purchased by RCA prior to their approval.[52][53] Charge storage remains a basic principle in the design of imaging devices for television to the present day.[50] On 25 December 1926, at Hamamatsu Industrial High School in Japan, Japanese inventor Kenjiro Takayanagi demonstrated a TV system with a 40-line resolution that employed a CRT display.[26] This was the first working example of a fully electronic television receiver. Takayanagi did not apply for a patent.[54]
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On 7 September 1927, American inventor Philo Farnsworth's image dissector camera tube transmitted its first image, a simple straight line, at his laboratory at 202 Green Street in San Francisco.[55][56] By 3 September 1928, Farnsworth had developed the system sufficiently to hold a demonstration for the press. This is widely regarded as the first electronic television demonstration.[56] In 1929, the system was improved further by the elimination of a motor generator, so that his television system now had no mechanical parts.[57] That year, Farnsworth transmitted the first live human images with his system, including a three and a half-inch image of his wife Elma ("Pem") with her eyes closed (possibly due to the bright lighting required).[58]
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Meanwhile, Vladimir Zworykin was also experimenting with the cathode ray tube to create and show images. While working for Westinghouse Electric in 1923, he began to develop an electronic camera tube. But in a 1925 demonstration, the image was dim, had low contrast, and poor definition, and was stationary.[59] Zworykin's imaging tube never got beyond the laboratory stage. But RCA, which acquired the Westinghouse patent, asserted that the patent for Farnsworth's 1927 image dissector was written so broadly that it would exclude any other electronic imaging device. Thus RCA, on the basis of Zworykin's 1923 patent application, filed a patent interference suit against Farnsworth. The U.S. Patent Office examiner disagreed in a 1935 decision, finding priority of invention for Farnsworth against Zworykin. Farnsworth claimed that Zworykin's 1923 system would be unable to produce an electrical image of the type to challenge his patent. Zworykin received a patent in 1928 for a color transmission version of his 1923 patent application;[60] he also divided his original application in 1931.[61] Zworykin was unable or unwilling to introduce evidence of a working model of his tube that was based on his 1923 patent application. In September 1939, after losing an appeal in the courts, and determined to go forward with the commercial manufacturing of television equipment, RCA agreed to pay Farnsworth US$1 million over a ten-year period, in addition to license payments, to use his patents.[62][63]
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In 1933, RCA introduced an improved camera tube that relied on Tihanyi's charge storage principle.[64] Dubbed the "Iconoscope" by Zworykin, the new tube had a light sensitivity of about 75,000 lux, and thus was claimed to be much more sensitive than Farnsworth's image dissector.[citation needed] However, Farnsworth had overcome his power problems with his Image Dissector through the invention of a completely unique "multipactor" device that he began work on in 1930, and demonstrated in 1931.[65][66] This small tube could amplify a signal reportedly to the 60th power or better[67] and showed great promise in all fields of electronics. Unfortunately, a problem with the multipactor was that it wore out at an unsatisfactory rate.[68]
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At the Berlin Radio Show in August 1931, Manfred von Ardenne gave a public demonstration of a television system using a CRT for both transmission and reception. However, Ardenne had not developed a camera tube, using the CRT instead as a flying-spot scanner to scan slides and film.[69] Philo Farnsworth gave the world's first public demonstration of an all-electronic television system, using a live camera, at the Franklin Institute of Philadelphia on 25 August 1934, and for ten days afterwards.[70][71] Mexican inventor Guillermo González Camarena also played an important role in early TV. His experiments with TV (known as telectroescopía at first) began in 1931 and led to a patent for the "trichromatic field sequential system" color television in 1940.[72] In Britain, the EMI engineering team led by Isaac Shoenberg applied in 1932 for a patent for a new device they dubbed "the Emitron",[73][74] which formed the heart of the cameras they designed for the BBC. On 2 November 1936, a 405-line broadcasting service employing the Emitron began at studios in Alexandra Palace, and transmitted from a specially built mast atop one of the Victorian building's towers. It alternated for a short time with Baird's mechanical system in adjoining studios, but was more reliable and visibly superior. This was the world's first regular "high-definition" television service.[75]
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The original American iconoscope was noisy, had a high ratio of interference to signal, and ultimately gave disappointing results, especially when compared to the high definition mechanical scanning systems then becoming available.[76][77] The EMI team, under the supervision of Isaac Shoenberg, analyzed how the iconoscope (or Emitron) produces an electronic signal and concluded that its real efficiency was only about 5% of the theoretical maximum.[78][79] They solved this problem by developing, and patenting in 1934, two new camera tubes dubbed super-Emitron and CPS Emitron.[80][81][82] The super-Emitron was between ten and fifteen times more sensitive than the original Emitron and iconoscope tubes and, in some cases, this ratio was considerably greater.[78] It was used for outside broadcasting by the BBC, for the first time, on Armistice Day 1937, when the general public could watch on a television set as the King laid a wreath at the Cenotaph.[83] This was the first time that anyone had broadcast a live street scene from cameras installed on the roof of neighboring buildings, because neither Farnsworth nor RCA would do the same until the 1939 New York World's Fair.
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On the other hand, in 1934, Zworykin shared some patent rights with the German licensee company Telefunken.[84] The "image iconoscope" ("Superikonoskop" in Germany) was produced as a result of the collaboration. This tube is essentially identical to the super-Emitron.[citation needed] The production and commercialization of the super-Emitron and image iconoscope in Europe were not affected by the patent war between Zworykin and Farnsworth, because Dieckmann and Hell had priority in Germany for the invention of the image dissector, having submitted a patent application for their Lichtelektrische Bildzerlegerröhre für Fernseher (Photoelectric Image Dissector Tube for Television) in Germany in 1925,[85] two years before Farnsworth did the same in the United States.[86] The image iconoscope (Superikonoskop) became the industrial standard for public broadcasting in Europe from 1936 until 1960, when it was replaced by the vidicon and plumbicon tubes. Indeed, it was the representative of the European tradition in electronic tubes competing against the American tradition represented by the image orthicon.[87][88] The German company Heimann produced the Superikonoskop for the 1936 Berlin Olympic Games,[89][90] later Heimann also produced and commercialized it from 1940 to 1955;[91] finally the Dutch company Philips produced and commercialized the image iconoscope and multicon from 1952 to 1958.[88][92]
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American television broadcasting, at the time, consisted of a variety of markets in a wide range of sizes, each competing for programming and dominance with separate technology, until deals were made and standards agreed upon in 1941.[93] RCA, for example, used only Iconoscopes in the New York area, but Farnsworth Image Dissectors in Philadelphia and San Francisco.[94] In September 1939, RCA agreed to pay the Farnsworth Television and Radio Corporation royalties over the next ten years for access to Farnsworth's patents.[95] With this historic agreement in place, RCA integrated much of what was best about the Farnsworth Technology into their systems.[94] In 1941, the United States implemented 525-line television.[96][97] Electrical engineer Benjamin Adler played a prominent role in the development of television.[98][99]
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The world's first 625-line television standard was designed in the Soviet Union in 1944 and became a national standard in 1946.[100] The first broadcast in 625-line standard occurred in Moscow in 1948.[101] The concept of 625 lines per frame was subsequently implemented in the European CCIR standard.[102] In 1936, Kálmán Tihanyi described the principle of plasma display, the first flat panel display system.[103][104]
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Early electronic television sets were large and bulky, with analog circuits made of vacuum tubes. Following the invention of the first working transistor at Bell Labs, Sony founder Masaru Ibuka predicted in 1952 that the transition to electronic circuits made of transistors would lead to smaller and more portable television sets.[105] The first fully transistorized, portable solid-state television set was the 8-inch Sony TV8-301, developed in 1959 and released in 1960.[106][107] This began the transformation of television viewership from a communal viewing experience to a solitary viewing experience.[108] By 1960, Sony had sold over 4 million portable television sets worldwide.[109]
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The basic idea of using three monochrome images to produce a color image had been experimented with almost as soon as black-and-white televisions had first been built. Although he gave no practical details, among the earliest published proposals for television was one by Maurice Le Blanc, in 1880, for a color system, including the first mentions in television literature of line and frame scanning.[110] Polish inventor Jan Szczepanik patented a color television system in 1897, using a selenium photoelectric cell at the transmitter and an electromagnet controlling an oscillating mirror and a moving prism at the receiver. But his system contained no means of analyzing the spectrum of colors at the transmitting end, and could not have worked as he described it.[111] Another inventor, Hovannes Adamian, also experimented with color television as early as 1907. The first color television project is claimed by him,[112] and was patented in Germany on 31 March 1908, patent No. 197183, then in Britain, on 1 April 1908, patent No. 7219,[113] in France (patent No. 390326) and in Russia in 1910 (patent No. 17912).[114]
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Scottish inventor John Logie Baird demonstrated the world's first color transmission on 3 July 1928, using scanning discs at the transmitting and receiving ends with three spirals of apertures, each spiral with filters of a different primary color; and three light sources at the receiving end, with a commutator to alternate their illumination.[115] Baird also made the world's first color broadcast on 4 February 1938, sending a mechanically scanned 120-line image from Baird's Crystal Palace studios to a projection screen at London's Dominion Theatre.[116] Mechanically scanned color television was also demonstrated by Bell Laboratories in June 1929 using three complete systems of photoelectric cells, amplifiers, glow-tubes, and color filters, with a series of mirrors to superimpose the red, green, and blue images into one full color image.
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The first practical hybrid system was again pioneered by John Logie Baird. In 1940 he publicly demonstrated a color television combining a traditional black-and-white display with a rotating colored disk. This device was very "deep", but was later improved with a mirror folding the light path into an entirely practical device resembling a large conventional console.[117] However, Baird was not happy with the design, and, as early as 1944, had commented to a British government committee that a fully electronic device would be better.
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In 1939, Hungarian engineer Peter Carl Goldmark introduced an electro-mechanical system while at CBS, which contained an Iconoscope sensor. The CBS field-sequential color system was partly mechanical, with a disc made of red, blue, and green filters spinning inside the television camera at 1,200 rpm, and a similar disc spinning in synchronization in front of the cathode ray tube inside the receiver set.[118] The system was first demonstrated to the Federal Communications Commission (FCC) on 29 August 1940, and shown to the press on 4 September.[119][120][121][122]
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CBS began experimental color field tests using film as early as 28 August 1940, and live cameras by 12 November.[120][123] NBC (owned by RCA) made its first field test of color television on 20 February 1941. CBS began daily color field tests on 1 June 1941.[124] These color systems were not compatible with existing black-and-white television sets, and, as no color television sets were available to the public at this time, viewing of the color field tests was restricted to RCA and CBS engineers and the invited press. The War Production Board halted the manufacture of television and radio equipment for civilian use from 22 April 1942 to 20 August 1945, limiting any opportunity to introduce color television to the general public.[125][126]
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As early as 1940, Baird had started work on a fully electronic system he called Telechrome. Early Telechrome devices used two electron guns aimed at either side of a phosphor plate. The phosphor was patterned so the electrons from the guns only fell on one side of the patterning or the other. Using cyan and magenta phosphors, a reasonable limited-color image could be obtained. He also demonstrated the same system using monochrome signals to produce a 3D image (called "stereoscopic" at the time). A demonstration on 16 August 1944 was the first example of a practical color television system. Work on the Telechrome continued and plans were made to introduce a three-gun version for full color. However, Baird's untimely death in 1946 ended development of the Telechrome system.[127][128]
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Similar concepts were common through the 1940s and 1950s, differing primarily in the way they re-combined the colors generated by the three guns. The Geer tube was similar to Baird's concept, but used small pyramids with the phosphors deposited on their outside faces, instead of Baird's 3D patterning on a flat surface. The Penetron used three layers of phosphor on top of each other and increased the power of the beam to reach the upper layers when drawing those colors. The Chromatron used a set of focusing wires to select the colored phosphors arranged in vertical stripes on the tube.
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One of the great technical challenges of introducing color broadcast television was the desire to conserve bandwidth, potentially three times that of the existing black-and-white standards, and not use an excessive amount of radio spectrum. In the United States, after considerable research, the National Television Systems Committee[129] approved an all-electronic system developed by RCA, which encoded the color information separately from the brightness information and greatly reduced the resolution of the color information in order to conserve bandwidth. As black-and-white TVs could receive the same transmission and display it in black-and-white, the color system adopted is [backwards] "compatible". ("Compatible Color", featured in RCA advertisements of the period, is mentioned in the song "America", of West Side Story, 1957.) The brightness image remained compatible with existing black-and-white television sets at slightly reduced resolution, while color televisions could decode the extra information in the signal and produce a limited-resolution color display. The higher resolution black-and-white and lower resolution color images combine in the brain to produce a seemingly high-resolution color image. The NTSC standard represented a major technical achievement.
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The first color broadcast (the first episode of the live program The Marriage (TV series)) occurred on 8 July 1954, but during the following ten years most network broadcasts, and nearly all local programming, continued to be in black-and-white. It was not until the mid-1960s that color sets started selling in large numbers, due in part to the color transition of 1965 in which it was announced that over half of all network prime-time programming would be broadcast in color that fall. The first all-color prime-time season came just one year later. In 1972, the last holdout among daytime network programs converted to color, resulting in the first completely all-color network season.
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Early color sets were either floor-standing console models or tabletop versions nearly as bulky and heavy, so in practice they remained firmly anchored in one place. GE's relatively compact and lightweight Porta-Color set was introduced in the spring of 1966. It used a transistor-based UHF tuner.[130] The first fully transistorized color television in the United States was the Quasar television introduced in 1967.[131] These developments made watching color television a more flexible and convenient proposition.
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The MOSFET (metal-oxide-semiconductor field-effect transistor, or MOS transistor) was invented by Mohamed M. Atalla and Dawon Kahng at Bell Labs in 1959,[132] and presented in 1960.[133] By the mid-1960s, RCA were using MOSFETs in their consumer television products.[134] RCA Laboratories researchers W.M. Austin, J.A. Dean, D.M. Griswold and O.P. Hart in 1966 described the use of the MOSFET in television circuits, including RF amplifier, low-level video, chroma and AGC circuits.[135] The power MOSFET was later widely adopted for television receiver circuits.[136]
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In 1972, sales of color sets finally surpassed sales of black-and-white sets. Color broadcasting in Europe was not standardized on the PAL format until the 1960s, and broadcasts did not start until 1967. By this point many of the technical problems in the early sets had been worked out, and the spread of color sets in Europe was fairly rapid. By the mid-1970s, the only stations broadcasting in black-and-white were a few high-numbered UHF stations in small markets, and a handful of low-power repeater stations in even smaller markets such as vacation spots. By 1979, even the last of these had converted to color and, by the early 1980s, B&W sets had been pushed into niche markets, notably low-power uses, small portable sets, or for use as video monitor screens in lower-cost consumer equipment. By the late 1980s even these areas switched to color sets.
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Digital television (DTV) is the transmission of audio and video by digitally processed and multiplexed signals, in contrast to the totally analog and channel separated signals used by analog television. Due to data compression, digital TV can support more than one program in the same channel bandwidth.[137] It is an innovative service that represents the most significant evolution in television broadcast technology since color television emerged in the 1950s.[138] Digital TV's roots have been tied very closely to the availability of inexpensive, high performance computers. It was not until the 1990s that digital TV became feasible.[139] Digital television was previously not practically feasible due to the impractically high bandwidth requirements of uncompressed digital video,[140][141] requiring around 200 Mbit/s bit-rate for a standard-definition television (SDTV) signal,[140] and over 1 Gbit/s for high-definition television (HDTV).[141]
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Digital TV became practically feasible in the early 1990s due to a major technological development, discrete cosine transform (DCT) video compression.[140][141] DCT coding is a lossy compression technique that was first proposed for image compression by Nasir Ahmed in 1972,[142] and was later adapted into a motion-compensated DCT video coding algorithm, for video coding standards such as the H.26x formats from 1988 onwards and the MPEG formats from 1991 onwards.[143][144] Motion-compensated DCT video compression significantly reduced the amount of bandwidth required for a digital TV signal.[140][141] DCT coding compressed down the bandwidth requirements of digital television signals to about 34 Mpps bit-rate for SDTV and around 70–140 Mbit/s for HDTV while maintaining near-studio-quality transmission, making digital television a practical reality in the 1990s.[141]
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A digital TV service was proposed in 1986 by Nippon Telegraph and Telephone (NTT) and the Ministry of Posts and Telecommunication (MPT) in Japan, where there were plans to develop an "Integrated Network System" service. However, it was not possible to practically implement such a digital TV service until the adoption of DCT video compression technology made it possible in the early 1990s.[140]
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In the mid-1980s, as Japanese consumer electronics firms forged ahead with the development of HDTV technology, the MUSE analog format proposed by NHK, a Japanese company, was seen as a pacesetter that threatened to eclipse U.S. electronics companies' technologies. Until June 1990, the Japanese MUSE standard, based on an analog system, was the front-runner among the more than 23 different technical concepts under consideration. Then, an American company, General Instrument, demonstrated the feasibility of a digital television signal. This breakthrough was of such significance that the FCC was persuaded to delay its decision on an ATV standard until a digitally based standard could be developed.
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In March 1990, when it became clear that a digital standard was feasible, the FCC made a number of critical decisions. First, the Commission declared that the new ATV standard must be more than an enhanced analog signal, but be able to provide a genuine HDTV signal with at least twice the resolution of existing television images.(7) Then, to ensure that viewers who did not wish to buy a new digital television set could continue to receive conventional television broadcasts, it dictated that the new ATV standard must be capable of being "simulcast" on different channels.(8)The new ATV standard also allowed the new DTV signal to be based on entirely new design principles. Although incompatible with the existing NTSC standard, the new DTV standard would be able to incorporate many improvements.
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The final standards adopted by the FCC did not require a single standard for scanning formats, aspect ratios, or lines of resolution. This compromise resulted from a dispute between the consumer electronics industry (joined by some broadcasters) and the computer industry (joined by the film industry and some public interest groups) over which of the two scanning processes—interlaced or progressive—would be best suited for the newer digital HDTV compatible display devices.[145] Interlaced scanning, which had been specifically designed for older analogue CRT display technologies, scans even-numbered lines first, then odd-numbered ones. In fact, interlaced scanning can be looked at as the first video compression model as it was partly designed in the 1940s to double the image resolution to exceed the limitations of the television broadcast bandwidth. Another reason for its adoption was to limit the flickering on early CRT screens whose phosphor coated screens could only retain the image from the electron scanning gun for a relatively short duration.[146] However interlaced scanning does not work as efficiently on newer display devices such as Liquid-crystal (LCD), for example, which are better suited to a more frequent progressive refresh rate.[145]
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Progressive scanning, the format that the computer industry had long adopted for computer display monitors, scans every line in sequence, from top to bottom. Progressive scanning in effect doubles the amount of data generated for every full screen displayed in comparison to interlaced scanning by painting the screen in one pass in 1/60-second, instead of two passes in 1/30-second. The computer industry argued that progressive scanning is superior because it does not "flicker" on the new standard of display devices in the manner of interlaced scanning. It also argued that progressive scanning enables easier connections with the Internet, and is more cheaply converted to interlaced formats than vice versa. The film industry also supported progressive scanning because it offered a more efficient means of converting filmed programming into digital formats. For their part, the consumer electronics industry and broadcasters argued that interlaced scanning was the only technology that could transmit the highest quality pictures then (and currently) feasible, i.e., 1,080 lines per picture and 1,920 pixels per line. Broadcasters also favored interlaced scanning because their vast archive of interlaced programming is not readily compatible with a progressive format. William F. Schreiber, who was director of the Advanced Television Research Program at the Massachusetts Institute of Technology from 1983 until his retirement in 1990, thought that the continued advocacy of interlaced equipment originated from consumer electronics companies that were trying to get back the substantial investments they made in the interlaced technology.[147]
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Digital television transition started in late 2000s. All governments across the world set the deadline for analog shutdown by 2010s. Initially the adoption rate was low, as the first digital tuner-equipped TVs were costly. But soon, as the price of digital-capable TVs dropped, more and more households were converting to digital televisions. The transition is expected to be completed worldwide by mid to late 2010s.
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The advent of digital television allowed innovations like smart TVs. A smart television, sometimes referred to as connected TV or hybrid TV, is a television set or set-top box with integrated Internet and Web 2.0 features, and is an example of technological convergence between computers, television sets and set-top boxes. Besides the traditional functions of television sets and set-top boxes provided through traditional Broadcasting media, these devices can also provide Internet TV, online interactive media, over-the-top content, as well as on-demand streaming media, and home networking access. These TVs come pre-loaded with an operating system.[9][148][149][150]
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Smart TV should not to be confused with Internet TV, Internet Protocol television (IPTV) or with Web TV. Internet television refers to the receiving of television content over the Internet instead of by traditional systems—terrestrial, cable and satellite (although internet itself is received by these methods). IPTV is one of the emerging Internet television technology standards for use by television broadcasters. Web television (WebTV) is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV. A first patent was filed in 1994[151] (and extended the following year)[152] for an "intelligent" television system, linked with data processing systems, by means of a digital or analog network. Apart from being linked to data networks, one key point is its ability to automatically download necessary software routines, according to a user's demand, and process their needs. Major TV manufacturers have announced production of smart TVs only, for middle-end and high-end TVs in 2015.[6][7][8] Smart TVs have gotten more affordable compared to when they were first introduced, with 46 million of U.S. households having at least one as of 2019.[153]
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3D television conveys depth perception to the viewer by employing techniques such as stereoscopic display, multi-view display, 2D-plus-depth, or any other form of 3D display. Most modern 3D television sets use an active shutter 3D system or a polarized 3D system, and some are autostereoscopic without the need of glasses. Stereoscopic 3D television was demonstrated for the first time on 10 August 1928, by John Logie Baird in his company's premises at 133 Long Acre, London.[154] Baird pioneered a variety of 3D television systems using electromechanical and cathode-ray tube techniques. The first 3D TV was produced in 1935. The advent of digital television in the 2000s greatly improved 3D TVs. Although 3D TV sets are quite popular for watching 3D home media such as on Blu-ray discs, 3D programming has largely failed to make inroads with the public. Many 3D television channels which started in the early 2010s were shut down by the mid-2010s. According to DisplaySearch 3D televisions shipments totaled 41.45 million units in 2012, compared with 24.14 in 2011 and 2.26 in 2010.[155] As of late 2013, the number of 3D TV viewers started to decline.[156][157][158][159][160]
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Programming is broadcast by television stations, sometimes called "channels", as stations are licensed by their governments to broadcast only over assigned channels in the television band. At first, terrestrial broadcasting was the only way television could be widely distributed, and because bandwidth was limited, i.e., there were only a small number of channels available, government regulation was the norm. In the U.S., the Federal Communications Commission (FCC) allowed stations to broadcast advertisements beginning in July 1941, but required public service programming commitments as a requirement for a license. By contrast, the United Kingdom chose a different route, imposing a television license fee on owners of television reception equipment to fund the British Broadcasting Corporation (BBC), which had public service as part of its Royal Charter.
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WRGB claims to be the world's oldest television station, tracing its roots to an experimental station founded on 13 January 1928, broadcasting from the General Electric factory in Schenectady, NY, under the call letters W2XB.[161] It was popularly known as "WGY Television" after its sister radio station. Later in 1928, General Electric started a second facility, this one in New York City, which had the call letters W2XBS and which today is known as WNBC. The two stations were experimental in nature and had no regular programming, as receivers were operated by engineers within the company. The image of a Felix the Cat doll rotating on a turntable was broadcast for 2 hours every day for several years as new technology was being tested by the engineers. On 2 November 1936, the BBC began transmitting the world's first public regular high-definition service from the Victorian Alexandra Palace in north London.[162] It therefore claims to be the birthplace of TV broadcasting as we know it today.
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With the widespread adoption of cable across the United States in the 1970s and 80s, terrestrial television broadcasts have been in decline; in 2013 it was estimated that about 7% of US households used an antenna.[163][164] A slight increase in use began around 2010 due to switchover to digital terrestrial television broadcasts, which offered pristine image quality over very large areas, and offered an alternate to cable television (CATV) for cord cutters. All other countries around the world are also in the process of either shutting down analog terrestrial television or switching over to digital terrestrial television.
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Cable television is a system of broadcasting television programming to paying subscribers via radio frequency (RF) signals transmitted through coaxial cables or light pulses through fiber-optic cables. This contrasts with traditional terrestrial television, in which the television signal is transmitted over the air by radio waves and received by a television antenna attached to the television. In the 2000s, FM radio programming, high-speed Internet, telephone service, and similar non-television services may also be provided through these cables. The abbreviation CATV is often used for cable television. It originally stood for Community Access Television or Community Antenna Television, from cable television's origins in 1948: in areas where over-the-air reception was limited by distance from transmitters or mountainous terrain, large "community antennas" were constructed, and cable was run from them to individual homes.[165] The origins of cable broadcasting are even older as radio programming was distributed by cable in some European cities as far back as 1924. Earlier cable television was analog, but since the 2000s, all cable operators have switched to, or are in the process of switching to, digital cable television.
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Satellite television is a system of supplying television programming using broadcast signals relayed from communication satellites. The signals are received via an outdoor parabolic reflector antenna usually referred to as a satellite dish and a low-noise block downconverter (LNB). A satellite receiver then decodes the desired television program for viewing on a television set. Receivers can be external set-top boxes, or a built-in television tuner. Satellite television provides a wide range of channels and services, especially to geographic areas without terrestrial television or cable television.
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The most common method of reception is direct-broadcast satellite television (DBSTV), also known as "direct to home" (DTH).[166] In DBSTV systems, signals are relayed from a direct broadcast satellite on the Ku wavelength and are completely digital.[167] Satellite TV systems formerly used systems known as television receive-only. These systems received analog signals transmitted in the C-band spectrum from FSS type satellites, and required the use of large dishes. Consequently, these systems were nicknamed "big dish" systems, and were more expensive and less popular.[168]
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The direct-broadcast satellite television signals were earlier analog signals and later digital signals, both of which require a compatible receiver. Digital signals may include high-definition television (HDTV). Some transmissions and channels are free-to-air or free-to-view, while many other channels are pay television requiring a subscription.[169]
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In 1945, British science fiction writer Arthur C. Clarke proposed a worldwide communications system which would function by means of three satellites equally spaced apart in earth orbit.[170][171] This was published in the October 1945 issue of the Wireless World magazine and won him the Franklin Institute's Stuart Ballantine Medal in 1963.[172][173]
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The first satellite television signals from Europe to North America were relayed via the Telstar satellite over the Atlantic ocean on 23 July 1962.[174] The signals were received and broadcast in North American and European countries and watched by over 100 million.[174] Launched in 1962, the Relay 1 satellite was the first satellite to transmit television signals from the US to Japan.[175] The first geosynchronous communication satellite, Syncom 2, was launched on 26 July 1963.[176]
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The world's first commercial communications satellite, called Intelsat I and nicknamed "Early Bird", was launched into geosynchronous orbit on 6 April 1965.[177] The first national network of television satellites, called Orbita, was created by the Soviet Union in October 1967, and was based on the principle of using the highly elliptical Molniya satellite for rebroadcasting and delivering of television signals to ground downlink stations.[178] The first commercial North American satellite to carry television transmissions was Canada's geostationary Anik 1, which was launched on 9 November 1972.[179] ATS-6, the world's first experimental educational and Direct Broadcast Satellite (DBS), was launched on 30 May 1974.[180] It transmitted at 860 MHz using wideband FM modulation and had two sound channels. The transmissions were focused on the Indian subcontinent but experimenters were able to receive the signal in Western Europe using home constructed equipment that drew on UHF television design techniques already in use.[181]
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The first in a series of Soviet geostationary satellites to carry Direct-To-Home television, Ekran 1, was launched on 26 October 1976.[182] It used a 714 MHz UHF downlink frequency so that the transmissions could be received with existing UHF television technology rather than microwave technology.[183]
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Internet television (Internet TV) (or online television) is the digital distribution of television content via the Internet as opposed to traditional systems like terrestrial, cable, and satellite, although the Internet itself is received by terrestrial, cable, or satellite methods. Internet television is a general term that covers the delivery of television shows, and other video content, over the Internet by video streaming technology, typically by major traditional television broadcasters. Internet television should not be confused with Smart TV, IPTV or with Web TV. Smart television refers to the TV set which has a built-in operating system. Internet Protocol television (IPTV) is one of the emerging Internet television technology standards for use by television broadcasters. Web television is a term used for programs created by a wide variety of companies and individuals for broadcast on Internet TV.
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A television set, also called a television receiver, television, TV set, TV, or "telly", is a device that combines a tuner, display, an amplifier, and speakers for the purpose of viewing television and hearing its audio components. Introduced in the late 1920s in mechanical form, television sets became a popular consumer product after World War II in electronic form, using cathode ray tubes. The addition of color to broadcast television after 1953 further increased the popularity of television sets and an outdoor antenna became a common feature of suburban homes. The ubiquitous television set became the display device for recorded media in the 1970s, such as Betamax and VHS, which enabled viewers to record TV shows and watch prerecorded movies. In the subsequent decades, TVs were used to watch DVDs and Blu-ray Discs of movies and other content. Major TV manufacturers announced the discontinuation of CRT, DLP, plasma and fluorescent-backlit LCDs by the mid-2010s. Televisions since 2010s mostly use LEDs.[3][4][184][185] LEDs are expected to be gradually replaced by OLEDs in the near future.[5]
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The earliest systems employed a spinning disk to create and reproduce images.[186] These usually had a low resolution and screen size and never became popular with the public.
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The cathode ray tube (CRT) is a vacuum tube containing one or more electron guns (a source of electrons or electron emitter) and a fluorescent screen used to view images.[33] It has a means to accelerate and deflect the electron beam(s) onto the screen to create the images. The images may represent electrical waveforms (oscilloscope), pictures (television, computer monitor), radar targets or others. The CRT uses an evacuated glass envelope which is large, deep (i.e. long from front screen face to rear end), fairly heavy, and relatively fragile. As a matter of safety, the face is typically made of thick lead glass so as to be highly shatter-resistant and to block most X-ray emissions, particularly if the CRT is used in a consumer product.
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In television sets and computer monitors, the entire front area of the tube is scanned repetitively and systematically in a fixed pattern called a raster. An image is produced by controlling the intensity of each of the three electron beams, one for each additive primary color (red, green, and blue) with a video signal as a reference.[187] In all modern CRT monitors and televisions, the beams are bent by magnetic deflection, a varying magnetic field generated by coils and driven by electronic circuits around the neck of the tube, although electrostatic deflection is commonly used in oscilloscopes, a type of diagnostic instrument.[187]
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Digital Light Processing (DLP) is a type of video projector technology that uses a digital micromirror device. Some DLPs have a TV tuner, which makes them a type of TV display. It was originally developed in 1987 by Dr. Larry Hornbeck of Texas Instruments. While the DLP imaging device was invented by Texas Instruments, the first DLP based projector was introduced by Digital Projection Ltd in 1997. Digital Projection and Texas Instruments were both awarded Emmy Awards in 1998 for invention of the DLP projector technology. DLP is used in a variety of display applications from traditional static displays to interactive displays and also non-traditional embedded applications including medical, security, and industrial uses. DLP technology is used in DLP front projectors (standalone projection units for classrooms and business primarily), but also in private homes; in these cases, the image is projected onto a projection screen. DLP is also used in DLP rear projection television sets and digital signs. It is also used in about 85% of digital cinema projection.[188]
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A plasma display panel (PDP) is a type of flat panel display common to large TV displays 30 inches (76 cm) or larger. They are called "plasma" displays because the technology utilizes small cells containing electrically charged ionized gases, or what are in essence chambers more commonly known as fluorescent lamps.
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Liquid-crystal-display televisions (LCD TV) are television sets that use LCD display technology to produce images. LCD televisions are much thinner and lighter than cathode ray tube (CRTs) of similar display size, and are available in much larger sizes (e.g., 90-inch diagonal). When manufacturing costs fell, this combination of features made LCDs practical for television receivers. LCDs come in two types: those using cold cathode fluorescent lamps, simply called LCDs and those using LED as backlight called as LEDs.
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In 2007, LCD televisions surpassed sales of CRT-based televisions worldwide for the first time, and their sales figures relative to other technologies accelerated. LCD TVs have quickly displaced the only major competitors in the large-screen market, the Plasma display panel and rear-projection television.[189] In mid 2010s LCDs especially LEDs became, by far, the most widely produced and sold television display type.[184][185] LCDs also have disadvantages. Other technologies address these weaknesses, including OLEDs, FED and SED, but as of 2014[update] none of these have entered widespread production.
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An OLED (organic light-emitting diode) is a light-emitting diode (LED) in which the emissive electroluminescent layer is a film of organic compound which emits light in response to an electric current. This layer of organic semiconductor is situated between two electrodes. Generally, at least one of these electrodes is transparent. OLEDs are used to create digital displays in devices such as television screens. It is also used for computer monitors, portable systems such as mobile phones, handheld game consoles and PDAs.
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There are two main families of OLED: those based on small molecules and those employing polymers. Adding mobile ions to an OLED creates a light-emitting electrochemical cell or LEC, which has a slightly different mode of operation. OLED displays can use either passive-matrix (PMOLED) or active-matrix (AMOLED) addressing schemes. Active-matrix OLEDs require a thin-film transistor backplane to switch each individual pixel on or off, but allow for higher resolution and larger display sizes.
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An OLED display works without a backlight. Thus, it can display deep black levels and can be thinner and lighter than a liquid crystal display (LCD). In low ambient light conditions such as a dark room an OLED screen can achieve a higher contrast ratio than an LCD, whether the LCD uses cold cathode fluorescent lamps or LED backlight. OLEDs are expected to replace other forms of display in near future.[5]
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Low-definition television or LDTV refers to television systems that have a lower screen resolution than standard-definition television systems such 240p (320*240). It is used in handheld television. The most common source of LDTV programming is the Internet, where mass distribution of higher-resolution video files could overwhelm computer servers and take too long to download. Many mobile phones and portable devices such as Apple's iPod Nano, or Sony's PlayStation Portable use LDTV video, as higher-resolution files would be excessive to the needs of their small screens (320×240 and 480×272 pixels respectively). The current generation of iPod Nanos have LDTV screens, as do the first three generations of iPod Touch and iPhone (480×320). For the first years of its existence, YouTube offered only one, low-definition resolution of 320x240p at 30fps or less. A standard, consumer grade VHS videotape can be considered SDTV due to its resolution (approximately 360 × 480i/576i).
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Standard-definition television or SDTV refers to two different resolutions: 576i, with 576 interlaced lines of resolution, derived from the European-developed PAL and SECAM systems; and 480i based on the American National Television System Committee NTSC system. SDTV is a television system that uses a resolution that is not considered to be either high-definition television (720p, 1080i, 1080p, 1440p, 4K UHDTV, and 8K UHD) or enhanced-definition television (EDTV 480p). In North America, digital SDTV is broadcast in the same 4:3 aspect ratio as NTSC signals with widescreen content being center cut.[190] However, in other parts of the world that used the PAL or SECAM color systems, standard-definition television is now usually shown with a 16:9 aspect ratio, with the transition occurring between the mid-1990s and mid-2000s. Older programs with a 4:3 aspect ratio are shown in the US as 4:3 with non-ATSC countries preferring to reduce the horizontal resolution by anamorphically scaling a pillarboxed image.
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High-definition television (HDTV) provides a resolution that is substantially higher than that of standard-definition television.
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HDTV may be transmitted in various formats:
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Ultra-high-definition television (also known as Super Hi-Vision, Ultra HD television, UltraHD, UHDTV, or UHD) includes 4K UHD (2160p) and 8K UHD (4320p), which are two digital video formats proposed by NHK Science & Technology Research Laboratories and defined and approved by the International Telecommunication Union (ITU). The Consumer Electronics Association announced on 17 October 2012, that "Ultra High Definition", or "Ultra HD", would be used for displays that have an aspect ratio of at least 16:9 and at least one digital input capable of carrying and presenting native video at a minimum resolution of 3840×2160 pixels.[191][192]
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North American consumers purchase a new television set on average every seven years, and the average household owns 2.8 televisions. As of 2011[update], 48 million are sold each year at an average price of $460 and size of 38 in (97 cm).[193]
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Getting TV programming shown to the public can happen in many different ways. After production, the next step is to market and deliver the product to whichever markets are open to using it. This typically happens on two levels:
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First-run programming is increasing on subscription services outside the US, but few domestically produced programs are syndicated on domestic free-to-air (FTA) elsewhere. This practice is increasing, however, generally on digital-only FTA channels or with subscriber-only, first-run material appearing on FTA. Unlike the US, repeat FTA screenings of an FTA network program usually only occur on that network. Also, affiliates rarely buy or produce non-network programming that is not centered on local programming.
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Television genres include a broad range of programming types that entertain, inform, and educate viewers. The most expensive entertainment genres to produce are usually dramas and dramatic miniseries. However, other genres, such as historical Western genres, may also have high production costs.
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Popular culture entertainment genres include action-oriented shows such as police, crime, detective dramas, horror, or thriller shows. As well, there are also other variants of the drama genre, such as medical dramas and daytime soap operas. Science fiction shows can fall into either the drama or action category, depending on whether they emphasize philosophical questions or high adventure. Comedy is a popular genre which includes situation comedy (sitcom) and animated shows for the adult demographic such as South Park.
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The least expensive forms of entertainment programming genres are game shows, talk shows, variety shows, and reality television. Game shows feature contestants answering questions and solving puzzles to win prizes. Talk shows contain interviews with film, television, music and sports celebrities and public figures. Variety shows feature a range of musical performers and other entertainers, such as comedians and magicians, introduced by a host or Master of Ceremonies. There is some crossover between some talk shows and variety shows because leading talk shows often feature performances by bands, singers, comedians, and other performers in between the interview segments. Reality TV shows "regular" people (i.e., not actors) facing unusual challenges or experiences ranging from arrest by police officers (COPS) to significant weight loss (The Biggest Loser). A variant version of reality shows depicts celebrities doing mundane activities such as going about their everyday life (The Osbournes, Snoop Dogg's Father Hood) or doing regular jobs (The Simple Life).
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Fictional television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include series such as Twin Peaks and The Sopranos. Kristin Thompson argues that some of these television series exhibit traits also found in art films, such as psychological realism, narrative complexity, and ambiguous plotlines. Nonfiction television programs that some television scholars and broadcasting advocacy groups argue are "quality television", include a range of serious, noncommercial, programming aimed at a niche audience, such as documentaries and public affairs shows.
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Around the globe, broadcast TV is financed by government, advertising, licensing (a form of tax), subscription, or any combination of these. To protect revenues, subscription TV channels are usually encrypted to ensure that only subscribers receive the decryption codes to see the signal. Unencrypted channels are known as free to air or FTA. In 2009, the global TV market represented 1,217.2 million TV households with at least one TV and total revenues of 268.9 billion EUR (declining 1.2% compared to 2008).[195] North America had the biggest TV revenue market share with 39% followed by Europe (31%), Asia-Pacific (21%), Latin America (8%), and Africa and the Middle East (2%).[196] Globally, the different TV revenue sources divide into 45–50% TV advertising revenues, 40–45% subscription fees and 10% public funding.[197][198]
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TV's broad reach makes it a powerful and attractive medium for advertisers. Many TV networks and stations sell blocks of broadcast time to advertisers ("sponsors") to fund their programming.[199] Television advertisements (variously called a television commercial, commercial or ad in American English, and known in British English as an advert) is a span of television programming produced and paid for by an organization, which conveys a message, typically to market a product or service. Advertising revenue provides a significant portion of the funding for most privately owned television networks. The vast majority of television advertisements today consist of brief advertising spots, ranging in length from a few seconds to several minutes (as well as program-length infomercials). Advertisements of this sort have been used to promote a wide variety of goods, services and ideas since the beginning of television.
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The effects of television advertising upon the viewing public (and the effects of mass media in general) have been the subject of philosophical discourse by such luminaries as Marshall McLuhan. The viewership of television programming, as measured by companies such as Nielsen Media Research, is often used as a metric for television advertisement placement, and consequently, for the rates charged to advertisers to air within a given network, television program, or time of day (called a "daypart"). In many countries, including the United States, television campaign advertisements are considered indispensable for a political campaign. In other countries, such as France, political advertising on television is heavily restricted,[200] while some countries, such as Norway, completely ban political advertisements.
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The first official, paid television advertisement was broadcast in the United States on 1 July 1941 over New York station WNBT (now WNBC) before a baseball game between the Brooklyn Dodgers and Philadelphia Phillies. The announcement for Bulova watches, for which the company paid anywhere from $4.00 to $9.00 (reports vary), displayed a WNBT test pattern modified to look like a clock with the hands showing the time. The Bulova logo, with the phrase "Bulova Watch Time", was shown in the lower right-hand quadrant of the test pattern while the second hand swept around the dial for one minute.[201][202] The first TV ad broadcast in the UK was on ITV on 22 September 1955, advertising Gibbs SR toothpaste. The first TV ad broadcast in Asia was on Nippon Television in Tokyo on 28 August 1953, advertising Seikosha (now Seiko), which also displayed a clock with the current time.[203]
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Since inception in the US in 1941,[204] television commercials have become one of the most effective, persuasive, and popular methods of selling products of many sorts, especially consumer goods. During the 1940s and into the 1950s, programs were hosted by single advertisers. This, in turn, gave great creative license to the advertisers over the content of the show. Perhaps due to the quiz show scandals in the 1950s,[205] networks shifted to the magazine concept, introducing advertising breaks with multiple advertisers.
|
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US advertising rates are determined primarily by Nielsen ratings. The time of the day and popularity of the channel determine how much a TV commercial can cost. For example, it can cost approximately $750,000 for a 30-second block of commercial time during the highly popular American Idol, while the same amount of time for the Super Bowl can cost several million dollars. Conversely, lesser-viewed time slots, such as early mornings and weekday afternoons, are often sold in bulk to producers of infomercials at far lower rates. In recent years, the paid program or infomercial has become common, usually in lengths of 30 minutes or one hour. Some drug companies and other businesses have even created "news" items for broadcast, known in the industry as video news releases, paying program directors to use them.[206]
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Some TV programs also deliberately place products into their shows as advertisements, a practice started in feature films[207] and known as product placement. For example, a character could be drinking a certain kind of soda, going to a particular chain restaurant, or driving a certain make of car. (This is sometimes very subtle, with shows having vehicles provided by manufacturers for low cost in exchange as a product placement). Sometimes, a specific brand or trade mark, or music from a certain artist or group, is used. (This excludes guest appearances by artists who perform on the show.)
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The TV regulator oversees TV advertising in the United Kingdom. Its restrictions have applied since the early days of commercially funded TV. Despite this, an early TV mogul, Roy Thomson, likened the broadcasting licence as being a "licence to print money".[208] Restrictions mean that the big three national commercial TV channels: ITV, Channel 4, and Channel 5 can show an average of only seven minutes of advertising per hour (eight minutes in the peak period). Other broadcasters must average no more than nine minutes (twelve in the peak). This means that many imported TV shows from the US have unnatural pauses where the UK company does not utilize the narrative breaks intended for more frequent US advertising. Advertisements must not be inserted in the course of certain specific proscribed types of programs which last less than half an hour in scheduled duration; this list includes any news or current affairs programs, documentaries, and programs for children; additionally, advertisements may not be carried in a program designed and broadcast for reception in schools or in any religious broadcasting service or other devotional program or during a formal Royal ceremony or occasion. There also must be clear demarcations in time between the programs and the advertisements. The BBC, being strictly non-commercial, is not allowed to show advertisements on television in the UK, although it has many advertising-funded channels abroad. The majority of its budget comes from television license fees (see below) and broadcast syndication, the sale of content to other broadcasters.
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Broadcast advertising is regulated by the Broadcasting Authority of Ireland.[209]
|
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Some TV channels are partly funded from subscriptions; therefore, the signals are encrypted during broadcast to ensure that only the paying subscribers have access to the decryption codes to watch pay television or specialty channels. Most subscription services are also funded by advertising.
|
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Television services in some countries may be funded by a television licence or a form of taxation, which means that advertising plays a lesser role or no role at all. For example, some channels may carry no advertising at all and some very little, including:
|
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|
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+
The BBC carries no television advertising on its UK channels and is funded by an annual television licence paid by premises receiving live TV broadcasts. Currently, it is estimated that approximately 26.8 million UK private domestic households own televisions, with approximately 25 million TV licences in all premises in force as of 2010.[210] This television license fee is set by the government, but the BBC is not answerable to or controlled by the government.
|
200 |
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The two main BBC TV channels are watched by almost 90% of the population each week and overall have 27% share of total viewing,[211] despite the fact that 85% of homes are multichannel, with 42% of these having access to 200 free to air channels via satellite and another 43% having access to 30 or more channels via Freeview.[212] The licence that funds the seven advertising-free BBC TV channels costs £147 a year (about US$200) as of 2018 regardless of the number of TV sets owned; the price is reduced by two-thirds if only black and white television is received.[213] When the same sporting event has been presented on both BBC and commercial channels, the BBC always attracts the lion's share of the audience, indicating that viewers prefer to watch TV uninterrupted by advertising.
|
202 |
+
|
203 |
+
Other than internal promotional material, the Australian Broadcasting Corporation (ABC) carries no advertising; it is banned under the ABC Act 1983. The ABC receives its funding from the Australian government every three years. In the 2014/15 federal budget, the ABC received A$1.11 billion.[214] The funds provide for the ABC's television, radio, online, and international outputs. The ABC also receives funds from its many ABC shops across Australia. Although funded by the Australian government, the editorial independence of the ABC is ensured through law.
|
204 |
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In France, government-funded channels carry advertisements, yet those who own television sets have to pay an annual tax ("la redevance audiovisuelle").[215]
|
206 |
+
|
207 |
+
In Japan, NHK is paid for by license fees (known in Japanese as reception fee (受信料, Jushinryō)). The broadcast law that governs NHK's funding stipulates that any television equipped to receive NHK is required to pay. The fee is standardized, with discounts for office workers and students who commute, as well a general discount for residents of Okinawa prefecture.
|
208 |
+
|
209 |
+
Broadcast programming, or TV listings in the United Kingdom, is the practice of organizing television programs in a schedule, with broadcast automation used to regularly change the scheduling of TV programs to build an audience for a new show, retain that audience, or compete with other broadcasters' programs.
|
210 |
+
|
211 |
+
Television has played a pivotal role in the socialization of the 20th and 21st centuries. There are many aspects of television that can be addressed, including negative issues such as media violence. Current research is discovering that individuals suffering from social isolation can employ television to create what is termed a parasocial or faux relationship with characters from their favorite television shows and movies as a way of deflecting feelings of loneliness and social deprivation.[216] Several studies have found that educational television has many advantages. The article "The Good Things about Television"[217] argues that television can be a very powerful and effective learning tool for children if used wisely.
|
212 |
+
|
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+
Methodist denominations in the conservative holiness movement, such as the Allegheny Wesleyan Methodist Connection and the Evangelical Wesleyan Church, eschew the use of the television.[218]
|
214 |
+
|
215 |
+
Children, especially those aged 5 or younger, are at risk of injury from falling televisions.[219] A CRT-style television that falls on a child will, because of its weight, hit with the equivalent force of falling multiple stories from a building.[220] Newer flat-screen televisions are "top-heavy and have narrow bases", which means that a small child can easily pull one over.[221] As of 2015[update], TV tip-overs were responsible for more than 10,000 injuries per year to children, at a cost of more than $8 million per year in emergency care.[219][221]
|
216 |
+
|
217 |
+
A 2017 study in The Journal of Human Resources found that exposure to cable television reduced cognitive ability and high school graduation rates for boys. This effect was stronger for boys from more educated families. The article suggests a mechanism where light television entertainment crowds out more cognitively stimulating activities.[222]
|
218 |
+
|
219 |
+
With high lead content in CRTs and the rapid diffusion of new flat-panel display technologies, some of which (LCDs) use lamps which contain mercury, there is growing concern about electronic waste from discarded televisions. Related occupational health concerns exist, as well, for disassemblers removing copper wiring and other materials from CRTs. Further environmental concerns related to television design and use relate to the devices' increasing electrical energy requirements.[223]
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1 |
+
|
2 |
+
|
3 |
+
Temperature is a physical property of matter that quantitatively expresses hot and cold. It is the manifestation of thermal energy, present in all matter, which is the source of the occurrence of heat, a flow of energy, when a body is in contact with another that is colder.
|
4 |
+
|
5 |
+
Temperature is measured with a thermometer. Thermometers are calibrated in various temperature scales that historically have used various reference points and thermometric substances for definition. The most common scales are the Celsius scale (formerly called centigrade, denoted °C), the Fahrenheit scale (denoted °F), and the Kelvin scale (denoted K), the last of which is predominantly used for scientific purposes by conventions of the International System of Units (SI).
|
6 |
+
|
7 |
+
The lowest theoretical temperature is absolute zero, at which no more thermal energy can be extracted from a body. Experimentally, it can only be approached very closely, but not reached, which is recognized in the third law of thermodynamics.
|
8 |
+
|
9 |
+
Temperature is important in all fields of natural science, including physics, chemistry, Earth science, medicine, and biology, as well as most aspects of daily life.
|
10 |
+
|
11 |
+
Many physical processes are related to temperature, such as:
|
12 |
+
|
13 |
+
Temperature scales differ in two ways: the point chosen as zero degrees, and the magnitudes of incremental units or degrees on the scale.
|
14 |
+
|
15 |
+
The Celsius scale (°C) is used for common temperature measurements in most of the world. It is an empirical scale that was developed by a historical progress, which led to its zero point 0 °C being defined by the freezing point of water, and additional degrees defined so that 100 °C was the boiling point of water, both at sea-level atmospheric pressure. Because of the 100-degree interval, it was called a centigrade scale.[3] Since the standardization of the kelvin in the International System of Units, it has subsequently been redefined in terms of the equivalent fixing points on the Kelvin scale, and so that a temperature increment of one degree Celsius is the same as an increment of one kelvin, though they differ by an additive offset of approximately 273.15.
|
16 |
+
|
17 |
+
The United States commonly uses the Fahrenheit scale, on which water freezes at 32 °F and boils at 212 °F at sea-level atmospheric pressure.
|
18 |
+
|
19 |
+
At the absolute zero of temperature, no more energy can be removed from matter as heat, a fact expressed in the third law of thermodynamics. At this temperature, matter contains no macroscopic thermal energy, but still has quantum-mechanical zero-point energy as predicted by the uncertainty principle. This does not enter into the definition of absolute temperature. Experimentally, absolute zero can only be approached very closely, but can never be actually reached. If it were possible to cool a system to absolute zero, all classical motion of its particles would cease and they would be at complete rest in this classical sense. The absolute zero, defined as 0 K, is approximately equal to −273.15 °C, or −459.67 °F.
|
20 |
+
|
21 |
+
Referring to the Boltzmann constant, to the Maxwell–Boltzmann distribution, and to the Boltzmann statistical mechanical definition of entropy, as distinct from the Gibbs definition,[4] for independently moving microscopic particles, disregarding interparticle potential energy, by international agreement, a temperature scale is defined and said to be absolute because it is independent of the characteristics of particular thermometric substances and thermometer mechanisms. Apart from the absolute zero, it does not have a reference temperature. It is known as the Kelvin scale, widely used in science and technology. The kelvin (the word is spelled with a lower-case k) is the unit of temperature in the International System of Units (SI). The temperature of a body in its own state of thermodynamic equilibrium is always positive, relative to the absolute zero.
|
22 |
+
|
23 |
+
Besides the internationally agreed Kelvin scale, there is also a thermodynamic temperature scale, invented by Kelvin, also with its numerical zero at the absolute zero of temperature, but directly relating to purely macroscopic thermodynamic concepts, including the macroscopic entropy, though microscopically referable to the Gibbs statistical mechanical definition of entropy for the canonical ensemble, that takes interparticle potential energy into account, as well as independent particle motion, so that it can account for measurements of temperatures near absolute zero.[4] This scale has a reference temperature at the triple point of water, the numerical value of which is defined by measurements using the aforementioned internationally agreed Kelvin scale.
|
24 |
+
|
25 |
+
Many scientific measurements use the Kelvin temperature scale (unit symbol: K), named in honor of the physicist who first defined it. It is an absolute scale. Its numerical zero point, 0 K, is at the absolute zero of temperature. Since May, 2019, its degrees have been defined through particle kinetic theory, and statistical mechanics. In the International System of Units (SI), the magnitude of the kelvin is defined through various empirical measurements of the average kinetic energies of microscopic particles. It is numerically evaluated in terms of the Boltzmann constant, the value of which is defined as fixed by international convention.[5][6]
|
26 |
+
|
27 |
+
Since May 2019, the magnitude of the kelvin is defined in relation to microscopic phenomena, characterized in terms of statistical mechanics. Previously, since 1954, the International System of Units defined a scale and unit for the kelvin as a thermodynamic temperature, by using the reliably reproducible temperature of the triple point of water as a second reference point, the first reference point being 0 K at absolute zero.
|
28 |
+
|
29 |
+
Historically, the triple point temperature of water was defined as exactly 273.16 units of the measurement increment. Today it is an empirically measured quantity. The freezing point of water at sea-level atmospheric pressure occurs at approximately 273.15 K = 0 °C.
|
30 |
+
|
31 |
+
There is a variety of kinds of temperature scale. It may be convenient to classify them as empirically and theoretically based. Empirical temperature scales are historically older, while theoretically based scales arose in the middle of the nineteenth century.[7][8]
|
32 |
+
|
33 |
+
Empirically based temperature scales rely directly on measurements of simple macroscopic physical properties of materials. For example, the length of a column of mercury, confined in a glass-walled capillary tube, is dependent largely on temperature, and is the basis of the very useful mercury-in-glass thermometer. Such scales are valid only within convenient ranges of temperature. For example, above the boiling point of mercury, a mercury-in-glass thermometer is impracticable. Most materials expand with temperature increase, but some materials, such as water, contract with temperature increase over some specific range, and then they are hardly useful as thermometric materials. A material is of no use as a thermometer near one of its phase-change temperatures, for example its boiling-point.
|
34 |
+
|
35 |
+
In spite of these limitations, most generally used practical thermometers are of the empirically based kind. Especially, it was used for calorimetry, which contributed greatly to the discovery of thermodynamics. Nevertheless, empirical thermometry has serious drawbacks when judged as a basis for theoretical physics. Empirically based thermometers, beyond their base as simple direct measurements of ordinary physical properties of thermometric materials, can be re-calibrated, by use of theoretical physical reasoning, and this can extend their range of adequacy.
|
36 |
+
|
37 |
+
Theoretically based temperature scales are based directly on theoretical arguments, especially those of kinetic theory and thermodynamics. They are more or less ideally realised in practically feasible physical devices and materials. Theoretically based temperature scales are used to provide calibrating standards for practical empirically based thermometers.
|
38 |
+
|
39 |
+
In physics, the internationally agreed conventional temperature scale is called the Kelvin scale. It is calibrated through the internationally agreed and prescribed value of the Boltzmann constant,[5][6] referring to motions of microscopic particles, such as atoms, molecules, and electrons, constituent in the body whose temperature is to be measured. In contrast with the thermodynamic temperature scale invented by Kelvin, the presently conventional Kelvin temperature is not defined through comparison with the temperature of a reference state of a standard body, nor in terms of macroscopic thermodynamics.
|
40 |
+
|
41 |
+
Apart from the absolute zero of temperature, Kelvin temperature of a body in a state of internal thermodynamic equilibrium is defined by measurements of suitably chosen of its physical properties, such as have precisely known theoretical explanations in terms of the Boltzmann constant. That constant refers to chosen kinds of motion of microscopic particles in the constitution of the body. In those kinds of motion, the particles move individually, without mutual interaction. Such motions are typically interrupted by inter-particle collisions, but for temperature measurement, the motions are chosen so that, between collisions, the non-interactive segments of their trajectories are known to be accessible to accurate measurement. For this purpose, interperticle potential energy is disregarded.
|
42 |
+
|
43 |
+
In an ideal gas, and in other theoretically understood bodies, the Kelvin temperature is defined to be proportional to the average kinetic energy of non-interactively moving microscopic particles, which can be measured by suitable techniques. The proportionality constant is a simple multiple of the Boltzmann constant. If molecules, atoms, or electrons,[9][10] are emitted from a material and their velocities are measured, the spectrum of their velocities often nearly obeys a theoretical law called the Maxwell–Boltzmann distribution, which gives a well-founded measurement of temperatures for which the law holds.[11] There have not yet been successful experiments of this same kind that directly use the Fermi–Dirac distribution for thermometry, but perhaps that will be achieved in future.[12]
|
44 |
+
|
45 |
+
The speed of sound in a gas can be calculated theoretically from the molecular character of the gas, from its temperature and pressure, and from the value of Boltzmann's constant. For a gas of known molecular character and pressure, this provides a relation between temperature and Boltzmann's constant. Those quantities can be known or measured more precisely than can the thermodynamic variables that define the state of a sample of water at its triple point. Consequently, taking the value of Boltzmann's constant as a primarily defined reference of exactly defined value, a measurement of the speed of sound can provide a more precise measurement of the temperature of the gas.[13]
|
46 |
+
|
47 |
+
Measurement of the spectrum of electromagnetic radiation from an ideal three-dimensional black body can provide an accurate temperature measurement because the frequency of maximum spectral radiance of black-body radiation is directly proportional to the temperature of the black body; this is known as Wien's displacement law and has a theoretical explanation in Planck's law and the Bose–Einstein law.
|
48 |
+
|
49 |
+
Measurement of the spectrum of noise-power produced by an electrical resistor can also provide an accurate temperature measurement. The resistor has two terminals and is in effect a one-dimensional body. The Bose-Einstein law for this case indicates that the noise-power is directly proportional to the temperature of the resistor and to the value of its resistance and to the noise band-width. In a given frequency band, the noise-power has equal contributions from every frequency and is called Johnson noise. If the value of the resistance is known then the temperature can be found.[14][15]
|
50 |
+
|
51 |
+
Historically, till May 2019, the definition of the Kelvin scale was that invented by Kelvin, based on a ratio of quantities of energy in processes in an ideal Carnot engine, entirely in terms of macroscopic thermodynamics. That Carnot engine was to work between two temperatures, that of the body whose temperature was to be measured, and a reference, that of a body at the temperature of the triple point of water. Then the reference temperature, that of the triple point, was defined to be exactly 273.16 K. Since May 2019, that value has not been fixed by definition, but is to be measured through microscopic phenomena, involving the Boltzmann constant, as described above. The microscopic statistical mechanical definition does not have a reference temperature.
|
52 |
+
|
53 |
+
A material on which a macroscopically defined temperature scale may be based is the ideal gas. The pressure exerted by a fixed volume and mass of an ideal gas is directly proportional to its temperature. Some natural gases show so nearly ideal properties over suitable temperature ranges that they can be used for thermometry; this was important during the development of thermodynamics and is still of practical importance today.[16][17] The ideal gas thermometer is, however, not theoretically perfect for thermodynamics. This is because the entropy of an ideal gas at its absolute zero of temperature is not a positive semi-definite quantity, which puts the gas in violation of the third law of thermodynamics. In contrast to real materials, the ideal gas does not liquefy or solidify, no matter how cold it is. Alternatively thinking, the ideal gas law, refers to the limit of infinitely high temperature and zero pressure; these conditions guarantee non-interactive motions of the constituent molecules.[18][19][20]
|
54 |
+
|
55 |
+
The magnitude of the kelvin is now defined in terms of kinetic theory, derived from the value of Boltzmann's constant.
|
56 |
+
|
57 |
+
Kinetic theory provides a microscopic account of temperature for some bodies of material, especially gases, based on macroscopic systems' being composed of many microscopic particles, such as molecules and ions of various species, the particles of a species being all alike. It explains macroscopic phenomena through the classical mechanics of the microscopic particles. The equipartition theorem of kinetic theory asserts that each classical degree of freedom of a freely moving particle has an average kinetic energy of kBT/2 where kB denotes Boltzmann's constant. The translational motion of the particle has three degrees of freedom, so that, except at very low temperatures where quantum effects predominate, the average translational kinetic energy of a freely moving particle in a system with temperature T will be 3kBT/2.
|
58 |
+
|
59 |
+
Molecules, such as oxygen (O2), have more degrees of freedom than single spherical atoms: they undergo rotational and vibrational motions as well as translations. Heating results in an increase in temperature due to an increase in the average translational kinetic energy of the molecules. Heating will also cause, through equipartitioning, the energy associated with vibrational and rotational modes to increase. Thus a diatomic gas will require more energy input to increase its temperature by a certain amount, i.e. it will have a greater heat capacity than a monatomic gas.
|
60 |
+
|
61 |
+
As noted above, the speed of sound in a gas can be calculated from the molecular character of the gas, from its temperature and pressure, and from the value of Boltzmann's constant. Taking the value of Boltzmann's constant as a primarily defined reference of exactly defined value, a measurement of the speed of sound can provide a more precise measurement of the temperature of the gas.[13]
|
62 |
+
|
63 |
+
It is possible to measure the average kinetic energy of constituent microscopic particles if they are allowed to escape from the bulk of the system, through a small hole in the containing wall. The spectrum of velocities has to be measured, and the average calculated from that. It is not necessarily the case that the particles that escape and are measured have the same velocity distribution as the particles that remain in the bulk of the system, but sometimes a good sample is possible.
|
64 |
+
|
65 |
+
Temperature is one of the principal quantities in the study of thermodynamics. Formerly, the magnitude of the kelvin was defined in thermodynamic terms, but nowadays, as mentioned above, it is defined in terms of kinetic theory.
|
66 |
+
|
67 |
+
The thermodynamic temperature is said to be absolute for two reasons. One is that its formal character is independent of the properties of particular materials. The other reason is that its zero is, in a sense, absolute, in that it indicates absence of microscopic classical motion of the constituent particles of matter, so that they have a limiting specific heat of zero for zero temperature, according to the third law of thermodynamics. Nevertheless, a thermodynamic temperature does in fact have a definite numerical value that has been arbitrarily chosen by tradition and is dependent on the property of a particular materials; it is simply less arbitrary than relative "degrees" scales such as Celsius and Fahrenheit. Being an absolute scale with one fixed point (zero), there is only one degree of freedom left to arbitrary choice, rather than two as in relative scales. For the Kelvin scale since May 2019, by international convention, the choice has been made to use knowledge of modes of operation of various thermometric devices, relying on microscopic kinetic theories about molecular motion. The numerical scale is settled by a conventional definition of the value of the Boltzmann constant, which relates macroscopic temperature to average microscopic kinetic energy of particles such as molecules. Its numerical value is arbitrary, and an alternate, less widely used absolute temperature scale exists called the Rankine scale, made to be aligned with the Fahrenheit scale as Kelvin is with Celsius.
|
68 |
+
|
69 |
+
The thermodynamic definition of temperature is due to Kelvin. It is framed in terms of an idealized device called a Carnot engine, imagined to run in a fictive continuous cycle of successive processes that traverse a cycle of states of its working body. The engine takes in a quantity of heat Q1 from a hot reservoir and passes out a lesser quantity of heat Q2 to a cold reservoir. The difference in energy is passed, as thermodynamic work, to a work reservoir, and is considered to be the output of the engine. The cycle is imagined to run so slowly that at each point of the cycle the working body is in a state of thermodynamic equilibrium. The successive processes the cycle are thus imagined to run reversibly with no entropy production. Then the quantity of entropy taken in from the hot reservoir when the working body is heated is equal to that passed to the cold reservoir when the working body is cooled. Then the absolute or thermodynamic temperatures, T1 and T2, of the reservoirs are defined so that to be such that
|
70 |
+
|
71 |
+
|
72 |
+
|
73 |
+
|
74 |
+
|
75 |
+
|
76 |
+
|
77 |
+
|
78 |
+
|
79 |
+
(1)
|
80 |
+
|
81 |
+
The zeroth law of thermodynamics allows this definition to be used to measure the absolute or thermodynamic temperature of an arbitrary body of interest, by making the other heat reservoir have the same temperature as the body of interest.
|
82 |
+
|
83 |
+
Kelvin's original work postulating absolute temperature was published in 1848. It was based on the work of Carnot, before the formulation of the first law of thermodynamics. Carnot had no sound understanding of heat, and no specific concept of entropy. He wrote of 'caloric', and said that all the caloric that passed from the hot reservoir was passed into the cold reservoir. Kelvin wrote in his 1848 paper that his scale was absolute in the sense that it was defined "independently of the properties of any particular kind of matter". His definitive publication, which sets out the definition just stated, was printed in 1853, a paper read in 1851.[21][22][23][24]
|
84 |
+
|
85 |
+
Numerical details were formerly settled by making one of the heat reservoirs a cell at the triple point of water, which was defined to have an absolute temperature of 273.16 K.[25] Nowadays, the numerical value is instead obtained from measurement through the microscopic statistical mechanical international definition, as above.
|
86 |
+
|
87 |
+
In thermodynamic terms, temperature is an intensive variable because it is equal to a differential coefficient of one extensive variable with respect to another, for a given body. It thus has the dimensions of a ratio of two extensive variables. In thermodynamics, two bodies are often considered as connected by contact with a common wall, which has some specific permeability properties. Such specific permeability can be referred to a specific intensive variable. An example is a diathermic wall that is permeable only to heat; the intensive variable for this case is temperature. When the two bodies have been in contact for a very long time, and have settled to a permanent steady state, the relevant intensive variables are equal in the two bodies; for a diathermal wall, this statement is sometimes called the zeroth law of thermodynamics.[26][27][28]
|
88 |
+
|
89 |
+
In particular, when the body is described by stating its internal energy U, an extensive variable, as a function of its entropy S, also an extensive variable, and other state variables V, N, with U = U (S, V, N), then the temperature is equal to the partial derivative of the internal energy with respect to the entropy:[27][28][29]
|
90 |
+
|
91 |
+
|
92 |
+
|
93 |
+
|
94 |
+
|
95 |
+
|
96 |
+
|
97 |
+
|
98 |
+
|
99 |
+
(2)
|
100 |
+
|
101 |
+
Likewise, when the body is described by stating its entropy S as a function of its internal energy U, and other state variables V, N, with S = S (U, V, N), then the reciprocal of the temperature is equal to the partial derivative of the entropy with respect to the internal energy:[27][29][30]
|
102 |
+
|
103 |
+
|
104 |
+
|
105 |
+
|
106 |
+
|
107 |
+
|
108 |
+
|
109 |
+
|
110 |
+
|
111 |
+
(3)
|
112 |
+
|
113 |
+
The above definition, equation (1), of the absolute temperature is due to Kelvin. It refers to systems closed to transfer of matter, and has special emphasis on directly experimental procedures. A presentation of thermodynamics by Gibbs starts at a more abstract level and deals with systems open to the transfer of matter; in this development of thermodynamics, the equations (2) and (3) above are actually alternative definitions of temperature.[31]
|
114 |
+
|
115 |
+
Real world bodies are often not in thermodynamic equilibrium and not homogeneous. For study by methods of classical irreversible thermodynamics, a body is usually spatially and temporally divided conceptually into 'cells' of small size. If classical thermodynamic equilibrium conditions for matter are fulfilled to good approximation in such a 'cell', then it is homogeneous and a temperature exists for it. If this is so for every 'cell' of the body, then local thermodynamic equilibrium is said to prevail throughout the body.[32][33][34][35][36]
|
116 |
+
|
117 |
+
It makes good sense, for example, to say of the extensive variable U, or of the extensive variable S, that it has a density per unit volume, or a quantity per unit mass of the system, but it makes no sense to speak of density of temperature per unit volume or quantity of temperature per unit mass of the system. On the other hand, it makes no sense to speak of the internal energy at a point, while when local thermodynamic equilibrium prevails, it makes good sense to speak of the temperature at a point. Consequently, temperature can vary from point to point in a medium that is not in global thermodynamic equilibrium, but in which there is local thermodynamic equilibrium.
|
118 |
+
|
119 |
+
Thus, when local thermodynamic equilibrium prevails in a body, temperature can be regarded as a spatially varying local property in that body, and this is because temperature is an intensive variable.
|
120 |
+
|
121 |
+
Temperature is a measure of a quality of a state of a material.[37] The quality may be regarded as a more abstract entity than any particular temperature scale that measures it, and is called hotness by some writers.[38] The quality of hotness refers to the state of material only in a particular locality, and in general, apart from bodies held in a steady state of thermodynamic equilibrium, hotness varies from place to place. It is not necessarily the case that a material in a particular place is in a state that is steady and nearly homogeneous enough to allow it to have a well-defined hotness or temperature. Hotness may be represented abstractly as a one-dimensional manifold. Every valid temperature scale has its own one-to-one map into the hotness manifold.[39][40]
|
122 |
+
|
123 |
+
When two systems in thermal contact are at the same temperature no heat transfers between them. When a temperature difference does exist heat flows spontaneously from the warmer system to the colder system until they are in thermal equilibrium. Such heat transfer occurs by conduction or by thermal radiation.[41][42][43][44][45][46][47][48]
|
124 |
+
|
125 |
+
Experimental physicists, for example Galileo and Newton,[49] found that there are indefinitely many empirical temperature scales. Nevertheless, the zeroth law of thermodynamics says that they all measure the same quality. This means that for a body in its own state of internal thermodynamic equilibrium, every correctly calibrated thermometer, of whatever kind, that measures the temperature of the body, records one and the same temperature. For a body that is not in its own state of internal thermodynamic equilibrium, different thermometers can record different temperatures, depending respectively on the mechanisms of operation of the thermometers.
|
126 |
+
|
127 |
+
For experimental physics, hotness means that, when comparing any two given bodies in their respective separate thermodynamic equilibria, any two suitably given empirical thermometers with numerical scale readings will agree as to which is the hotter of the two given bodies, or that they have the same temperature.[50] This does not require the two thermometers to have a linear relation between their numerical scale readings, but it does require that the relation between their numerical readings shall be strictly monotonic.[51][52] A definite sense of greater hotness can be had, independently of calorimetry, of thermodynamics, and of properties of particular materials, from Wien's displacement law of thermal radiation: the temperature of a bath of thermal radiation is proportional, by a universal constant, to the frequency of the maximum of its frequency spectrum; this frequency is always positive, but can have values that tend to zero. Thermal radiation is initially defined for a cavity in thermodynamic equilibrium. These physical facts justify a mathematical statement that hotness exists on an ordered one-dimensional manifold. This is a fundamental character of temperature and thermometers for bodies in their own thermodynamic equilibrium.[7][39][40][53][54]
|
128 |
+
|
129 |
+
Except for a system undergoing a first-order phase change such as the melting of ice, as a closed system receives heat, without change in its volume and without change in external force fields acting on it, its temperature rises. For a system undergoing such a phase change so slowly that departure from thermodynamic equilibrium can be neglected, its temperature remains constant as the system is supplied with latent heat. Conversely, a loss of heat from a closed system, without phase change, without change of volume, and without change in external force fields acting on it, decreases its temperature.[55]
|
130 |
+
|
131 |
+
While for bodies in their own thermodynamic equilibrium states, the notion of temperature requires that all empirical thermometers must agree as to which of two bodies is the hotter or that they are at the same temperature, this requirement is not safe for bodies that are in steady states though not in thermodynamic equilibrium. It can then well be that different empirical thermometers disagree about which is the hotter, and if this is so, then at least one of the bodies does not have a well defined absolute thermodynamic temperature. Nevertheless, any one given body and any one suitable empirical thermometer can still support notions of empirical, non-absolute, hotness and temperature, for a suitable range of processes. This is a matter for study in non-equilibrium thermodynamics.
|
132 |
+
|
133 |
+
When a body is not in a steady state, then the notion of temperature becomes even less safe than for a body in a steady state not in thermodynamic equilibrium. This is also a matter for study in non-equilibrium thermodynamics.
|
134 |
+
|
135 |
+
For axiomatic treatment of thermodynamic equilibrium, since the 1930s, it has become customary to refer to a zeroth law of thermodynamics. The customarily stated minimalist version of such a law postulates only that all bodies, which when thermally connected would be in thermal equilibrium, should be said to have the same temperature by definition, but by itself does not establish temperature as a quantity expressed as a real number on a scale. A more physically informative version of such a law views empirical temperature as a chart on a hotness manifold.[39][54][56] While the zeroth law permits the definitions of many different empirical scales of temperature, the second law of thermodynamics selects the definition of a single preferred, absolute temperature, unique up to an arbitrary scale factor, whence called the thermodynamic temperature.[7][39][57][58][59][60] If internal energy is considered as a function of the volume and entropy of a homogeneous system in thermodynamic equilibrium, thermodynamic absolute temperature appears as the partial derivative of internal energy with respect the entropy at constant volume. Its natural, intrinsic origin or null point is absolute zero at which the entropy of any system is at a minimum. Although this is the lowest absolute temperature described by the model, the third law of thermodynamics postulates that absolute zero cannot be attained by any physical system.
|
136 |
+
|
137 |
+
When an energy transfer to or from a body is only as heat, the state of the body changes. Depending on the surroundings and the walls separating them from the body, various changes are possible in the body. They include chemical reactions, increase of pressure, increase of temperature, and phase change. For each kind of change under specified conditions, the heat capacity is the ratio of the quantity of heat transferred to the magnitude of the change. For example, if the change is an increase in temperature at constant volume, with no phase change and no chemical change, then the temperature of the body rises and its pressure increases. The quantity of heat transferred, ΔQ, divided by the observed temperature change, ΔT, is the body's heat capacity at constant volume:
|
138 |
+
|
139 |
+
If heat capacity is measured for a well defined amount of substance, the specific heat is the measure of the heat required to increase the temperature of such a unit quantity by one unit of temperature. For example, to raise the temperature of water by one kelvin (equal to one degree Celsius) requires 4186 joules per kilogram (J/kg).
|
140 |
+
|
141 |
+
Temperature measurement using modern scientific thermometers and temperature scales goes back at least as far as the early 18th century, when Gabriel Fahrenheit adapted a thermometer (switching to mercury) and a scale both developed by Ole Christensen Rømer. Fahrenheit's scale is still in use in the United States for non-scientific applications.
|
142 |
+
|
143 |
+
Temperature is measured with thermometers that may be calibrated to a variety of temperature scales. In most of the world (except for Belize, Myanmar, Liberia and the United States), the Celsius scale is used for most temperature measuring purposes. Most scientists measure temperature using the Celsius scale and thermodynamic temperature using the Kelvin scale, which is the Celsius scale offset so that its null point is 0 K = −273.15 °C, or absolute zero. Many engineering fields in the US, notably high-tech and US federal specifications (civil and military), also use the Kelvin and Celsius scales. Other engineering fields in the US also rely upon the Rankine scale (a shifted Fahrenheit scale) when working in thermodynamic-related disciplines such as combustion.
|
144 |
+
|
145 |
+
The basic unit of temperature in the International System of Units (SI) is the Kelvin. It has the symbol K.
|
146 |
+
|
147 |
+
For everyday applications, it is often convenient to use the Celsius scale, in which 0 °C corresponds very closely to the freezing point of water and 100 °C is its boiling point at sea level. Because liquid droplets commonly exist in clouds at sub-zero temperatures, 0 °C is better defined as the melting point of ice. In this scale a temperature difference of 1 degree Celsius is the same as a 1kelvin increment, but the scale is offset by the temperature at which ice melts (273.15 K).
|
148 |
+
|
149 |
+
By international agreement,[61] until May 2019, the Kelvin and Celsius scales were defined by two fixing points: absolute zero and the triple point of Vienna Standard Mean Ocean Water, which is water specially prepared with a specified blend of hydrogen and oxygen isotopes. Absolute zero was defined as precisely 0 K and −273.15 °C. It is the temperature at which all classical translational motion of the particles comprising matter ceases and they are at complete rest in the classical model. Quantum-mechanically, however, zero-point motion remains and has an associated energy, the zero-point energy. Matter is in its ground state,[62] and contains no thermal energy. The triple point of water is defined as 273.16 K and 0.01 °C. This definition served the following purposes: it fixed the magnitude of the kelvin as being precisely 1 part in 273.16 parts of the difference between absolute zero and the triple point of water; it establishes that one kelvin has precisely the same magnitude as one degree on the Celsius scale; and it established the difference between the null points of these scales as being 273.15 K (0 K = −273.15 °C and 273.16 K = 0.01 °C). Since 2019, there has been a new definition based on the Boltzmann constant,[63] but the scales are scarcely changed.
|
150 |
+
|
151 |
+
In the United States, the Fahrenheit scale is widely used. On this scale the freezing point of water corresponds to 32 °F and the boiling point to 212 °F. The Rankine scale, still used in fields of chemical engineering in the US, is an absolute scale based on the Fahrenheit increment.
|
152 |
+
|
153 |
+
The following table shows the temperature conversion formulas for conversions to and from the Celsius scale.
|
154 |
+
|
155 |
+
The field of plasma physics deals with phenomena of electromagnetic nature that involve very high temperatures. It is customary to express temperature as energy in units of electronvolts (eV) or kiloelectronvolts (keV). The energy, which has a different dimension from temperature, is then calculated as the product of the Boltzmann constant and temperature,
|
156 |
+
|
157 |
+
|
158 |
+
|
159 |
+
E
|
160 |
+
=
|
161 |
+
|
162 |
+
k
|
163 |
+
|
164 |
+
B
|
165 |
+
|
166 |
+
|
167 |
+
T
|
168 |
+
|
169 |
+
|
170 |
+
{\displaystyle E=k_{\text{B}}T}
|
171 |
+
|
172 |
+
. Then, 1 eV corresponds to 11605 K. In the study of QCD matter one routinely encounters temperatures of the order of a few hundred MeV, equivalent to about 1012 K.
|
173 |
+
|
174 |
+
Historically, there are several scientific approaches to the explanation of temperature: the classical thermodynamic description based on macroscopic empirical variables that can be measured in a laboratory; the kinetic theory of gases which relates the macroscopic description to the probability distribution of the energy of motion of gas particles; and a microscopic explanation based on statistical physics and quantum mechanics. In addition, rigorous and purely mathematical treatments have provided an axiomatic approach to classical thermodynamics and temperature.[64] Statistical physics provides a deeper understanding by describing the atomic behavior of matter, and derives macroscopic properties from statistical averages of microscopic states, including both classical and quantum states. In the fundamental physical description, using natural units, temperature may be measured directly in units of energy. However, in the practical systems of measurement for science, technology, and commerce, such as the modern metric system of units, the macroscopic and the microscopic descriptions are interrelated by the Boltzmann constant, a proportionality factor that scales temperature to the microscopic mean kinetic energy.
|
175 |
+
|
176 |
+
The microscopic description in statistical mechanics is based on a model that analyzes a system into its fundamental particles of matter or into a set of classical or quantum-mechanical oscillators and considers the system as a statistical ensemble of microstates. As a collection of classical material particles, temperature is a measure of the mean energy of motion, called kinetic energy, of the particles, whether in solids, liquids, gases, or plasmas. The kinetic energy, a concept of classical mechanics, is half the mass of a particle times its speed squared. In this mechanical interpretation of thermal motion, the kinetic energies of material particles may reside in the velocity of the particles of their translational or vibrational motion or in the inertia of their rotational modes. In monatomic perfect gases and, approximately, in most gases, temperature is a measure of the mean particle kinetic energy. It also determines the probability distribution function of the energy. In condensed matter, and particularly in solids, this purely mechanical description is often less useful and the oscillator model provides a better description to account for quantum mechanical phenomena. Temperature determines the statistical occupation of the microstates of the ensemble. The microscopic definition of temperature is only meaningful in the thermodynamic limit, meaning for large ensembles of states or particles, to fulfill the requirements of the statistical model.
|
177 |
+
|
178 |
+
The kinetic energy is also considered as a component of thermal energy. The thermal energy may be partitioned into independent components attributed to the degrees of freedom of the particles or to the modes of oscillators in a thermodynamic system. In general, the number of these degrees of freedom that are available for the equipartitioning of energy depends on the temperature, i.e. the energy region of the interactions under consideration. For solids, the thermal energy is associated primarily with the vibrations of its atoms or molecules about their equilibrium position. In an ideal monatomic gas, the kinetic energy is found exclusively in the purely translational motions of the particles. In other systems, vibrational and rotational motions also contribute degrees of freedom.
|
179 |
+
|
180 |
+
Maxwell and Boltzmann developed a kinetic theory that yields a fundamental understanding of temperature in gases.[65]
|
181 |
+
This theory also explains the ideal gas law and the observed heat capacity of monatomic (or 'noble') gases.[66][67][68]
|
182 |
+
|
183 |
+
The ideal gas law is based on observed empirical relationships between pressure (p), volume (V), and temperature (T), and was recognized long before the kinetic theory of gases was developed (see Boyle's and Charles's laws). The ideal gas law states:[69]
|
184 |
+
|
185 |
+
where n is the number of moles of gas and R = 8.314462618... J⋅mol−1⋅K−1[70] is the gas constant.
|
186 |
+
|
187 |
+
This relationship gives us our first hint that there is an absolute zero on the temperature scale, because it only holds if the temperature is measured on an absolute scale such as Kelvin's. The ideal gas law allows one to measure temperature on this absolute scale using the gas thermometer. The temperature in kelvins can be defined as the pressure in pascals of one mole of gas in a container of one cubic meter, divided by the gas constant.
|
188 |
+
|
189 |
+
Although it is not a particularly convenient device, the gas thermometer provides an essential theoretical basis by which all thermometers can be calibrated. As a practical matter, it is not possible to use a gas thermometer to measure absolute zero temperature since the gases tend to condense into a liquid long before the temperature reaches zero. It is possible, however, to extrapolate to absolute zero by using the ideal gas law, as shown in the figure.
|
190 |
+
|
191 |
+
The kinetic theory assumes that pressure is caused by the force associated with individual atoms striking the walls, and that all energy is translational kinetic energy. Using a sophisticated symmetry argument,[71] Boltzmann deduced what is now called the Maxwell–Boltzmann probability distribution function for the velocity of particles in an ideal gas. From that probability distribution function, the average kinetic energy (per particle) of a monatomic ideal gas is[67][72]
|
192 |
+
|
193 |
+
where the Boltzmann constant kB is the ideal gas constant divided by the Avogadro number, and
|
194 |
+
|
195 |
+
|
196 |
+
|
197 |
+
|
198 |
+
v
|
199 |
+
|
200 |
+
rms
|
201 |
+
|
202 |
+
|
203 |
+
=
|
204 |
+
|
205 |
+
|
206 |
+
|
207 |
+
⟨
|
208 |
+
|
209 |
+
v
|
210 |
+
|
211 |
+
2
|
212 |
+
|
213 |
+
|
214 |
+
⟩
|
215 |
+
|
216 |
+
|
217 |
+
|
218 |
+
|
219 |
+
|
220 |
+
{\displaystyle v_{\text{rms}}={\sqrt {\left\langle v^{2}\right\rangle }}}
|
221 |
+
|
222 |
+
is the root-mean-square speed. Thus the ideal gas law states that internal energy is directly proportional to temperature.[73] This direct proportionality between temperature and internal energy is a special case of the equipartition theorem, and holds only in the classical limit of an ideal gas. It does not hold for most substances, although it is true that temperature is a monotonic (non-decreasing) function of internal energy.
|
223 |
+
|
224 |
+
When two otherwise isolated bodies are connected together by a rigid physical path impermeable to matter, there is spontaneous transfer of energy as heat from the hotter to the colder of them. Eventually, they reach a state of mutual thermal equilibrium, in which heat transfer has ceased, and the bodies' respective state variables have settled to become unchanging.
|
225 |
+
|
226 |
+
One statement of the zeroth law of thermodynamics is that if two systems are each in thermal equilibrium with a third system, then they are also in thermal equilibrium with each other.
|
227 |
+
|
228 |
+
This statement helps to define temperature but it does not, by itself, complete the definition. An empirical temperature is a numerical scale for the hotness of a thermodynamic system. Such hotness may be defined as existing on a one-dimensional manifold, stretching between hot and cold. Sometimes the zeroth law is stated to include the existence of a unique universal hotness manifold, and of numerical scales on it, so as to provide a complete definition of empirical temperature.[56] To be suitable for empirical thermometry, a material must have a monotonic relation between hotness and some easily measured state variable, such as pressure or volume, when all other relevant coordinates are fixed. An exceptionally suitable system is the ideal gas, which can provide a temperature scale that matches the absolute Kelvin scale. The Kelvin scale is defined on the basis of the second law of thermodynamics.
|
229 |
+
|
230 |
+
As an alternative to considering or defining the zeroth law of thermodynamics, it was the historical development in thermodynamics to define temperature in terms of the second law of thermodynamics which deals with entropy. The second law states that any process will result in either no change or a net increase in the entropy of the universe. This can be understood in terms of probability.
|
231 |
+
|
232 |
+
For example, in a series of coin tosses, a perfectly ordered system would be one in which either every toss comes up heads or every toss comes up tails. This means the outcome is always 100% the same result. In contrast, many mixed (disordered) outcomes are possible, and their number increases with each toss. Eventually, the combinations of ~50% heads and ~50% tails dominate and obtaining an outcome significantly different from 50/50 becomes increasingly unlikely. Thus the system naturally progresses to a state of maximum disorder or entropy.
|
233 |
+
|
234 |
+
As temperature governs the transfer of heat between two systems and the universe tends to progress toward a maximum of entropy, it is expected that there is some relationship between temperature and entropy. A heat engine is a device for converting thermal energy into mechanical energy, resulting in the performance of work. and analysis of the Carnot heat engine provides the necessary relationships. The work from a heat engine corresponds to the difference between the heat put into the system at high temperature, qH and the heat extracted at the low temperature, qC. The efficiency is the work divided by the heat input:
|
235 |
+
|
236 |
+
|
237 |
+
|
238 |
+
|
239 |
+
|
240 |
+
|
241 |
+
|
242 |
+
|
243 |
+
|
244 |
+
(4)
|
245 |
+
|
246 |
+
where wcy is the work done per cycle. The efficiency depends only on qC/qH. Because qC and qH correspond to heat transfer at the temperatures TC and TH respectively, qC/qH should be some function of these temperatures:
|
247 |
+
|
248 |
+
|
249 |
+
|
250 |
+
|
251 |
+
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
|
256 |
+
(5)
|
257 |
+
|
258 |
+
Carnot's theorem states that all reversible engines operating between the same heat reservoirs are equally efficient. Thus, a heat engine operating between T1 and T3 must have the same efficiency as one consisting of two cycles, one between T1 and T2, and the second between T2 and T3. This can only be the case if
|
259 |
+
|
260 |
+
which implies
|
261 |
+
|
262 |
+
Since the first function is independent of T2, this temperature must cancel on the right side, meaning f(T1, T3) is of the form g(T1)/g(T3) (i.e. f(T1, T3) = f(T1, T2)f(T2, T3) = g(T1)/g(T2) · g(T2)/g(T3) = g(T1)/g(T3)), where g is a function of a single temperature. A temperature scale can now be chosen with the property that
|
263 |
+
|
264 |
+
|
265 |
+
|
266 |
+
|
267 |
+
|
268 |
+
|
269 |
+
|
270 |
+
|
271 |
+
|
272 |
+
(6)
|
273 |
+
|
274 |
+
Substituting (6) back into (4) gives a relationship for the efficiency in terms of temperature:
|
275 |
+
|
276 |
+
|
277 |
+
|
278 |
+
|
279 |
+
|
280 |
+
|
281 |
+
|
282 |
+
|
283 |
+
|
284 |
+
(7)
|
285 |
+
|
286 |
+
For TC = 0 K the efficiency is 100% and that efficiency becomes greater than 100% below 0 K. Since an efficiency greater than 100% violates the first law of thermodynamics, this implies that 0 K is the minimum possible temperature. In fact the lowest temperature ever obtained in a macroscopic system was 20 nK, which was achieved in 1995 at NIST. Subtracting the right hand side of (5) from the middle portion and rearranging gives
|
287 |
+
|
288 |
+
where the negative sign indicates heat ejected from the system. This relationship suggests the existence of a state function, S, defined by
|
289 |
+
|
290 |
+
|
291 |
+
|
292 |
+
|
293 |
+
|
294 |
+
|
295 |
+
|
296 |
+
|
297 |
+
|
298 |
+
(8)
|
299 |
+
|
300 |
+
where the subscript indicates a reversible process. The change of this state function around any cycle is zero, as is necessary for any state function. This function corresponds to the entropy of the system, which was described previously. Rearranging (8) gives a formula for temperature in terms of fictive infinitesimal quasi-reversible elements of entropy and heat:
|
301 |
+
|
302 |
+
|
303 |
+
|
304 |
+
|
305 |
+
|
306 |
+
|
307 |
+
|
308 |
+
|
309 |
+
|
310 |
+
(9)
|
311 |
+
|
312 |
+
For a system, where entropy S(E) is a function of its energy E, the temperature T is given by
|
313 |
+
|
314 |
+
|
315 |
+
|
316 |
+
|
317 |
+
|
318 |
+
|
319 |
+
|
320 |
+
|
321 |
+
|
322 |
+
(10)
|
323 |
+
|
324 |
+
i.e. the reciprocal of the temperature is the rate of increase of entropy with respect to energy.
|
325 |
+
|
326 |
+
Statistical mechanics defines temperature based on a system's fundamental degrees of freedom. Eq.(10) is the defining relation of temperature, where the entropy
|
327 |
+
|
328 |
+
|
329 |
+
|
330 |
+
S
|
331 |
+
|
332 |
+
|
333 |
+
{\displaystyle S}
|
334 |
+
|
335 |
+
is defined (up to a constant) by the logarithm of the number of microstates of the system in the given macrostate (as specified in the microcanonical ensemble):
|
336 |
+
|
337 |
+
where
|
338 |
+
|
339 |
+
|
340 |
+
|
341 |
+
|
342 |
+
k
|
343 |
+
|
344 |
+
|
345 |
+
B
|
346 |
+
|
347 |
+
|
348 |
+
|
349 |
+
|
350 |
+
|
351 |
+
{\displaystyle k_{\mathrm {B} }}
|
352 |
+
|
353 |
+
is Boltzmann's constant and N is the number of microstates.
|
354 |
+
|
355 |
+
When two systems with different temperatures are put into purely thermal connection, heat will flow from the higher temperature system to the lower temperature one; thermodynamically this is understood by the second law of thermodynamics: The total change in entropy following a transfer of energy
|
356 |
+
|
357 |
+
|
358 |
+
|
359 |
+
Δ
|
360 |
+
E
|
361 |
+
|
362 |
+
|
363 |
+
{\displaystyle \Delta E}
|
364 |
+
|
365 |
+
from system 1 to system 2 is:
|
366 |
+
|
367 |
+
and is thus positive if
|
368 |
+
|
369 |
+
|
370 |
+
|
371 |
+
|
372 |
+
T
|
373 |
+
|
374 |
+
1
|
375 |
+
|
376 |
+
|
377 |
+
>
|
378 |
+
|
379 |
+
T
|
380 |
+
|
381 |
+
2
|
382 |
+
|
383 |
+
|
384 |
+
|
385 |
+
|
386 |
+
{\displaystyle T_{1}>T_{2}}
|
387 |
+
|
388 |
+
From the point of view of statistical mechanics, the total number of micsostates in the combined system 1 + system 2 is
|
389 |
+
|
390 |
+
|
391 |
+
|
392 |
+
|
393 |
+
N
|
394 |
+
|
395 |
+
1
|
396 |
+
|
397 |
+
|
398 |
+
⋅
|
399 |
+
|
400 |
+
N
|
401 |
+
|
402 |
+
2
|
403 |
+
|
404 |
+
|
405 |
+
|
406 |
+
|
407 |
+
{\displaystyle N_{1}\cdot N_{2}}
|
408 |
+
|
409 |
+
, the logarithm of which (times Boltzmann's constant) is the sum of their entropies; thus a flow of heat from high to low temperature, which brings an increase in total entropy, is more likely than any other scenario (normally it is much more likely), as there are more microstates in the resulting macrostate.
|
410 |
+
|
411 |
+
It is possible to extend the definition of temperature even to systems of few particles, like in a quantum dot. The generalized temperature is obtained by considering time ensembles instead of configuration-space ensembles given in statistical mechanics in the case of thermal and particle exchange between a small system of fermions (N even less than 10) with a single/double-occupancy system. The finite quantum grand canonical ensemble,[74] obtained under the hypothesis of ergodicity and orthodicity,[75] allows expressing the generalized temperature from the ratio of the average time of occupation
|
412 |
+
|
413 |
+
|
414 |
+
|
415 |
+
|
416 |
+
τ
|
417 |
+
|
418 |
+
1
|
419 |
+
|
420 |
+
|
421 |
+
|
422 |
+
|
423 |
+
{\displaystyle \tau _{1}}
|
424 |
+
|
425 |
+
and
|
426 |
+
|
427 |
+
|
428 |
+
|
429 |
+
|
430 |
+
τ
|
431 |
+
|
432 |
+
2
|
433 |
+
|
434 |
+
|
435 |
+
|
436 |
+
|
437 |
+
{\displaystyle \tau _{2}}
|
438 |
+
|
439 |
+
of the single/double-occupancy system:[76]
|
440 |
+
|
441 |
+
where EF is the Fermi energy. This generalized temperature tends to the ordinary temperature when N goes to infinity.
|
442 |
+
|
443 |
+
On the empirical temperature scales that are not referenced to absolute zero, a negative temperature is one below the zero-point of the scale used. For example, dry ice has a sublimation temperature of −78.5 °C which is equivalent to −109.3 °F. On the absolute kelvin scale this temperature is 194.6 K. No body can be brought to exactly 0 K (the temperature of the ideally coldest possible body) by any finite practicable process; this is a consequence of the third law of thermodynamics.
|
444 |
+
|
445 |
+
The international kinetic theory temperature of a body cannot take negative values. The thermodynamic temperature scale, however, is not so constrained.
|
446 |
+
|
447 |
+
For a body of matter, there can sometimes be conceptually defined, in terms of microscopic degrees of freedom, namely particle spins, a subsystem, with a temperature other than that of the whole body. When the body is in its own state of internal thermodynamic equilibrium, the temperatures of the whole body and of the subsystem must be the same. The two temperatures can differ when, by work through externally imposed force fields, energy can be transferred to and from the subsystem, separately from the rest of the body; then the whole body is not in its own state of internal thermodynamic equilibrium. There is an upper limit of energy such a spin subsystem can attain.
|
448 |
+
|
449 |
+
Considering the subsystem to be in a temporary state of virtual thermodynamic equilibrium, it is possible to obtain a negative temperature on the thermodynamic scale. Thermodynamic temperature is the inverse of the derivative of the subsystem's entropy with respect to its internal energy. As the subsystem's internal energy increases, the entropy increases for some range, but eventually attains a maximum value and then begins to decrease as the highest energy states begin to fill. At the point of maximum entropy, the temperature function shows the behavior of a singularity, because the slope of the entropy function decreases to zero and then turns negative. As the subsystem's entropy reaches its maximum, its thermodynamic temperature goes to positive infinity, switching to negative infinity as the slope turns negative. Such negative temperatures are hotter than any positive temperature. Over time, when the subsystem is exposed to the rest of the body, which has a positive temperature, energy is transferred as heat from the negative temperature subsystem to the positive temperature system.[77] The kinetic theory temperature is not defined for such subsystems.
|
en/5644.html.txt
ADDED
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|
1 |
+
|
2 |
+
|
3 |
+
Temperature is a physical property of matter that quantitatively expresses hot and cold. It is the manifestation of thermal energy, present in all matter, which is the source of the occurrence of heat, a flow of energy, when a body is in contact with another that is colder.
|
4 |
+
|
5 |
+
Temperature is measured with a thermometer. Thermometers are calibrated in various temperature scales that historically have used various reference points and thermometric substances for definition. The most common scales are the Celsius scale (formerly called centigrade, denoted °C), the Fahrenheit scale (denoted °F), and the Kelvin scale (denoted K), the last of which is predominantly used for scientific purposes by conventions of the International System of Units (SI).
|
6 |
+
|
7 |
+
The lowest theoretical temperature is absolute zero, at which no more thermal energy can be extracted from a body. Experimentally, it can only be approached very closely, but not reached, which is recognized in the third law of thermodynamics.
|
8 |
+
|
9 |
+
Temperature is important in all fields of natural science, including physics, chemistry, Earth science, medicine, and biology, as well as most aspects of daily life.
|
10 |
+
|
11 |
+
Many physical processes are related to temperature, such as:
|
12 |
+
|
13 |
+
Temperature scales differ in two ways: the point chosen as zero degrees, and the magnitudes of incremental units or degrees on the scale.
|
14 |
+
|
15 |
+
The Celsius scale (°C) is used for common temperature measurements in most of the world. It is an empirical scale that was developed by a historical progress, which led to its zero point 0 °C being defined by the freezing point of water, and additional degrees defined so that 100 °C was the boiling point of water, both at sea-level atmospheric pressure. Because of the 100-degree interval, it was called a centigrade scale.[3] Since the standardization of the kelvin in the International System of Units, it has subsequently been redefined in terms of the equivalent fixing points on the Kelvin scale, and so that a temperature increment of one degree Celsius is the same as an increment of one kelvin, though they differ by an additive offset of approximately 273.15.
|
16 |
+
|
17 |
+
The United States commonly uses the Fahrenheit scale, on which water freezes at 32 °F and boils at 212 °F at sea-level atmospheric pressure.
|
18 |
+
|
19 |
+
At the absolute zero of temperature, no more energy can be removed from matter as heat, a fact expressed in the third law of thermodynamics. At this temperature, matter contains no macroscopic thermal energy, but still has quantum-mechanical zero-point energy as predicted by the uncertainty principle. This does not enter into the definition of absolute temperature. Experimentally, absolute zero can only be approached very closely, but can never be actually reached. If it were possible to cool a system to absolute zero, all classical motion of its particles would cease and they would be at complete rest in this classical sense. The absolute zero, defined as 0 K, is approximately equal to −273.15 °C, or −459.67 °F.
|
20 |
+
|
21 |
+
Referring to the Boltzmann constant, to the Maxwell–Boltzmann distribution, and to the Boltzmann statistical mechanical definition of entropy, as distinct from the Gibbs definition,[4] for independently moving microscopic particles, disregarding interparticle potential energy, by international agreement, a temperature scale is defined and said to be absolute because it is independent of the characteristics of particular thermometric substances and thermometer mechanisms. Apart from the absolute zero, it does not have a reference temperature. It is known as the Kelvin scale, widely used in science and technology. The kelvin (the word is spelled with a lower-case k) is the unit of temperature in the International System of Units (SI). The temperature of a body in its own state of thermodynamic equilibrium is always positive, relative to the absolute zero.
|
22 |
+
|
23 |
+
Besides the internationally agreed Kelvin scale, there is also a thermodynamic temperature scale, invented by Kelvin, also with its numerical zero at the absolute zero of temperature, but directly relating to purely macroscopic thermodynamic concepts, including the macroscopic entropy, though microscopically referable to the Gibbs statistical mechanical definition of entropy for the canonical ensemble, that takes interparticle potential energy into account, as well as independent particle motion, so that it can account for measurements of temperatures near absolute zero.[4] This scale has a reference temperature at the triple point of water, the numerical value of which is defined by measurements using the aforementioned internationally agreed Kelvin scale.
|
24 |
+
|
25 |
+
Many scientific measurements use the Kelvin temperature scale (unit symbol: K), named in honor of the physicist who first defined it. It is an absolute scale. Its numerical zero point, 0 K, is at the absolute zero of temperature. Since May, 2019, its degrees have been defined through particle kinetic theory, and statistical mechanics. In the International System of Units (SI), the magnitude of the kelvin is defined through various empirical measurements of the average kinetic energies of microscopic particles. It is numerically evaluated in terms of the Boltzmann constant, the value of which is defined as fixed by international convention.[5][6]
|
26 |
+
|
27 |
+
Since May 2019, the magnitude of the kelvin is defined in relation to microscopic phenomena, characterized in terms of statistical mechanics. Previously, since 1954, the International System of Units defined a scale and unit for the kelvin as a thermodynamic temperature, by using the reliably reproducible temperature of the triple point of water as a second reference point, the first reference point being 0 K at absolute zero.
|
28 |
+
|
29 |
+
Historically, the triple point temperature of water was defined as exactly 273.16 units of the measurement increment. Today it is an empirically measured quantity. The freezing point of water at sea-level atmospheric pressure occurs at approximately 273.15 K = 0 °C.
|
30 |
+
|
31 |
+
There is a variety of kinds of temperature scale. It may be convenient to classify them as empirically and theoretically based. Empirical temperature scales are historically older, while theoretically based scales arose in the middle of the nineteenth century.[7][8]
|
32 |
+
|
33 |
+
Empirically based temperature scales rely directly on measurements of simple macroscopic physical properties of materials. For example, the length of a column of mercury, confined in a glass-walled capillary tube, is dependent largely on temperature, and is the basis of the very useful mercury-in-glass thermometer. Such scales are valid only within convenient ranges of temperature. For example, above the boiling point of mercury, a mercury-in-glass thermometer is impracticable. Most materials expand with temperature increase, but some materials, such as water, contract with temperature increase over some specific range, and then they are hardly useful as thermometric materials. A material is of no use as a thermometer near one of its phase-change temperatures, for example its boiling-point.
|
34 |
+
|
35 |
+
In spite of these limitations, most generally used practical thermometers are of the empirically based kind. Especially, it was used for calorimetry, which contributed greatly to the discovery of thermodynamics. Nevertheless, empirical thermometry has serious drawbacks when judged as a basis for theoretical physics. Empirically based thermometers, beyond their base as simple direct measurements of ordinary physical properties of thermometric materials, can be re-calibrated, by use of theoretical physical reasoning, and this can extend their range of adequacy.
|
36 |
+
|
37 |
+
Theoretically based temperature scales are based directly on theoretical arguments, especially those of kinetic theory and thermodynamics. They are more or less ideally realised in practically feasible physical devices and materials. Theoretically based temperature scales are used to provide calibrating standards for practical empirically based thermometers.
|
38 |
+
|
39 |
+
In physics, the internationally agreed conventional temperature scale is called the Kelvin scale. It is calibrated through the internationally agreed and prescribed value of the Boltzmann constant,[5][6] referring to motions of microscopic particles, such as atoms, molecules, and electrons, constituent in the body whose temperature is to be measured. In contrast with the thermodynamic temperature scale invented by Kelvin, the presently conventional Kelvin temperature is not defined through comparison with the temperature of a reference state of a standard body, nor in terms of macroscopic thermodynamics.
|
40 |
+
|
41 |
+
Apart from the absolute zero of temperature, Kelvin temperature of a body in a state of internal thermodynamic equilibrium is defined by measurements of suitably chosen of its physical properties, such as have precisely known theoretical explanations in terms of the Boltzmann constant. That constant refers to chosen kinds of motion of microscopic particles in the constitution of the body. In those kinds of motion, the particles move individually, without mutual interaction. Such motions are typically interrupted by inter-particle collisions, but for temperature measurement, the motions are chosen so that, between collisions, the non-interactive segments of their trajectories are known to be accessible to accurate measurement. For this purpose, interperticle potential energy is disregarded.
|
42 |
+
|
43 |
+
In an ideal gas, and in other theoretically understood bodies, the Kelvin temperature is defined to be proportional to the average kinetic energy of non-interactively moving microscopic particles, which can be measured by suitable techniques. The proportionality constant is a simple multiple of the Boltzmann constant. If molecules, atoms, or electrons,[9][10] are emitted from a material and their velocities are measured, the spectrum of their velocities often nearly obeys a theoretical law called the Maxwell–Boltzmann distribution, which gives a well-founded measurement of temperatures for which the law holds.[11] There have not yet been successful experiments of this same kind that directly use the Fermi–Dirac distribution for thermometry, but perhaps that will be achieved in future.[12]
|
44 |
+
|
45 |
+
The speed of sound in a gas can be calculated theoretically from the molecular character of the gas, from its temperature and pressure, and from the value of Boltzmann's constant. For a gas of known molecular character and pressure, this provides a relation between temperature and Boltzmann's constant. Those quantities can be known or measured more precisely than can the thermodynamic variables that define the state of a sample of water at its triple point. Consequently, taking the value of Boltzmann's constant as a primarily defined reference of exactly defined value, a measurement of the speed of sound can provide a more precise measurement of the temperature of the gas.[13]
|
46 |
+
|
47 |
+
Measurement of the spectrum of electromagnetic radiation from an ideal three-dimensional black body can provide an accurate temperature measurement because the frequency of maximum spectral radiance of black-body radiation is directly proportional to the temperature of the black body; this is known as Wien's displacement law and has a theoretical explanation in Planck's law and the Bose–Einstein law.
|
48 |
+
|
49 |
+
Measurement of the spectrum of noise-power produced by an electrical resistor can also provide an accurate temperature measurement. The resistor has two terminals and is in effect a one-dimensional body. The Bose-Einstein law for this case indicates that the noise-power is directly proportional to the temperature of the resistor and to the value of its resistance and to the noise band-width. In a given frequency band, the noise-power has equal contributions from every frequency and is called Johnson noise. If the value of the resistance is known then the temperature can be found.[14][15]
|
50 |
+
|
51 |
+
Historically, till May 2019, the definition of the Kelvin scale was that invented by Kelvin, based on a ratio of quantities of energy in processes in an ideal Carnot engine, entirely in terms of macroscopic thermodynamics. That Carnot engine was to work between two temperatures, that of the body whose temperature was to be measured, and a reference, that of a body at the temperature of the triple point of water. Then the reference temperature, that of the triple point, was defined to be exactly 273.16 K. Since May 2019, that value has not been fixed by definition, but is to be measured through microscopic phenomena, involving the Boltzmann constant, as described above. The microscopic statistical mechanical definition does not have a reference temperature.
|
52 |
+
|
53 |
+
A material on which a macroscopically defined temperature scale may be based is the ideal gas. The pressure exerted by a fixed volume and mass of an ideal gas is directly proportional to its temperature. Some natural gases show so nearly ideal properties over suitable temperature ranges that they can be used for thermometry; this was important during the development of thermodynamics and is still of practical importance today.[16][17] The ideal gas thermometer is, however, not theoretically perfect for thermodynamics. This is because the entropy of an ideal gas at its absolute zero of temperature is not a positive semi-definite quantity, which puts the gas in violation of the third law of thermodynamics. In contrast to real materials, the ideal gas does not liquefy or solidify, no matter how cold it is. Alternatively thinking, the ideal gas law, refers to the limit of infinitely high temperature and zero pressure; these conditions guarantee non-interactive motions of the constituent molecules.[18][19][20]
|
54 |
+
|
55 |
+
The magnitude of the kelvin is now defined in terms of kinetic theory, derived from the value of Boltzmann's constant.
|
56 |
+
|
57 |
+
Kinetic theory provides a microscopic account of temperature for some bodies of material, especially gases, based on macroscopic systems' being composed of many microscopic particles, such as molecules and ions of various species, the particles of a species being all alike. It explains macroscopic phenomena through the classical mechanics of the microscopic particles. The equipartition theorem of kinetic theory asserts that each classical degree of freedom of a freely moving particle has an average kinetic energy of kBT/2 where kB denotes Boltzmann's constant. The translational motion of the particle has three degrees of freedom, so that, except at very low temperatures where quantum effects predominate, the average translational kinetic energy of a freely moving particle in a system with temperature T will be 3kBT/2.
|
58 |
+
|
59 |
+
Molecules, such as oxygen (O2), have more degrees of freedom than single spherical atoms: they undergo rotational and vibrational motions as well as translations. Heating results in an increase in temperature due to an increase in the average translational kinetic energy of the molecules. Heating will also cause, through equipartitioning, the energy associated with vibrational and rotational modes to increase. Thus a diatomic gas will require more energy input to increase its temperature by a certain amount, i.e. it will have a greater heat capacity than a monatomic gas.
|
60 |
+
|
61 |
+
As noted above, the speed of sound in a gas can be calculated from the molecular character of the gas, from its temperature and pressure, and from the value of Boltzmann's constant. Taking the value of Boltzmann's constant as a primarily defined reference of exactly defined value, a measurement of the speed of sound can provide a more precise measurement of the temperature of the gas.[13]
|
62 |
+
|
63 |
+
It is possible to measure the average kinetic energy of constituent microscopic particles if they are allowed to escape from the bulk of the system, through a small hole in the containing wall. The spectrum of velocities has to be measured, and the average calculated from that. It is not necessarily the case that the particles that escape and are measured have the same velocity distribution as the particles that remain in the bulk of the system, but sometimes a good sample is possible.
|
64 |
+
|
65 |
+
Temperature is one of the principal quantities in the study of thermodynamics. Formerly, the magnitude of the kelvin was defined in thermodynamic terms, but nowadays, as mentioned above, it is defined in terms of kinetic theory.
|
66 |
+
|
67 |
+
The thermodynamic temperature is said to be absolute for two reasons. One is that its formal character is independent of the properties of particular materials. The other reason is that its zero is, in a sense, absolute, in that it indicates absence of microscopic classical motion of the constituent particles of matter, so that they have a limiting specific heat of zero for zero temperature, according to the third law of thermodynamics. Nevertheless, a thermodynamic temperature does in fact have a definite numerical value that has been arbitrarily chosen by tradition and is dependent on the property of a particular materials; it is simply less arbitrary than relative "degrees" scales such as Celsius and Fahrenheit. Being an absolute scale with one fixed point (zero), there is only one degree of freedom left to arbitrary choice, rather than two as in relative scales. For the Kelvin scale since May 2019, by international convention, the choice has been made to use knowledge of modes of operation of various thermometric devices, relying on microscopic kinetic theories about molecular motion. The numerical scale is settled by a conventional definition of the value of the Boltzmann constant, which relates macroscopic temperature to average microscopic kinetic energy of particles such as molecules. Its numerical value is arbitrary, and an alternate, less widely used absolute temperature scale exists called the Rankine scale, made to be aligned with the Fahrenheit scale as Kelvin is with Celsius.
|
68 |
+
|
69 |
+
The thermodynamic definition of temperature is due to Kelvin. It is framed in terms of an idealized device called a Carnot engine, imagined to run in a fictive continuous cycle of successive processes that traverse a cycle of states of its working body. The engine takes in a quantity of heat Q1 from a hot reservoir and passes out a lesser quantity of heat Q2 to a cold reservoir. The difference in energy is passed, as thermodynamic work, to a work reservoir, and is considered to be the output of the engine. The cycle is imagined to run so slowly that at each point of the cycle the working body is in a state of thermodynamic equilibrium. The successive processes the cycle are thus imagined to run reversibly with no entropy production. Then the quantity of entropy taken in from the hot reservoir when the working body is heated is equal to that passed to the cold reservoir when the working body is cooled. Then the absolute or thermodynamic temperatures, T1 and T2, of the reservoirs are defined so that to be such that
|
70 |
+
|
71 |
+
|
72 |
+
|
73 |
+
|
74 |
+
|
75 |
+
|
76 |
+
|
77 |
+
|
78 |
+
|
79 |
+
(1)
|
80 |
+
|
81 |
+
The zeroth law of thermodynamics allows this definition to be used to measure the absolute or thermodynamic temperature of an arbitrary body of interest, by making the other heat reservoir have the same temperature as the body of interest.
|
82 |
+
|
83 |
+
Kelvin's original work postulating absolute temperature was published in 1848. It was based on the work of Carnot, before the formulation of the first law of thermodynamics. Carnot had no sound understanding of heat, and no specific concept of entropy. He wrote of 'caloric', and said that all the caloric that passed from the hot reservoir was passed into the cold reservoir. Kelvin wrote in his 1848 paper that his scale was absolute in the sense that it was defined "independently of the properties of any particular kind of matter". His definitive publication, which sets out the definition just stated, was printed in 1853, a paper read in 1851.[21][22][23][24]
|
84 |
+
|
85 |
+
Numerical details were formerly settled by making one of the heat reservoirs a cell at the triple point of water, which was defined to have an absolute temperature of 273.16 K.[25] Nowadays, the numerical value is instead obtained from measurement through the microscopic statistical mechanical international definition, as above.
|
86 |
+
|
87 |
+
In thermodynamic terms, temperature is an intensive variable because it is equal to a differential coefficient of one extensive variable with respect to another, for a given body. It thus has the dimensions of a ratio of two extensive variables. In thermodynamics, two bodies are often considered as connected by contact with a common wall, which has some specific permeability properties. Such specific permeability can be referred to a specific intensive variable. An example is a diathermic wall that is permeable only to heat; the intensive variable for this case is temperature. When the two bodies have been in contact for a very long time, and have settled to a permanent steady state, the relevant intensive variables are equal in the two bodies; for a diathermal wall, this statement is sometimes called the zeroth law of thermodynamics.[26][27][28]
|
88 |
+
|
89 |
+
In particular, when the body is described by stating its internal energy U, an extensive variable, as a function of its entropy S, also an extensive variable, and other state variables V, N, with U = U (S, V, N), then the temperature is equal to the partial derivative of the internal energy with respect to the entropy:[27][28][29]
|
90 |
+
|
91 |
+
|
92 |
+
|
93 |
+
|
94 |
+
|
95 |
+
|
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+
|
97 |
+
|
98 |
+
|
99 |
+
(2)
|
100 |
+
|
101 |
+
Likewise, when the body is described by stating its entropy S as a function of its internal energy U, and other state variables V, N, with S = S (U, V, N), then the reciprocal of the temperature is equal to the partial derivative of the entropy with respect to the internal energy:[27][29][30]
|
102 |
+
|
103 |
+
|
104 |
+
|
105 |
+
|
106 |
+
|
107 |
+
|
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+
|
109 |
+
|
110 |
+
|
111 |
+
(3)
|
112 |
+
|
113 |
+
The above definition, equation (1), of the absolute temperature is due to Kelvin. It refers to systems closed to transfer of matter, and has special emphasis on directly experimental procedures. A presentation of thermodynamics by Gibbs starts at a more abstract level and deals with systems open to the transfer of matter; in this development of thermodynamics, the equations (2) and (3) above are actually alternative definitions of temperature.[31]
|
114 |
+
|
115 |
+
Real world bodies are often not in thermodynamic equilibrium and not homogeneous. For study by methods of classical irreversible thermodynamics, a body is usually spatially and temporally divided conceptually into 'cells' of small size. If classical thermodynamic equilibrium conditions for matter are fulfilled to good approximation in such a 'cell', then it is homogeneous and a temperature exists for it. If this is so for every 'cell' of the body, then local thermodynamic equilibrium is said to prevail throughout the body.[32][33][34][35][36]
|
116 |
+
|
117 |
+
It makes good sense, for example, to say of the extensive variable U, or of the extensive variable S, that it has a density per unit volume, or a quantity per unit mass of the system, but it makes no sense to speak of density of temperature per unit volume or quantity of temperature per unit mass of the system. On the other hand, it makes no sense to speak of the internal energy at a point, while when local thermodynamic equilibrium prevails, it makes good sense to speak of the temperature at a point. Consequently, temperature can vary from point to point in a medium that is not in global thermodynamic equilibrium, but in which there is local thermodynamic equilibrium.
|
118 |
+
|
119 |
+
Thus, when local thermodynamic equilibrium prevails in a body, temperature can be regarded as a spatially varying local property in that body, and this is because temperature is an intensive variable.
|
120 |
+
|
121 |
+
Temperature is a measure of a quality of a state of a material.[37] The quality may be regarded as a more abstract entity than any particular temperature scale that measures it, and is called hotness by some writers.[38] The quality of hotness refers to the state of material only in a particular locality, and in general, apart from bodies held in a steady state of thermodynamic equilibrium, hotness varies from place to place. It is not necessarily the case that a material in a particular place is in a state that is steady and nearly homogeneous enough to allow it to have a well-defined hotness or temperature. Hotness may be represented abstractly as a one-dimensional manifold. Every valid temperature scale has its own one-to-one map into the hotness manifold.[39][40]
|
122 |
+
|
123 |
+
When two systems in thermal contact are at the same temperature no heat transfers between them. When a temperature difference does exist heat flows spontaneously from the warmer system to the colder system until they are in thermal equilibrium. Such heat transfer occurs by conduction or by thermal radiation.[41][42][43][44][45][46][47][48]
|
124 |
+
|
125 |
+
Experimental physicists, for example Galileo and Newton,[49] found that there are indefinitely many empirical temperature scales. Nevertheless, the zeroth law of thermodynamics says that they all measure the same quality. This means that for a body in its own state of internal thermodynamic equilibrium, every correctly calibrated thermometer, of whatever kind, that measures the temperature of the body, records one and the same temperature. For a body that is not in its own state of internal thermodynamic equilibrium, different thermometers can record different temperatures, depending respectively on the mechanisms of operation of the thermometers.
|
126 |
+
|
127 |
+
For experimental physics, hotness means that, when comparing any two given bodies in their respective separate thermodynamic equilibria, any two suitably given empirical thermometers with numerical scale readings will agree as to which is the hotter of the two given bodies, or that they have the same temperature.[50] This does not require the two thermometers to have a linear relation between their numerical scale readings, but it does require that the relation between their numerical readings shall be strictly monotonic.[51][52] A definite sense of greater hotness can be had, independently of calorimetry, of thermodynamics, and of properties of particular materials, from Wien's displacement law of thermal radiation: the temperature of a bath of thermal radiation is proportional, by a universal constant, to the frequency of the maximum of its frequency spectrum; this frequency is always positive, but can have values that tend to zero. Thermal radiation is initially defined for a cavity in thermodynamic equilibrium. These physical facts justify a mathematical statement that hotness exists on an ordered one-dimensional manifold. This is a fundamental character of temperature and thermometers for bodies in their own thermodynamic equilibrium.[7][39][40][53][54]
|
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+
|
129 |
+
Except for a system undergoing a first-order phase change such as the melting of ice, as a closed system receives heat, without change in its volume and without change in external force fields acting on it, its temperature rises. For a system undergoing such a phase change so slowly that departure from thermodynamic equilibrium can be neglected, its temperature remains constant as the system is supplied with latent heat. Conversely, a loss of heat from a closed system, without phase change, without change of volume, and without change in external force fields acting on it, decreases its temperature.[55]
|
130 |
+
|
131 |
+
While for bodies in their own thermodynamic equilibrium states, the notion of temperature requires that all empirical thermometers must agree as to which of two bodies is the hotter or that they are at the same temperature, this requirement is not safe for bodies that are in steady states though not in thermodynamic equilibrium. It can then well be that different empirical thermometers disagree about which is the hotter, and if this is so, then at least one of the bodies does not have a well defined absolute thermodynamic temperature. Nevertheless, any one given body and any one suitable empirical thermometer can still support notions of empirical, non-absolute, hotness and temperature, for a suitable range of processes. This is a matter for study in non-equilibrium thermodynamics.
|
132 |
+
|
133 |
+
When a body is not in a steady state, then the notion of temperature becomes even less safe than for a body in a steady state not in thermodynamic equilibrium. This is also a matter for study in non-equilibrium thermodynamics.
|
134 |
+
|
135 |
+
For axiomatic treatment of thermodynamic equilibrium, since the 1930s, it has become customary to refer to a zeroth law of thermodynamics. The customarily stated minimalist version of such a law postulates only that all bodies, which when thermally connected would be in thermal equilibrium, should be said to have the same temperature by definition, but by itself does not establish temperature as a quantity expressed as a real number on a scale. A more physically informative version of such a law views empirical temperature as a chart on a hotness manifold.[39][54][56] While the zeroth law permits the definitions of many different empirical scales of temperature, the second law of thermodynamics selects the definition of a single preferred, absolute temperature, unique up to an arbitrary scale factor, whence called the thermodynamic temperature.[7][39][57][58][59][60] If internal energy is considered as a function of the volume and entropy of a homogeneous system in thermodynamic equilibrium, thermodynamic absolute temperature appears as the partial derivative of internal energy with respect the entropy at constant volume. Its natural, intrinsic origin or null point is absolute zero at which the entropy of any system is at a minimum. Although this is the lowest absolute temperature described by the model, the third law of thermodynamics postulates that absolute zero cannot be attained by any physical system.
|
136 |
+
|
137 |
+
When an energy transfer to or from a body is only as heat, the state of the body changes. Depending on the surroundings and the walls separating them from the body, various changes are possible in the body. They include chemical reactions, increase of pressure, increase of temperature, and phase change. For each kind of change under specified conditions, the heat capacity is the ratio of the quantity of heat transferred to the magnitude of the change. For example, if the change is an increase in temperature at constant volume, with no phase change and no chemical change, then the temperature of the body rises and its pressure increases. The quantity of heat transferred, ΔQ, divided by the observed temperature change, ΔT, is the body's heat capacity at constant volume:
|
138 |
+
|
139 |
+
If heat capacity is measured for a well defined amount of substance, the specific heat is the measure of the heat required to increase the temperature of such a unit quantity by one unit of temperature. For example, to raise the temperature of water by one kelvin (equal to one degree Celsius) requires 4186 joules per kilogram (J/kg).
|
140 |
+
|
141 |
+
Temperature measurement using modern scientific thermometers and temperature scales goes back at least as far as the early 18th century, when Gabriel Fahrenheit adapted a thermometer (switching to mercury) and a scale both developed by Ole Christensen Rømer. Fahrenheit's scale is still in use in the United States for non-scientific applications.
|
142 |
+
|
143 |
+
Temperature is measured with thermometers that may be calibrated to a variety of temperature scales. In most of the world (except for Belize, Myanmar, Liberia and the United States), the Celsius scale is used for most temperature measuring purposes. Most scientists measure temperature using the Celsius scale and thermodynamic temperature using the Kelvin scale, which is the Celsius scale offset so that its null point is 0 K = −273.15 °C, or absolute zero. Many engineering fields in the US, notably high-tech and US federal specifications (civil and military), also use the Kelvin and Celsius scales. Other engineering fields in the US also rely upon the Rankine scale (a shifted Fahrenheit scale) when working in thermodynamic-related disciplines such as combustion.
|
144 |
+
|
145 |
+
The basic unit of temperature in the International System of Units (SI) is the Kelvin. It has the symbol K.
|
146 |
+
|
147 |
+
For everyday applications, it is often convenient to use the Celsius scale, in which 0 °C corresponds very closely to the freezing point of water and 100 °C is its boiling point at sea level. Because liquid droplets commonly exist in clouds at sub-zero temperatures, 0 °C is better defined as the melting point of ice. In this scale a temperature difference of 1 degree Celsius is the same as a 1kelvin increment, but the scale is offset by the temperature at which ice melts (273.15 K).
|
148 |
+
|
149 |
+
By international agreement,[61] until May 2019, the Kelvin and Celsius scales were defined by two fixing points: absolute zero and the triple point of Vienna Standard Mean Ocean Water, which is water specially prepared with a specified blend of hydrogen and oxygen isotopes. Absolute zero was defined as precisely 0 K and −273.15 °C. It is the temperature at which all classical translational motion of the particles comprising matter ceases and they are at complete rest in the classical model. Quantum-mechanically, however, zero-point motion remains and has an associated energy, the zero-point energy. Matter is in its ground state,[62] and contains no thermal energy. The triple point of water is defined as 273.16 K and 0.01 °C. This definition served the following purposes: it fixed the magnitude of the kelvin as being precisely 1 part in 273.16 parts of the difference between absolute zero and the triple point of water; it establishes that one kelvin has precisely the same magnitude as one degree on the Celsius scale; and it established the difference between the null points of these scales as being 273.15 K (0 K = −273.15 °C and 273.16 K = 0.01 °C). Since 2019, there has been a new definition based on the Boltzmann constant,[63] but the scales are scarcely changed.
|
150 |
+
|
151 |
+
In the United States, the Fahrenheit scale is widely used. On this scale the freezing point of water corresponds to 32 °F and the boiling point to 212 °F. The Rankine scale, still used in fields of chemical engineering in the US, is an absolute scale based on the Fahrenheit increment.
|
152 |
+
|
153 |
+
The following table shows the temperature conversion formulas for conversions to and from the Celsius scale.
|
154 |
+
|
155 |
+
The field of plasma physics deals with phenomena of electromagnetic nature that involve very high temperatures. It is customary to express temperature as energy in units of electronvolts (eV) or kiloelectronvolts (keV). The energy, which has a different dimension from temperature, is then calculated as the product of the Boltzmann constant and temperature,
|
156 |
+
|
157 |
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|
158 |
+
|
159 |
+
E
|
160 |
+
=
|
161 |
+
|
162 |
+
k
|
163 |
+
|
164 |
+
B
|
165 |
+
|
166 |
+
|
167 |
+
T
|
168 |
+
|
169 |
+
|
170 |
+
{\displaystyle E=k_{\text{B}}T}
|
171 |
+
|
172 |
+
. Then, 1 eV corresponds to 11605 K. In the study of QCD matter one routinely encounters temperatures of the order of a few hundred MeV, equivalent to about 1012 K.
|
173 |
+
|
174 |
+
Historically, there are several scientific approaches to the explanation of temperature: the classical thermodynamic description based on macroscopic empirical variables that can be measured in a laboratory; the kinetic theory of gases which relates the macroscopic description to the probability distribution of the energy of motion of gas particles; and a microscopic explanation based on statistical physics and quantum mechanics. In addition, rigorous and purely mathematical treatments have provided an axiomatic approach to classical thermodynamics and temperature.[64] Statistical physics provides a deeper understanding by describing the atomic behavior of matter, and derives macroscopic properties from statistical averages of microscopic states, including both classical and quantum states. In the fundamental physical description, using natural units, temperature may be measured directly in units of energy. However, in the practical systems of measurement for science, technology, and commerce, such as the modern metric system of units, the macroscopic and the microscopic descriptions are interrelated by the Boltzmann constant, a proportionality factor that scales temperature to the microscopic mean kinetic energy.
|
175 |
+
|
176 |
+
The microscopic description in statistical mechanics is based on a model that analyzes a system into its fundamental particles of matter or into a set of classical or quantum-mechanical oscillators and considers the system as a statistical ensemble of microstates. As a collection of classical material particles, temperature is a measure of the mean energy of motion, called kinetic energy, of the particles, whether in solids, liquids, gases, or plasmas. The kinetic energy, a concept of classical mechanics, is half the mass of a particle times its speed squared. In this mechanical interpretation of thermal motion, the kinetic energies of material particles may reside in the velocity of the particles of their translational or vibrational motion or in the inertia of their rotational modes. In monatomic perfect gases and, approximately, in most gases, temperature is a measure of the mean particle kinetic energy. It also determines the probability distribution function of the energy. In condensed matter, and particularly in solids, this purely mechanical description is often less useful and the oscillator model provides a better description to account for quantum mechanical phenomena. Temperature determines the statistical occupation of the microstates of the ensemble. The microscopic definition of temperature is only meaningful in the thermodynamic limit, meaning for large ensembles of states or particles, to fulfill the requirements of the statistical model.
|
177 |
+
|
178 |
+
The kinetic energy is also considered as a component of thermal energy. The thermal energy may be partitioned into independent components attributed to the degrees of freedom of the particles or to the modes of oscillators in a thermodynamic system. In general, the number of these degrees of freedom that are available for the equipartitioning of energy depends on the temperature, i.e. the energy region of the interactions under consideration. For solids, the thermal energy is associated primarily with the vibrations of its atoms or molecules about their equilibrium position. In an ideal monatomic gas, the kinetic energy is found exclusively in the purely translational motions of the particles. In other systems, vibrational and rotational motions also contribute degrees of freedom.
|
179 |
+
|
180 |
+
Maxwell and Boltzmann developed a kinetic theory that yields a fundamental understanding of temperature in gases.[65]
|
181 |
+
This theory also explains the ideal gas law and the observed heat capacity of monatomic (or 'noble') gases.[66][67][68]
|
182 |
+
|
183 |
+
The ideal gas law is based on observed empirical relationships between pressure (p), volume (V), and temperature (T), and was recognized long before the kinetic theory of gases was developed (see Boyle's and Charles's laws). The ideal gas law states:[69]
|
184 |
+
|
185 |
+
where n is the number of moles of gas and R = 8.314462618... J⋅mol−1⋅K−1[70] is the gas constant.
|
186 |
+
|
187 |
+
This relationship gives us our first hint that there is an absolute zero on the temperature scale, because it only holds if the temperature is measured on an absolute scale such as Kelvin's. The ideal gas law allows one to measure temperature on this absolute scale using the gas thermometer. The temperature in kelvins can be defined as the pressure in pascals of one mole of gas in a container of one cubic meter, divided by the gas constant.
|
188 |
+
|
189 |
+
Although it is not a particularly convenient device, the gas thermometer provides an essential theoretical basis by which all thermometers can be calibrated. As a practical matter, it is not possible to use a gas thermometer to measure absolute zero temperature since the gases tend to condense into a liquid long before the temperature reaches zero. It is possible, however, to extrapolate to absolute zero by using the ideal gas law, as shown in the figure.
|
190 |
+
|
191 |
+
The kinetic theory assumes that pressure is caused by the force associated with individual atoms striking the walls, and that all energy is translational kinetic energy. Using a sophisticated symmetry argument,[71] Boltzmann deduced what is now called the Maxwell–Boltzmann probability distribution function for the velocity of particles in an ideal gas. From that probability distribution function, the average kinetic energy (per particle) of a monatomic ideal gas is[67][72]
|
192 |
+
|
193 |
+
where the Boltzmann constant kB is the ideal gas constant divided by the Avogadro number, and
|
194 |
+
|
195 |
+
|
196 |
+
|
197 |
+
|
198 |
+
v
|
199 |
+
|
200 |
+
rms
|
201 |
+
|
202 |
+
|
203 |
+
=
|
204 |
+
|
205 |
+
|
206 |
+
|
207 |
+
⟨
|
208 |
+
|
209 |
+
v
|
210 |
+
|
211 |
+
2
|
212 |
+
|
213 |
+
|
214 |
+
⟩
|
215 |
+
|
216 |
+
|
217 |
+
|
218 |
+
|
219 |
+
|
220 |
+
{\displaystyle v_{\text{rms}}={\sqrt {\left\langle v^{2}\right\rangle }}}
|
221 |
+
|
222 |
+
is the root-mean-square speed. Thus the ideal gas law states that internal energy is directly proportional to temperature.[73] This direct proportionality between temperature and internal energy is a special case of the equipartition theorem, and holds only in the classical limit of an ideal gas. It does not hold for most substances, although it is true that temperature is a monotonic (non-decreasing) function of internal energy.
|
223 |
+
|
224 |
+
When two otherwise isolated bodies are connected together by a rigid physical path impermeable to matter, there is spontaneous transfer of energy as heat from the hotter to the colder of them. Eventually, they reach a state of mutual thermal equilibrium, in which heat transfer has ceased, and the bodies' respective state variables have settled to become unchanging.
|
225 |
+
|
226 |
+
One statement of the zeroth law of thermodynamics is that if two systems are each in thermal equilibrium with a third system, then they are also in thermal equilibrium with each other.
|
227 |
+
|
228 |
+
This statement helps to define temperature but it does not, by itself, complete the definition. An empirical temperature is a numerical scale for the hotness of a thermodynamic system. Such hotness may be defined as existing on a one-dimensional manifold, stretching between hot and cold. Sometimes the zeroth law is stated to include the existence of a unique universal hotness manifold, and of numerical scales on it, so as to provide a complete definition of empirical temperature.[56] To be suitable for empirical thermometry, a material must have a monotonic relation between hotness and some easily measured state variable, such as pressure or volume, when all other relevant coordinates are fixed. An exceptionally suitable system is the ideal gas, which can provide a temperature scale that matches the absolute Kelvin scale. The Kelvin scale is defined on the basis of the second law of thermodynamics.
|
229 |
+
|
230 |
+
As an alternative to considering or defining the zeroth law of thermodynamics, it was the historical development in thermodynamics to define temperature in terms of the second law of thermodynamics which deals with entropy. The second law states that any process will result in either no change or a net increase in the entropy of the universe. This can be understood in terms of probability.
|
231 |
+
|
232 |
+
For example, in a series of coin tosses, a perfectly ordered system would be one in which either every toss comes up heads or every toss comes up tails. This means the outcome is always 100% the same result. In contrast, many mixed (disordered) outcomes are possible, and their number increases with each toss. Eventually, the combinations of ~50% heads and ~50% tails dominate and obtaining an outcome significantly different from 50/50 becomes increasingly unlikely. Thus the system naturally progresses to a state of maximum disorder or entropy.
|
233 |
+
|
234 |
+
As temperature governs the transfer of heat between two systems and the universe tends to progress toward a maximum of entropy, it is expected that there is some relationship between temperature and entropy. A heat engine is a device for converting thermal energy into mechanical energy, resulting in the performance of work. and analysis of the Carnot heat engine provides the necessary relationships. The work from a heat engine corresponds to the difference between the heat put into the system at high temperature, qH and the heat extracted at the low temperature, qC. The efficiency is the work divided by the heat input:
|
235 |
+
|
236 |
+
|
237 |
+
|
238 |
+
|
239 |
+
|
240 |
+
|
241 |
+
|
242 |
+
|
243 |
+
|
244 |
+
(4)
|
245 |
+
|
246 |
+
where wcy is the work done per cycle. The efficiency depends only on qC/qH. Because qC and qH correspond to heat transfer at the temperatures TC and TH respectively, qC/qH should be some function of these temperatures:
|
247 |
+
|
248 |
+
|
249 |
+
|
250 |
+
|
251 |
+
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
|
256 |
+
(5)
|
257 |
+
|
258 |
+
Carnot's theorem states that all reversible engines operating between the same heat reservoirs are equally efficient. Thus, a heat engine operating between T1 and T3 must have the same efficiency as one consisting of two cycles, one between T1 and T2, and the second between T2 and T3. This can only be the case if
|
259 |
+
|
260 |
+
which implies
|
261 |
+
|
262 |
+
Since the first function is independent of T2, this temperature must cancel on the right side, meaning f(T1, T3) is of the form g(T1)/g(T3) (i.e. f(T1, T3) = f(T1, T2)f(T2, T3) = g(T1)/g(T2) · g(T2)/g(T3) = g(T1)/g(T3)), where g is a function of a single temperature. A temperature scale can now be chosen with the property that
|
263 |
+
|
264 |
+
|
265 |
+
|
266 |
+
|
267 |
+
|
268 |
+
|
269 |
+
|
270 |
+
|
271 |
+
|
272 |
+
(6)
|
273 |
+
|
274 |
+
Substituting (6) back into (4) gives a relationship for the efficiency in terms of temperature:
|
275 |
+
|
276 |
+
|
277 |
+
|
278 |
+
|
279 |
+
|
280 |
+
|
281 |
+
|
282 |
+
|
283 |
+
|
284 |
+
(7)
|
285 |
+
|
286 |
+
For TC = 0 K the efficiency is 100% and that efficiency becomes greater than 100% below 0 K. Since an efficiency greater than 100% violates the first law of thermodynamics, this implies that 0 K is the minimum possible temperature. In fact the lowest temperature ever obtained in a macroscopic system was 20 nK, which was achieved in 1995 at NIST. Subtracting the right hand side of (5) from the middle portion and rearranging gives
|
287 |
+
|
288 |
+
where the negative sign indicates heat ejected from the system. This relationship suggests the existence of a state function, S, defined by
|
289 |
+
|
290 |
+
|
291 |
+
|
292 |
+
|
293 |
+
|
294 |
+
|
295 |
+
|
296 |
+
|
297 |
+
|
298 |
+
(8)
|
299 |
+
|
300 |
+
where the subscript indicates a reversible process. The change of this state function around any cycle is zero, as is necessary for any state function. This function corresponds to the entropy of the system, which was described previously. Rearranging (8) gives a formula for temperature in terms of fictive infinitesimal quasi-reversible elements of entropy and heat:
|
301 |
+
|
302 |
+
|
303 |
+
|
304 |
+
|
305 |
+
|
306 |
+
|
307 |
+
|
308 |
+
|
309 |
+
|
310 |
+
(9)
|
311 |
+
|
312 |
+
For a system, where entropy S(E) is a function of its energy E, the temperature T is given by
|
313 |
+
|
314 |
+
|
315 |
+
|
316 |
+
|
317 |
+
|
318 |
+
|
319 |
+
|
320 |
+
|
321 |
+
|
322 |
+
(10)
|
323 |
+
|
324 |
+
i.e. the reciprocal of the temperature is the rate of increase of entropy with respect to energy.
|
325 |
+
|
326 |
+
Statistical mechanics defines temperature based on a system's fundamental degrees of freedom. Eq.(10) is the defining relation of temperature, where the entropy
|
327 |
+
|
328 |
+
|
329 |
+
|
330 |
+
S
|
331 |
+
|
332 |
+
|
333 |
+
{\displaystyle S}
|
334 |
+
|
335 |
+
is defined (up to a constant) by the logarithm of the number of microstates of the system in the given macrostate (as specified in the microcanonical ensemble):
|
336 |
+
|
337 |
+
where
|
338 |
+
|
339 |
+
|
340 |
+
|
341 |
+
|
342 |
+
k
|
343 |
+
|
344 |
+
|
345 |
+
B
|
346 |
+
|
347 |
+
|
348 |
+
|
349 |
+
|
350 |
+
|
351 |
+
{\displaystyle k_{\mathrm {B} }}
|
352 |
+
|
353 |
+
is Boltzmann's constant and N is the number of microstates.
|
354 |
+
|
355 |
+
When two systems with different temperatures are put into purely thermal connection, heat will flow from the higher temperature system to the lower temperature one; thermodynamically this is understood by the second law of thermodynamics: The total change in entropy following a transfer of energy
|
356 |
+
|
357 |
+
|
358 |
+
|
359 |
+
Δ
|
360 |
+
E
|
361 |
+
|
362 |
+
|
363 |
+
{\displaystyle \Delta E}
|
364 |
+
|
365 |
+
from system 1 to system 2 is:
|
366 |
+
|
367 |
+
and is thus positive if
|
368 |
+
|
369 |
+
|
370 |
+
|
371 |
+
|
372 |
+
T
|
373 |
+
|
374 |
+
1
|
375 |
+
|
376 |
+
|
377 |
+
>
|
378 |
+
|
379 |
+
T
|
380 |
+
|
381 |
+
2
|
382 |
+
|
383 |
+
|
384 |
+
|
385 |
+
|
386 |
+
{\displaystyle T_{1}>T_{2}}
|
387 |
+
|
388 |
+
From the point of view of statistical mechanics, the total number of micsostates in the combined system 1 + system 2 is
|
389 |
+
|
390 |
+
|
391 |
+
|
392 |
+
|
393 |
+
N
|
394 |
+
|
395 |
+
1
|
396 |
+
|
397 |
+
|
398 |
+
⋅
|
399 |
+
|
400 |
+
N
|
401 |
+
|
402 |
+
2
|
403 |
+
|
404 |
+
|
405 |
+
|
406 |
+
|
407 |
+
{\displaystyle N_{1}\cdot N_{2}}
|
408 |
+
|
409 |
+
, the logarithm of which (times Boltzmann's constant) is the sum of their entropies; thus a flow of heat from high to low temperature, which brings an increase in total entropy, is more likely than any other scenario (normally it is much more likely), as there are more microstates in the resulting macrostate.
|
410 |
+
|
411 |
+
It is possible to extend the definition of temperature even to systems of few particles, like in a quantum dot. The generalized temperature is obtained by considering time ensembles instead of configuration-space ensembles given in statistical mechanics in the case of thermal and particle exchange between a small system of fermions (N even less than 10) with a single/double-occupancy system. The finite quantum grand canonical ensemble,[74] obtained under the hypothesis of ergodicity and orthodicity,[75] allows expressing the generalized temperature from the ratio of the average time of occupation
|
412 |
+
|
413 |
+
|
414 |
+
|
415 |
+
|
416 |
+
τ
|
417 |
+
|
418 |
+
1
|
419 |
+
|
420 |
+
|
421 |
+
|
422 |
+
|
423 |
+
{\displaystyle \tau _{1}}
|
424 |
+
|
425 |
+
and
|
426 |
+
|
427 |
+
|
428 |
+
|
429 |
+
|
430 |
+
τ
|
431 |
+
|
432 |
+
2
|
433 |
+
|
434 |
+
|
435 |
+
|
436 |
+
|
437 |
+
{\displaystyle \tau _{2}}
|
438 |
+
|
439 |
+
of the single/double-occupancy system:[76]
|
440 |
+
|
441 |
+
where EF is the Fermi energy. This generalized temperature tends to the ordinary temperature when N goes to infinity.
|
442 |
+
|
443 |
+
On the empirical temperature scales that are not referenced to absolute zero, a negative temperature is one below the zero-point of the scale used. For example, dry ice has a sublimation temperature of −78.5 °C which is equivalent to −109.3 °F. On the absolute kelvin scale this temperature is 194.6 K. No body can be brought to exactly 0 K (the temperature of the ideally coldest possible body) by any finite practicable process; this is a consequence of the third law of thermodynamics.
|
444 |
+
|
445 |
+
The international kinetic theory temperature of a body cannot take negative values. The thermodynamic temperature scale, however, is not so constrained.
|
446 |
+
|
447 |
+
For a body of matter, there can sometimes be conceptually defined, in terms of microscopic degrees of freedom, namely particle spins, a subsystem, with a temperature other than that of the whole body. When the body is in its own state of internal thermodynamic equilibrium, the temperatures of the whole body and of the subsystem must be the same. The two temperatures can differ when, by work through externally imposed force fields, energy can be transferred to and from the subsystem, separately from the rest of the body; then the whole body is not in its own state of internal thermodynamic equilibrium. There is an upper limit of energy such a spin subsystem can attain.
|
448 |
+
|
449 |
+
Considering the subsystem to be in a temporary state of virtual thermodynamic equilibrium, it is possible to obtain a negative temperature on the thermodynamic scale. Thermodynamic temperature is the inverse of the derivative of the subsystem's entropy with respect to its internal energy. As the subsystem's internal energy increases, the entropy increases for some range, but eventually attains a maximum value and then begins to decrease as the highest energy states begin to fill. At the point of maximum entropy, the temperature function shows the behavior of a singularity, because the slope of the entropy function decreases to zero and then turns negative. As the subsystem's entropy reaches its maximum, its thermodynamic temperature goes to positive infinity, switching to negative infinity as the slope turns negative. Such negative temperatures are hotter than any positive temperature. Over time, when the subsystem is exposed to the rest of the body, which has a positive temperature, energy is transferred as heat from the negative temperature subsystem to the positive temperature system.[77] The kinetic theory temperature is not defined for such subsystems.
|
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|
1 |
+
|
2 |
+
|
3 |
+
Time is the indefinite continued progress of existence and events that occur in an apparently irreversible succession from the past, through the present, into the future.[1][2][3] It is a component quantity of various measurements used to sequence events, to compare the duration of events or the intervals between them, and to quantify rates of change of quantities in material reality or in the conscious experience.[4][5][6][7] Time is often referred to as a fourth dimension, along with three spatial dimensions.[8]
|
4 |
+
|
5 |
+
Time has long been an important subject of study in religion, philosophy, and science, but defining it in a manner applicable to all fields without circularity has consistently eluded scholars.[2][6][7][9][10][11]
|
6 |
+
Nevertheless, diverse fields such as business, industry, sports, the sciences, and the performing arts all incorporate some notion of time into their respective measuring systems.[12][13][14]
|
7 |
+
|
8 |
+
Time in physics is operationally defined as "what a clock reads".[6][15][16] Time is one of the seven fundamental physical quantities in both the International System of Units (SI) and International System of Quantities. The SI base unit of time is the second. Time is used to define other quantities – such as velocity – so defining time in terms of such quantities would result in circularity of definition.[17] An operational definition of time, wherein one says that observing a certain number of repetitions of one or another standard cyclical event (such as the passage of a free-swinging pendulum) constitutes one standard unit such as the second, is highly useful in the conduct of both advanced experiments and everyday affairs of life. To describe observations of an event, a location (position in space) and time are typically noted.
|
9 |
+
|
10 |
+
The operational definition of time does not address what the fundamental nature of it is. It does not address why events can happen forward and backwards in space, whereas events only happen in the forward progress of time. Investigations into the relationship between space and time led physicists to define the spacetime continuum. General Relativity is the primary framework for understanding how spacetime works. Through advances in both theoretical and experimental investigations of space-time, it has been shown that time can be distorted, particularly at the edges of black holes.
|
11 |
+
|
12 |
+
Temporal measurement has occupied scientists and technologists, and was a prime motivation in navigation and astronomy. Periodic events and periodic motion have long served as standards for units of time. Examples include the apparent motion of the sun across the sky, the phases of the moon, the swing of a pendulum, and the beat of a heart. Currently, the international unit of time, the second, is defined by measuring the electronic transition frequency of caesium atoms (see below). Time is also of significant social importance, having economic value ("time is money") as well as personal value, due to an awareness of the limited time in each day and in human life spans.
|
13 |
+
|
14 |
+
Generally speaking, methods of temporal measurement, or chronometry, take two distinct forms: the calendar, a mathematical tool for organising intervals of time,[18]
|
15 |
+
and the clock, a physical mechanism that counts the passage of time. In day-to-day life, the clock is consulted for periods less than a day whereas the calendar is consulted for periods longer than a day. Increasingly, personal electronic devices display both calendars and clocks simultaneously. The number (as on a clock dial or calendar) that marks the occurrence of a specified event as to hour or date is obtained by counting from a fiducial epoch – a central reference point.
|
16 |
+
|
17 |
+
Artifacts from the Paleolithic suggest that the moon was used to reckon time as early as 6,000 years ago.[19]
|
18 |
+
Lunar calendars were among the first to appear, with years of either 12 or 13 lunar months (either 354 or 384 days). Without intercalation to add days or months to some years, seasons quickly drift in a calendar based solely on twelve lunar months. Lunisolar calendars have a thirteenth month added to some years to make up for the difference between a full year (now known to be about 365.24 days) and a year of just twelve lunar months. The numbers twelve and thirteen came to feature prominently in many cultures, at least partly due to this relationship of months to years. Other early forms of calendars originated in Mesoamerica, particularly in ancient Mayan civilization. These calendars were religiously and astronomically based, with 18 months in a year and 20 days in a month, plus five epagomenal days at the end of the year.[20]
|
19 |
+
|
20 |
+
The reforms of Julius Caesar in 45 BC put the Roman world on a solar calendar. This Julian calendar was faulty in that its intercalation still allowed the astronomical solstices and equinoxes to advance against it by about 11 minutes per year. Pope Gregory XIII introduced a correction in 1582; the Gregorian calendar was only slowly adopted by different nations over a period of centuries, but it is now by far the most commonly used calendar around the world.
|
21 |
+
|
22 |
+
During the French Revolution, a new clock and calendar were invented in an attempt to de-Christianize time and create a more rational system in order to replace the Gregorian calendar. The French Republican Calendar's days consisted of ten hours of a hundred minutes of a hundred seconds, which marked a deviation from the base 12 (duodecimal) system used in many other devices by many cultures. The system was abolished in 1806.[21]
|
23 |
+
|
24 |
+
A large variety of devices have been invented to measure time. The study of these devices is called horology.[22]
|
25 |
+
|
26 |
+
An Egyptian device that dates to c. 1500 BC, similar in shape to a bent T-square, measured the passage of time from the shadow cast by its crossbar on a nonlinear rule. The T was oriented eastward in the mornings. At noon, the device was turned around so that it could cast its shadow in the evening direction.[23]
|
27 |
+
|
28 |
+
A sundial uses a gnomon to cast a shadow on a set of markings calibrated to the hour. The position of the shadow marks the hour in local time. The idea to separate the day into smaller parts is credited to Egyptians because of their sundials, which operated on a duodecimal system. The importance of the number 12 is due to the number of lunar cycles in a year and the number of stars used to count the passage of night.[24]
|
29 |
+
|
30 |
+
The most precise timekeeping device of the ancient world was the water clock, or clepsydra, one of which was found in the tomb of Egyptian pharaoh Amenhotep I. They could be used to measure the hours even at night, but required manual upkeep to replenish the flow of water. The ancient Greeks and the people from Chaldea (southeastern Mesopotamia) regularly maintained timekeeping records as an essential part of their astronomical observations. Arab inventors and engineers in particular made improvements on the use of water clocks up to the Middle Ages.[25] In the 11th century, Chinese inventors and engineers invented the first mechanical clocks driven by an escapement mechanism.
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The hourglass uses the flow of sand to measure the flow of time. They were used in navigation. Ferdinand Magellan used 18 glasses on each ship for his circumnavigation of the globe (1522).[26]
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Incense sticks and candles were, and are, commonly used to measure time in temples and churches across the globe. Waterclocks, and later, mechanical clocks, were used to mark the events of the abbeys and monasteries of the Middle Ages. Richard of Wallingford (1292–1336), abbot of St. Alban's abbey, famously built a mechanical clock as an astronomical orrery about 1330.[27][28]
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Great advances in accurate time-keeping were made by Galileo Galilei and especially Christiaan Huygens with the invention of pendulum driven clocks along with the invention of the minute hand by Jost Burgi.[29]
|
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+
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The English word clock probably comes from the Middle Dutch word klocke which, in turn, derives from the medieval Latin word clocca, which ultimately derives from Celtic and is cognate with French, Latin, and German words that mean bell. The passage of the hours at sea were marked by bells, and denoted the time (see ship's bell). The hours were marked by bells in abbeys as well as at sea.
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Clocks can range from watches, to more exotic varieties such as the Clock of the Long Now. They can be driven by a variety of means, including gravity, springs, and various forms of electrical power, and regulated by a variety of means such as a pendulum.
|
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Alarm clocks first appeared in ancient Greece around 250 BC with a water clock that would set off a whistle. This idea was later mechanized by Levi Hutchins and Seth E. Thomas.[29]
|
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A chronometer is a portable timekeeper that meets certain precision standards. Initially, the term was used to refer to the marine chronometer, a timepiece used to determine longitude by means of celestial navigation, a precision firstly achieved by John Harrison. More recently, the term has also been applied to the chronometer watch, a watch that meets precision standards set by the Swiss agency COSC.
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The most accurate timekeeping devices are atomic clocks, which are accurate to seconds in many millions of years,[31] and are used to calibrate other clocks and timekeeping instruments.
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Atomic clocks use the frequency of electronic transitions in certain atoms to measure the second. One of the atoms used is caesium, most modern atomic clocks probe caesium with microwaves to determine the frequency of these electron vibrations.[32] Since 1967, the International System of Measurements bases its unit of time, the second, on the properties of caesium atoms. SI defines the second as 9,192,631,770 cycles of the radiation that corresponds to the transition between two electron spin energy levels of the ground state of the 133Cs atom.
|
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Today, the Global Positioning System in coordination with the Network Time Protocol can be used to synchronize timekeeping systems across the globe.
|
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In medieval philosophical writings, the atom was a unit of time referred to as the smallest possible division of time. The earliest known occurrence in English is in Byrhtferth's Enchiridion (a science text) of 1010–1012,[33] where it was defined as 1/564 of a momentum (1½ minutes),[34] and thus equal to 15/94 of a second. It was used in the computus, the process of calculating the date of Easter.
|
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As of May 2010[update], the smallest time interval uncertainty in direct measurements is on the order of 12 attoseconds (1.2 × 10−17 seconds), about 3.7 × 1026 Planck times.[35]
|
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+
|
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The second (s) is the SI base unit. A minute (min) is 60 seconds in length, and an hour is 60 minutes or 3600 seconds in length. A day is usually 24 hours or 86,400 seconds in length; however, the duration of a calendar day can vary due to Daylight saving time and Leap seconds.
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|
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The Mean Solar Time system defines the second as 1/86,400 of the mean solar day, which is the year-average of the solar day. The solar day is the time interval between two successive solar noons, i.e., the time interval between two successive passages of the Sun across the local meridian. The local meridian is an imaginary line that runs from celestial north pole to celestial south pole passing directly over the head of the observer. At the local meridian the Sun reaches its highest point on its daily arc across the sky.
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|
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In 1874 the British Association for the Advancement of Science introduced the CGS (centimetre/gramme/second system) combining fundamental units of length, mass and time. The second is "elastic", because tidal friction is slowing the earth's rotation rate. For use in calculating ephemerides of celestial motion, therefore, in 1952 astronomers introduced the "ephemeris second", currently defined as
|
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the fraction 1/31,556,925.9747 of the tropical year for 1900 January 0 at 12 hours ephemeris time.[36]
|
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+
|
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The CGS system has been superseded by the Système international. The SI base unit for time is the SI second. The International System of Quantities, which incorporates the SI, also defines larger units of time equal to fixed integer multiples of one second (1 s), such as the minute, hour and day. These are not part of the SI, but may be used alongside the SI. Other units of time such as the month and the year are not equal to fixed multiples of 1 s, and instead exhibit significant variations in duration.[37]
|
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|
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The official SI definition of the second is as follows:[37][38]
|
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+
|
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The second is the duration of 9,192,631,770 periods of the radiation corresponding to the transition between the two hyperfine levels of the ground state of the caesium 133 atom.
|
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+
|
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At its 1997 meeting, the CIPM affirmed that this definition refers to a caesium atom in its ground state at a temperature of 0 K.[37]
|
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+
|
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The current definition of the second, coupled with the current definition of the meter, is based on the special theory of relativity, which affirms our spacetime to be a Minkowski space. The definition of the second in mean solar time, however, is unchanged.
|
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+
|
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While in theory, the concept of a single worldwide universal time-scale may have been conceived of many centuries ago, in practicality the technical ability to create and maintain such a time-scale did not become possible until the mid-19th century. The timescale adopted was Greenwich Mean Time, created in 1847. A few countries have replaced it with Coordinated Universal Time, UTC.
|
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|
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With the advent of the industrial revolution, a greater understanding and agreement on the nature of time itself became increasingly necessary and helpful. In 1847 in Britain, Greenwich Mean Time (GMT) was first created for use by the British railways, the British navy, and the British shipping industry. Using telescopes, GMT was calibrated to the mean solar time at the Royal Observatory, Greenwich in the UK.
|
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|
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As international commerce continued to increase throughout Europe, in order to achieve a more efficiently functioning modern society, an agreed upon, and highly accurate international standard of time measurement became necessary. In order to find or determine such a time-standard, three steps had to be followed:
|
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|
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The development of what is now known as UTC time began as a collaboration between 41 nations, officially agreed and signed at the International Meridian Conference, in Washington D.C. in 1884. At this conference, the local mean solar time at the Royal Observatory, Greenwich in England was chosen to define the "universal day", counted from 0 hours at Greenwich mean midnight. This agreed with the civil Greenwich Mean Time used on the island of Great Britain since 1847. In contrast astronomical GMT began at mean noon, i.e. astronomical day X began at noon of civil day X. The purpose of this was to keep one night's observations under one date. The civil system was adopted as of 0 hours (civil) 1 January 1925. Nautical GMT began 24 hours before astronomical GMT, at least until 1805 in the Royal Navy, but persisted much later elsewhere because it was mentioned at the 1884 conference. In 1884, the Greenwich meridian was used for two-thirds of all charts and maps as their Prime Meridian.[39]
|
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|
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Among the 41 nations represented at the conference, the advanced time-technologies that had already come into use in Britain were fundamental components of the agreed method of arriving at a universal and agreed international time. In 1928 Greenwich Mean Time was rebranded for scientific purposes by the International Astronomical Union as Universal Time (UT). This was to avoid confusion with the previous system in which the day had begun at noon. As the general public had always begun the day at midnight, the timescale continued to be presented to them as Greenwich Mean Time. By 1956, universal time had been split into various versions: UT2, which smoothed for polar motion and seasonal effects, was presented to the public as Greenwich Mean Time. Later, UT1 (which smooths only for polar motion) became the default form of UT used by astronomers and hence the form used in navigation, sunrise and sunset and moonrise and moonset tables where the name Greenwich Mean Time continues to be employed. Greenwich Mean Time is also the preferred method of describing the timescale used by legislators. Even to the present day, UT is still based on an international telescopic system. Observations at the Greenwich Observatory itself ceased in 1954, though the location is still used as the basis for the coordinate system. Because the rotational period of Earth is not perfectly constant, the duration of a second would vary if calibrated to a telescope-based standard like GMT, where the second is defined as 1/86 400 of the mean solar day.
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|
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Until 1960, the methods and definitions of time-keeping that had been laid out at the International Meridian Conference proved to be adequate to meet time tracking needs of science. Still, with the advent of the "electronic revolution" in the latter half of the 20th century, the technologies that had been available at the time of the Convention of the Metre proved to be in need of further refinement in order to meet the needs of the ever-increasing precision that the "electronic revolution" had begun to require.
|
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|
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An invariable second (the "ephemeris second") had been defined, use of which removed the errors in ephemerides resulting from the use of the variable mean solar second as the time argument. In 1960 this ephemeris second was made the basis of the "coordinated universal time" which was being derived from atomic clocks. It is a specified fraction of the mean tropical year as at 1900 and, being based on historical telescope observations, corresponds roughly to the mean solar second of the early nineteenth century.[40]
|
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|
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In 1967 a further step was taken with the introduction of the SI second, essentially the ephemeris second as measured by atomic clocks and formally defined in atomic terms.[41]
|
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The SI second (Standard Internationale second) is based directly on the measurement of the atomic-clock observation of the frequency oscillation of caesium atoms. It is the basis of all atomic timescales, e.g. coordinated universal time, GPS time, International Atomic Time, etc. Atomic clocks do not measure nuclear decay rates, which is a common misconception, but rather measure a certain natural vibrational frequency of caesium-133.[42] Coordinated universal time is subject to one constraint which does not affect the other atomic timescales. As it has been adopted as the civil timescale by some countries (most countries have opted to retain mean solar time) it is not permitted to deviate from GMT by more than 0.9 second. This is achieved by the occasional insertion of a leap second.
|
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|
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Most countries use mean solar time. Australia, Canada (Quebec only), Colombia, France, Germany, New Zealand, Papua New Guinea (Bougainville only), Paraguay, Portugal, Switzerland, the United States and Venezuela use UTC. However, UTC is widely used by the scientific community in countries where mean solar time is official. UTC time is based on the SI second, which was first defined in 1967, and is based on the use of atomic clocks. Some other less used but closely related time-standards include International Atomic Time (TAI), Terrestrial Time, and Barycentric Dynamical Time.
|
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+
|
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Between 1967 and 1971, UTC was periodically adjusted by fractional amounts of a second in order to adjust and refine for variations in mean solar time, with which it is aligned. After 1 January 1972, UTC time has been defined as being offset from atomic time by a whole number of seconds, changing only when a leap second is added to keep radio-controlled clocks synchronized with the rotation of the Earth.
|
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|
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The Global Positioning System also broadcasts a very precise time signal worldwide, along with instructions for converting GPS time to UTC. GPS-time is based on, and regularly synchronized with or from, UTC-time.
|
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+
|
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The surface of the Earth is split up into a number of time zones. Most time zones are exactly one hour apart, and by convention compute their local time as an offset from GMT. For example, time zones at sea are based on GMT. In many locations (but not at sea) these offsets vary twice yearly due to daylight saving time transitions.
|
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|
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These conversions are accurate at the millisecond level for time systems based on the rotation of the Earth (UT1 and TT). Conversions between atomic time systems (TAI, GPS, and UTC) are accurate at the microsecond level.
|
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+
|
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Definitions:
|
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|
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Unlike solar time, which is relative to the apparent position of the Sun, sidereal time is the measurement of time relative to that of a distant star. In astronomy, sidereal time is used to predict when a star will reach its highest point in the sky. Due to Earth's orbital motion around the Sun, a mean solar day is about 3 minutes 56 seconds longer than a mean sidereal day, or 1⁄366 more than a mean sidereal day.
|
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|
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Another form of time measurement consists of studying the past. Events in the past can be ordered in a sequence (creating a chronology), and can be put into chronological groups (periodization). One of the most important systems of periodization is the geologic time scale, which is a system of periodizing the events that shaped the Earth and its life. Chronology, periodization, and interpretation of the past are together known as the study of history.
|
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|
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The term "time" is generally used for many close but different concepts, including:
|
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|
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Ancient cultures such as Incan, Mayan, Hopi, and other Native American Tribes – plus the Babylonians, ancient Greeks, Hinduism, Buddhism, Jainism, and others – have a concept of a wheel of time: they regard time as cyclical and quantic,[clarification needed] consisting of repeating ages that happen to every being of the Universe between birth and extinction[48].
|
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+
|
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In general, the Islamic and Judeo-Christian world-view regards time as linear[49]
|
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and directional,[50]
|
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beginning with the act of creation by God. The traditional Christian view sees time ending, teleologically,[51]
|
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with the eschatological end of the present order of things, the "end time".
|
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+
|
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In the Old Testament book Ecclesiastes, traditionally ascribed to Solomon (970–928 BC), time (as the Hebrew word עידן, זמן 'iddan(age, as in "Ice age") zĕman(time) is often translated) was traditionally regarded[by whom?] as a medium for the passage of predestined events.[citation needed] (Another word, زمان" זמן" zamān, meant time fit for an event, and is used as the modern Arabic, Persian, and Hebrew equivalent to the English word "time".)
|
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|
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The Greek language denotes two distinct principles, Chronos and Kairos. The former refers to numeric, or chronological, time. The latter, literally "the right or opportune moment", relates specifically to metaphysical or Divine time. In theology, Kairos is qualitative, as opposed to quantitative.[52]
|
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|
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In Greek mythology, Chronos (ancient Greek: Χρόνος) is identified as the Personification of Time. His name in Greek means "time" and is alternatively spelled Chronus (Latin spelling) or Khronos. Chronos is usually portrayed as an old, wise man with a long, gray beard, such as "Father Time". Some English words whose etymological root is khronos/chronos include chronology, chronometer, chronic, anachronism, synchronise, and chronicle.
|
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|
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According to Kabbalists, "time" is a paradox[53] and an illusion.[54] Both the future and the past are recognised to be combined and simultaneously present.[clarification needed]
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+
|
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Two contrasting viewpoints on time divide prominent philosophers. One view is that time is part of the fundamental structure of the universe – a dimension independent of events, in which events occur in sequence. Isaac Newton subscribed to this realist view, and hence it is sometimes referred to as Newtonian time.[55][56]
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The opposing view is that time does not refer to any kind of "container" that events and objects "move through", nor to any entity that "flows", but that it is instead part of a fundamental intellectual structure (together with space and number) within which humans sequence and compare events. This second view, in the tradition of Gottfried Leibniz[15] and Immanuel Kant,[57][58] holds that time is neither an event nor a thing, and thus is not itself measurable nor can it be travelled.
|
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|
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Furthermore, it may be that there is a subjective component to time, but whether or not time itself is "felt", as a sensation, or is a judgment, is a matter of debate.[2][6][7][59][60]
|
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+
|
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In Philosophy, time was questioned throughout the centuries; what time is and if it is real or not. Ancient Greek philosophers asked if time was linear or cyclical and if time was endless or finite.[61] These philosophers had different ways of explaining time; for instance, ancient Indian philosophers had something called the Wheel of Time. It is believed that there was repeating ages over the lifespan of the universe.[62] This led to beliefs like cycles of rebirth and reincarnation.[62] The Greek philosophers believe that the universe was infinite, and was an illusion to humans.[62] Plato believed that time was made by the Creator at the same instant as the heavens.[62] He also says that time is a period of motion of the heavenly bodies.[62] Aristotle believed that time correlated to movement, that time did not exist on its own but was relative to motion of objects.[62] he also believed that time was related to the motion of celestial bodies; the reason that humans can tell time was because of orbital periods and therefore there was a duration on time.[63]
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The Vedas, the earliest texts on Indian philosophy and Hindu philosophy dating back to the late 2nd millennium BC, describe ancient Hindu cosmology, in which the universe goes through repeated cycles of creation, destruction and rebirth, with each cycle lasting 4,320 million years.[64]
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Ancient Greek philosophers, including Parmenides and Heraclitus, wrote essays on the nature of time.[65]
|
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Plato, in the Timaeus, identified time with the period of motion of the heavenly bodies. Aristotle, in Book IV of his Physica defined time as 'number of movement in respect of the before and after'.[66]
|
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In Book 11 of his Confessions, St. Augustine of Hippo ruminates on the nature of time, asking, "What then is time? If no one asks me, I know: if I wish to explain it to one that asketh, I know not." He begins to define time by what it is not rather than what it is,[67]
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an approach similar to that taken in other negative definitions. However, Augustine ends up calling time a "distention" of the mind (Confessions 11.26) by which we simultaneously grasp the past in memory, the present by attention, and the future by expectation.
|
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|
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Isaac Newton believed in absolute space and absolute time; Leibniz believed that time and space are relational.[68]
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The differences between Leibniz's and Newton's interpretations came to a head in the famous Leibniz–Clarke correspondence.
|
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+
|
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Philosophers in the 17th and 18th century questioned if time was real and absolute, or if it was an intellectual concept that humans use to understand and sequence events.[61] These questions lead to realism vs anti-realism; the realists believed that time is a fundamental part of the universe, and be perceived by events happening in a sequence, in a dimension.[69] Isaac Newton said that we are merely occupying time, he also says that humans can only understand relative time.[69] Relative time is a measurement of objects in motion.[69] The anti-realists believed that time is merely a convenient intellectual concept for humans to understand events.[69] This means that time was useless unless there were objects that it could interact with, this was called relational time.[69] René Descartes, John Locke, and David Hume said that one's mind needs to acknowledge time, in order to understand what time is.[63] Immanuel Kant believed that we can not know what something is unless we experience it first hand.[70]
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|
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Immanuel Kant, Critique of Pure Reason (1781), trans. Vasilis Politis (London: Dent., 1991), p.54.
|
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Immanuel Kant, in the Critique of Pure Reason, described time as an a priori intuition that allows us (together with the other a priori intuition, space) to comprehend sense experience.[71]
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With Kant, neither space nor time are conceived as substances, but rather both are elements of a systematic mental framework that necessarily structures the experiences of any rational agent, or observing subject. Kant thought of time as a fundamental part of an abstract conceptual framework, together with space and number, within which we sequence events, quantify their duration, and compare the motions of objects. In this view, time does not refer to any kind of entity that "flows," that objects "move through," or that is a "container" for events. Spatial measurements are used to quantify the extent of and distances between objects, and temporal measurements are used to quantify the durations of and between events. Time was designated by Kant as the purest possible schema of a pure concept or category.
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|
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Henri Bergson believed that time was neither a real homogeneous medium nor a mental construct, but possesses what he referred to as Duration. Duration, in Bergson's view, was creativity and memory as an essential component of reality.[72]
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According to Martin Heidegger we do not exist inside time, we are time. Hence, the relationship to the past is a present awareness of having been, which allows the past to exist in the present. The relationship to the future is the state of anticipating a potential possibility, task, or engagement. It is related to the human propensity for caring and being concerned, which causes "being ahead of oneself" when thinking of a pending occurrence. Therefore, this concern for a potential occurrence also allows the future to exist in the present. The present becomes an experience, which is qualitative instead of quantitative. Heidegger seems to think this is the way that a linear relationship with time, or temporal existence, is broken or transcended.[73]
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We are not stuck in sequential time. We are able to remember the past and project into the future – we have a kind of random access to our representation of temporal existence; we can, in our thoughts, step out of (ecstasis) sequential time.[74]
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Modern philosophers asked: is time real or unreal, is time happening all at once or a duration, If time tensed or tenseless, and is there a future to be?[61] There is a theory called the tenseless or B-theory; this theory says that any tensed terminology can be replaced with tenseless terminology.[75] For example, "we will win the game" can be replaced with "we do win the game", taking out the future tense. On the other hand, there is a theory called the tense or A-theory; this theory says that our language has tense verbs for a reason and that the future can not be determined.[75] There is also something called imaginary time, this was from Stephen Hawking, he says that space and imaginary time are finite but have no boundaries.[75] Imaginary time is not real or unreal, it is something that is hard to visualize.[75] Philosophers can agree that physical time exists outside of the human mind and is objective, and psychological time is mind dependent and subjective.[63]
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|
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In 5th century BC Greece, Antiphon the Sophist, in a fragment preserved from his chief work On Truth, held that: "Time is not a reality (hypostasis), but a concept (noêma) or a measure (metron)." Parmenides went further, maintaining that time, motion, and change were illusions, leading to the paradoxes of his follower Zeno.[76] Time as an illusion is also a common theme in Buddhist thought.[77][78]
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|
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J. M. E. McTaggart's 1908 The Unreality of Time argues that, since every event has the characteristic of being both present and not present (i.e., future or past), that time is a self-contradictory idea (see also The flow of time).
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|
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These arguments often center on what it means for something to be unreal. Modern physicists generally believe that time is as real as space – though others, such as Julian Barbour in his book The End of Time, argue that quantum equations of the universe take their true form when expressed in the timeless realm containing every possible now or momentary configuration of the universe, called "platonia" by Barbour.[79]
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|
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A modern philosophical theory called presentism views the past and the future as human-mind interpretations of movement instead of real parts of time (or "dimensions") which coexist with the present. This theory rejects the existence of all direct interaction with the past or the future, holding only the present as tangible. This is one of the philosophical arguments against time travel. This contrasts with eternalism (all time: present, past and future, is real) and the growing block theory (the present and the past are real, but the future is not).
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|
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Until Einstein's reinterpretation of the physical concepts associated with time and space, time was considered to be the same everywhere in the universe, with all observers measuring the same time interval for any event.[80]
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Non-relativistic classical mechanics is based on this Newtonian idea of time.
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|
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Einstein, in his special theory of relativity,[81]
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postulated the constancy and finiteness of the speed of light for all observers. He showed that this postulate, together with a reasonable definition for what it means for two events to be simultaneous, requires that distances appear compressed and time intervals appear lengthened for events associated with objects in motion relative to an inertial observer.
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|
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The theory of special relativity finds a convenient formulation in Minkowski spacetime, a mathematical structure that combines three dimensions of space with a single dimension of time. In this formalism, distances in space can be measured by how long light takes to travel that distance, e.g., a light-year is a measure of distance, and a meter is now defined in terms of how far light travels in a certain amount of time. Two events in Minkowski spacetime are separated by an invariant interval, which can be either space-like, light-like, or time-like. Events that have a time-like separation cannot be simultaneous in any frame of reference, there must be a temporal component (and possibly a spatial one) to their separation. Events that have a space-like separation will be simultaneous in some frame of reference, and there is no frame of reference in which they do not have a spatial separation. Different observers may calculate different distances and different time intervals between two events, but the invariant interval between the events is independent of the observer (and his or her velocity).
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In non-relativistic classical mechanics, Newton's concept of "relative, apparent, and common time" can be used in the formulation of a prescription for the synchronization of clocks. Events seen by two different observers in motion relative to each other produce a mathematical concept of time that works sufficiently well for describing the everyday phenomena of most people's experience. In the late nineteenth century, physicists encountered problems with the classical understanding of time, in connection with the behavior of electricity and magnetism. Einstein resolved these problems by invoking a method of synchronizing clocks using the constant, finite speed of light as the maximum signal velocity. This led directly to the conclusion that observers in motion relative to one another measure different elapsed times for the same event.
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Time has historically been closely related with space, the two together merging into spacetime in Einstein's special relativity and general relativity. According to these theories, the concept of time depends on the spatial reference frame of the observer, and the human perception as well as the measurement by instruments such as clocks are different for observers in relative motion. For example, if a spaceship carrying a clock flies through space at (very nearly) the speed of light, its crew does not notice a change in the speed of time on board their vessel because everything traveling at the same speed slows down at the same rate (including the clock, the crew's thought processes, and the functions of their bodies). However, to a stationary observer watching the spaceship fly by, the spaceship appears flattened in the direction it is traveling and the clock on board the spaceship appears to move very slowly.
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On the other hand, the crew on board the spaceship also perceives the observer as slowed down and flattened along the spaceship's direction of travel, because both are moving at very nearly the speed of light relative to each other. Because the outside universe appears flattened to the spaceship, the crew perceives themselves as quickly traveling between regions of space that (to the stationary observer) are many light years apart. This is reconciled by the fact that the crew's perception of time is different from the stationary observer's; what seems like seconds to the crew might be hundreds of years to the stationary observer. In either case, however, causality remains unchanged: the past is the set of events that can send light signals to an entity and the future is the set of events to which an entity can send light signals.[82][83]
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Einstein showed in his thought experiments that people travelling at different speeds, while agreeing on cause and effect, measure different time separations between events, and can even observe different chronological orderings between non-causally related events. Though these effects are typically minute in the human experience, the effect becomes much more pronounced for objects moving at speeds approaching the speed of light. Subatomic particles exist for a well known average fraction of a second in a lab relatively at rest, but when travelling close to the speed of light they are measured to travel farther and exist for much longer than when at rest. According to the special theory of relativity, in the high-speed particle's frame of reference, it exists, on the average, for a standard amount of time known as its mean lifetime, and the distance it travels in that time is zero, because its velocity is zero. Relative to a frame of reference at rest, time seems to "slow down" for the particle. Relative to the high-speed particle, distances seem to shorten. Einstein showed how both temporal and spatial dimensions can be altered (or "warped") by high-speed motion.
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Einstein (The Meaning of Relativity): "Two events taking place at the points A and B of a system K are simultaneous if they appear at the same instant when observed from the middle point, M, of the interval AB. Time is then defined as the ensemble of the indications of similar clocks, at rest relative to K, which register the same simultaneously."
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Einstein wrote in his book, Relativity, that simultaneity is also relative, i.e., two events that appear simultaneous to an observer in a particular inertial reference frame need not be judged as simultaneous by a second observer in a different inertial frame of reference.
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The animations visualise the different treatments of time in the Newtonian and the relativistic descriptions. At the heart of these differences are the Galilean and Lorentz transformations applicable in the Newtonian and relativistic theories, respectively.
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In the figures, the vertical direction indicates time. The horizontal direction indicates distance (only one spatial dimension is taken into account), and the thick dashed curve is the spacetime trajectory ("world line") of the observer. The small dots indicate specific (past and future) events in spacetime.
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The slope of the world line (deviation from being vertical) gives the relative velocity to the observer. Note how in both pictures the view of spacetime changes when the observer accelerates.
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In the Newtonian description these changes are such that time is absolute:[84] the movements of the observer do not influence whether an event occurs in the 'now' (i.e., whether an event passes the horizontal line through the observer).
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However, in the relativistic description the observability of events is absolute: the movements of the observer do not influence whether an event passes the "light cone" of the observer. Notice that with the change from a Newtonian to a relativistic description, the concept of absolute time is no longer applicable: events move up-and-down in the figure depending on the acceleration of the observer.
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Time appears to have a direction – the past lies behind, fixed and immutable, while the future lies ahead and is not necessarily fixed. Yet for the most part the laws of physics do not specify an arrow of time, and allow any process to proceed both forward and in reverse. This is generally a consequence of time being modelled by a parameter in the system being analysed, where there is no "proper time": the direction of the arrow of time is sometimes arbitrary. Examples of this include the cosmological arrow of time, which points away from the Big Bang, CPT symmetry, and the radiative arrow of time, caused by light only travelling forwards in time (see light cone). In particle physics, the violation of CP symmetry implies that there should be a small counterbalancing time asymmetry to preserve CPT symmetry as stated above. The standard description of measurement in quantum mechanics is also time asymmetric (see Measurement in quantum mechanics). The second law of thermodynamics states that entropy must increase over time (see Entropy). This can be in either direction – Brian Greene theorizes that, according to the equations, the change in entropy occurs symmetrically whether going forward or backward in time. So entropy tends to increase in either direction, and our current low-entropy universe is a statistical aberration, in the similar manner as tossing a coin often enough that eventually heads will result ten times in a row. However, this theory is not supported empirically in local experiment.[85]
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Time quantization is a hypothetical concept. In the modern established physical theories (the Standard Model of Particles and Interactions and General Relativity) time is not quantized.
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Planck time (~ 5.4 × 10−44 seconds) is the unit of time in the system of natural units known as Planck units. Current established physical theories are believed to fail at this time scale, and many physicists expect that the Planck time might be the smallest unit of time that could ever be measured, even in principle. Tentative physical theories that describe this time scale exist; see for instance loop quantum gravity.
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Time travel is the concept of moving backwards or forwards to different points in time, in a manner analogous to moving through space, and different from the normal "flow" of time to an earthbound observer. In this view, all points in time (including future times) "persist" in some way. Time travel has been a plot device in fiction since the 19th century. Travelling backwards in time has never been verified, presents many theoretical problems, and may be an impossibility.[86] Any technological device, whether fictional or hypothetical, that is used to achieve time travel is known as a time machine.
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A central problem with time travel to the past is the violation of causality; should an effect precede its cause, it would give rise to the possibility of a temporal paradox. Some interpretations of time travel resolve this by accepting the possibility of travel between branch points, parallel realities, or universes.
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Another solution to the problem of causality-based temporal paradoxes is that such paradoxes cannot arise simply because they have not arisen. As illustrated in numerous works of fiction, free will either ceases to exist in the past or the outcomes of such decisions are predetermined. As such, it would not be possible to enact the grandfather paradox because it is a historical fact that one's grandfather was not killed before his child (one's parent) was conceived. This view does not simply hold that history is an unchangeable constant, but that any change made by a hypothetical future time traveller would already have happened in his or her past, resulting in the reality that the traveller moves from. More elaboration on this view can be found in the Novikov self-consistency principle.
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The specious present refers to the time duration wherein one's perceptions are considered to be in the present. The experienced present is said to be 'specious' in that, unlike the objective present, it is an interval and not a durationless instant. The term specious present was first introduced by the psychologist E.R. Clay, and later developed by William James.[87]
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The brain's judgment of time is known to be a highly distributed system, including at least the cerebral cortex, cerebellum and basal ganglia as its components. One particular component, the suprachiasmatic nuclei, is responsible for the circadian (or daily) rhythm, while other cell clusters appear capable of shorter-range (ultradian) timekeeping.
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Psychoactive drugs can impair the judgment of time. Stimulants can lead both humans and rats to overestimate time intervals,[88][89] while depressants can have the opposite effect.[90] The level of activity in the brain of neurotransmitters such as dopamine and norepinephrine may be the reason for this.[91] Such chemicals will either excite or inhibit the firing of neurons in the brain, with a greater firing rate allowing the brain to register the occurrence of more events within a given interval (speed up time) and a decreased firing rate reducing the brain's capacity to distinguish events occurring within a given interval (slow down time).[92]
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Mental chronometry is the use of response time in perceptual-motor tasks to infer the content, duration, and temporal sequencing of cognitive operations.
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Children's expanding cognitive abilities allow them to understand time more clearly. Two- and three-year-olds' understanding of time is mainly limited to "now and not now". Five- and six-year-olds can grasp the ideas of past, present, and future. Seven- to ten-year-olds can use clocks and calendars.[93]
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In addition to psychoactive drugs, judgments of time can be altered by temporal illusions (like the kappa effect),[94] age,[95] and hypnosis.[96] The sense of time is impaired in some people with neurological diseases such as Parkinson's disease and attention deficit disorder.
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Psychologists assert that time seems to go faster with age, but the literature on this age-related perception of time remains controversial.[97] Those who support this notion argue that young people, having more excitatory neurotransmitters, are able to cope with faster external events.[92]
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In sociology and anthropology, time discipline is the general name given to social and economic rules, conventions, customs, and expectations governing the measurement of time, the social currency and awareness of time measurements, and people's expectations concerning the observance of these customs by others. Arlie Russell Hochschild[98][99] and Norbert Elias[100] have written on the use of time from a sociological perspective.
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The use of time is an important issue in understanding human behavior, education, and travel behavior. Time-use research is a developing field of study. The question concerns how time is allocated across a number of activities (such as time spent at home, at work, shopping, etc.). Time use changes with technology, as the television or the Internet created new opportunities to use time in different ways. However, some aspects of time use are relatively stable over long periods of time, such as the amount of time spent traveling to work, which despite major changes in transport, has been observed to be about 20–30 minutes one-way for a large number of cities over a long period.
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Time management is the organization of tasks or events by first estimating how much time a task requires and when it must be completed, and adjusting events that would interfere with its completion so it is done in the appropriate amount of time. Calendars and day planners are common examples of time management tools.
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A sequence of events, or series of events, is a sequence of items, facts, events, actions, changes, or procedural steps, arranged in time order (chronological order), often with causality relationships among the items.[101][102][103]
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Because of causality, cause precedes effect, or cause and effect may appear together in a single item, but effect never precedes cause. A sequence of events can be presented in text, tables, charts, or timelines. The description of the items or events may include a timestamp. A sequence of events that includes the time along with place or location information to describe a sequential path may be referred to as a world line.
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Uses of a sequence of events include stories,[104]
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historical events (chronology), directions and steps in procedures,[105]
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and timetables for scheduling activities. A sequence of events may also be used to help describe processes in science, technology, and medicine. A sequence of events may be focused on past events (e.g., stories, history, chronology), on future events that must be in a predetermined order (e.g., plans, schedules, procedures, timetables), or focused on the observation of past events with the expectation that the events will occur in the future (e.g., processes, projections). The use of a sequence of events occurs in fields as diverse as machines (cam timer), documentaries (Seconds From Disaster), law (choice of law), finance (directional-change intrinsic time), computer simulation (discrete event simulation), and electric power transmission[106]
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(sequence of events recorder). A specific example of a sequence of events is the timeline of the Fukushima Daiichi nuclear disaster.
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Although time is regarded as an abstract concept, there is increasing evidence that time is conceptualized in the mind in terms of space.[107] That is, instead of thinking about time in a general, abstract way, humans think about time in a spatial way and mentally organize it as such. Using space to think about time allows humans to mentally organize temporal events in a specific way.
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This spatial representation of time is often represented in the mind as a Mental Time Line (MTL).[108] Using space to think about time allows humans to mentally organize temporal order. These origins are shaped by many environmental factors[107]––for example, literacy appears to play a large role in the different types of MTLs, as reading/writing direction provides an everyday temporal orientation that differs from culture to culture.[108] In western cultures, the MTL may unfold rightward (with the past on the left and the future on the right) since people read and write from left to right.[108] Western calendars also continue this trend by placing the past on the left with the future progressing toward the right. Conversely, Arabic, Farsi, Urdu and Israeli-Hebrew speakers read from right to left, and their MTLs unfold leftward (past on the right with future on the left), and evidence suggests these speakers organize time events in their minds like this as well.[108]
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This linguistic evidence that abstract concepts are based in spatial concepts also reveals that the way humans mentally organize time events varies across cultures––that is, a certain specific mental organization system is not universal. So, although Western cultures typically associate past events with the left and future events with the right according to a certain MTL, this kind of horizontal, egocentric MTL is not the spatial organization of all cultures. Although most developed nations use an egocentric spatial system, there is recent evidence that some cultures use an allocentric spatialization, often based on environmental features.[107]
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A recent study of the indigenous Yupno people of Papua New Guinea focused on the directional gestures used when individuals used time-related words.[107] When speaking of the past (such as "last year" or "past times"), individuals gestured downhill, where the river of the valley flowed into the ocean. When speaking of the future, they gestured uphill, toward the source of the river. This was common regardless of which direction the person faced, revealing that the Yupno people may use an allocentric MTL, in which time flows uphill.[107]
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A similar study of the Pormpuraawans, an aboriginal group in Australia, revealed a similar distinction in which when asked to organize photos of a man aging "in order," individuals consistently placed the youngest photos to the east and the oldest photos to the west, regardless of which direction they faced.[109] This directly clashed with an American group which consistently organized the photos from left to right. Therefore, this group also appears to have an allocentric MTL, but based on the cardinal directions instead of geographical features.[109]
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The wide array of distinctions in the way different groups think about time leads to the broader question that different groups may also think about other abstract concepts in different ways as well, such as causality and number.[107]
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A tendon or sinew is a tough band of fibrous connective tissue that connects muscle to bone and is capable of withstanding tension.
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Tendons are similar to ligaments; both are made of collagen. Ligaments connect one bone to another, while tendons connect muscle to bone.
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Histologically, tendons consist of dense regular connective tissue. The main cellular component of tendons are specialized fibroblasts called tenocytes. Tenocytes synthesize the extracellular matrix of tendons, abundant in densely packed collagen fibers. The collagen fibers are parallel to each other and organized into fascicles. Individual fascicles are bound by the endotendineum, which is a delicate loose connective tissue containing thin collagen fibrils[1][2] and elastic fibres.[3] Groups of fascicles are bounded by the epitenon, which is a sheath of dense irregular connective tissue. The whole tendon is enclosed by a fascia. The space between the fascia and the tendon tissue is filled with the paratenon, a fatty areolar tissue.[4] Normal healthy tendons are anchored to bone by Sharpey's fibres.
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The dry mass of normal tendons, which makes up 30-45% of their total mass, is composed of:
|
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While collagen I makes up most of the collagen in tendon, many minor collagens are present that play vital roles in proper tendon development and function. These include type II collagen in the cartilaginous zones, type III collagen in the reticulin fibres of the vascular walls, type IX collagen, type IV collagen in the basement membranes of the capillaries, type V collagen in the vascular walls, and type X collagen in the mineralized fibrocartilage near the interface with the bone.[5][9]
|
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Collagen fibres coalesce into macroaggregates. After secretion from the cell, cleaved by procollagen N- and C-proteases, the tropocollagen molecules spontaneously assemble into insoluble fibrils. A collagen molecule is about 300 nm long and 1–2 nm wide, and the diameter of the fibrils that are formed can range from 50–500 nm. In tendons, the fibrils then assemble further to form fascicles, which are about 10 mm in length with a diameter of 50–300 μm, and finally into a tendon fibre with a diameter of 100–500 μm.[10]
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The collagen in tendons are held together with proteoglycan (a compound consisting of a protein bonded to glycosaminoglycan groups, present especially in connective tissue) components including decorin and, in compressed regions of tendon, aggrecan, which are capable of binding to the collagen fibrils at specific locations.[11] The proteoglycans are interwoven with the collagen fibrils – their glycosaminoglycan (GAG) side chains have multiple interactions with the surface of the fibrils – showing that the proteoglycans are important structurally in the interconnection of the fibrils.[12] The major GAG components of the tendon are dermatan sulfate and chondroitin sulfate, which associate with collagen and are involved in the fibril assembly process during tendon development. Dermatan sulfate is thought to be responsible for forming associations between fibrils, while chondroitin sulfate is thought to be more involved with occupying volume between the fibrils to keep them separated and help withstand deformation.[13] The dermatan sulfate side chains of decorin aggregate in solution, and this behavior can assist with the assembly of the collagen fibrils. When decorin molecules are bound to a collagen fibril, their dermatan sulfate chains may extend and associate with other dermatan sulfate chains on decorin that is bound to separate fibrils, therefore creating interfibrillar bridges and eventually causing parallel alignment of the fibrils.[14]
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The tenocytes produce the collagen molecules, which aggregate end-to-end and side-to-side to produce collagen fibrils. Fibril bundles are organized to form fibres with the elongated tenocytes closely packed between them. There is a three-dimensional network of cell processes associated with collagen in the tendon. The cells communicate with each other through gap junctions, and this signalling gives them the ability to detect and respond to mechanical loading.[15]
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Blood vessels may be visualized within the endotendon running parallel to collagen fibres, with occasional branching transverse anastomoses.
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The internal tendon bulk is thought to contain no nerve fibres, but the epitenon and paratenon contain nerve endings, while Golgi tendon organs are present at the junction between tendon and muscle.
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Tendon length varies in all major groups and from person to person. Tendon length is, in practice, the deciding factor regarding actual and potential muscle size. For example, all other relevant biological factors being equal, a man with a shorter tendons and a longer biceps muscle will have greater potential for muscle mass than a man with a longer tendon and a shorter muscle. Successful bodybuilders will generally have shorter tendons. Conversely, in sports requiring athletes to excel in actions such as running or jumping, it is beneficial to have longer than average Achilles tendon and a shorter calf muscle.[16]
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Tendon length is determined by genetic predisposition, and has not been shown to either increase or decrease in response to environment, unlike muscles, which can be shortened by trauma, use imbalances and a lack of recovery and stretching.[17]
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Traditionally, tendons have been considered to be a mechanism by which muscles connect to bone as well as muscles itself, functioning to transmit forces. This connection allows tendons to passively modulate forces during locomotion, providing additional stability with no active work. However, over the past two decades, much research focused on the elastic properties of some tendons and their ability to function as springs. Not all tendons are required to perform the same functional role, with some predominantly positioning limbs, such as the fingers when writing (positional tendons) and others acting as springs to make locomotion more efficient (energy storing tendons).[18] Energy storing tendons can store and recover energy at high efficiency. For example, during a human stride, the Achilles tendon stretches as the ankle joint dorsiflexes. During the last portion of the stride, as the foot plantar-flexes (pointing the toes down), the stored elastic energy is released. Furthermore, because the tendon stretches, the muscle is able to function with less or even no change in length, allowing the muscle to generate greater force.
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The mechanical properties of the tendon are dependent on the collagen fiber diameter and orientation. The collagen fibrils are parallel to each other and closely packed, but show a wave-like appearance due to planar undulations, or crimps, on a scale of several micrometers.[19] In tendons, the collagen fibres have some flexibility due to the absence of hydroxyproline and proline residues at specific locations in the amino acid sequence, which allows the formation of other conformations such as bends or internal loops in the triple helix and results in the development of crimps.[20] The crimps in the collagen fibrils allow the tendons to have some flexibility as well as a low compressive stiffness. In addition, because the tendon is a multi-stranded structure made up of many partially independent fibrils and fascicles, it does not behave as a single rod, and this property also contributes to its flexibility.[21]
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The proteoglycan components of tendons also are important to the mechanical properties. While the collagen fibrils allow tendons to resist tensile stress, the proteoglycans allow them to resist compressive stress. These molecules are very hydrophilic, meaning that they can absorb a large amount of water and therefore have a high swelling ratio. Since they are noncovalently bound to the fibrils, they may reversibly associate and disassociate so that the bridges between fibrils can be broken and reformed. This process may be involved in allowing the fibril to elongate and decrease in diameter under tension.[22] However, the proteoglycans may also have a role in the tensile properties of tendon. The structure of tendon is effectively a fibre composite material, built as a series of hierarchical levels. At each level of the hierarchy, the collagen units are bound together by either collagen crosslinks, or the proteoglycans, to create a structure highly resistant to tensile load.[23] The elongation and the strain of the collagen fibrils alone have been shown to be much lower than the total elongation and strain of the entire tendon under the same amount of stress, demonstrating that the proteoglycan-rich matrix must also undergo deformation, and stiffening of the matrix occurs at high strain rates.[24] This deformation of the non-collagenous matrix occurs at all levels of the tendon hierarchy, and by modulating the organisation and structure of this matrix, the different mechanical properties required by different tendons can be achieved.[25] Energy storing tendons have been shown to utilise significant amounts of sliding between fascicles to enable the high strain characteristics they require, whilst positional tendons rely more heavily on sliding between collagen fibres and fibrils.[26] However, recent data suggests that energy storing tendons may also contain fascicles which are twisted, or helical, in nature - an arrangement that would be highly beneficial for providing the spring-like behaviour required in these tendons.[27]
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Tendons are viscoelastic structures, which means they exhibit both elastic and viscous behaviour. When stretched, tendons exhibit typical "soft tissue" behavior. The force-extension, or stress-strain curve starts with a very low stiffness region, as the crimp structure straightens and the collagen fibres align suggesting negative Poisson's ratio in the fibres of the tendon. More recently, tests carried out in vivo (through MRI) and ex vivo (through mechanical testing of various cadaveric tendon tissue) have shown that healthy tendons are highly anisotropic and exhibit a negative Poisson's ratio (auxetic) in some planes when stretched up to 2% along their length, i.e. within their normal range of motion.[28] After this 'toe' region, the structure becomes significantly stiffer, and has a linear stress-strain curve until it begins to fail. The mechanical properties of tendons vary widely, as they are matched to the functional requirements of the tendon. The energy storing tendons tend to be more elastic, or less stiff, so they can more easily store energy, whilst the stiffer positional tendons tend to be a little more viscoelastic, and less elastic, so they can provide finer control of movement. A typical energy storing tendon will fail at around 12-15% strain, and a stress in the region of 100-150 MPa, although some tendons are notably more extensible than this, for example the superficial digital flexor in the horse, which stretches in excess of 20% when galloping.[29] Positional tendons can fail at strains as low as 6-8%, but can have moduli in the region of 700-1000 MPa.[30]
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Several studies have demonstrated that tendons respond to changes in mechanical loading with growth and remodeling processes, much like bones. In particular, a study showed that disuse of the Achilles tendon in rats resulted in a decrease in the average thickness of the collagen fiber bundles comprising the tendon.[31] In humans, an experiment in which people were subjected to a simulated micro-gravity environment found that tendon stiffness decreased significantly, even when subjects were required to perform restiveness exercises.[32] These effects have implications in areas ranging from treatment of bedridden patients to the design of more effective exercises for astronauts.
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The tendons in the foot are highly complex and intricate. Therefore, the healing process for a broken tendon is long and painful. Most people who do not receive medical attention within the first 48 hours of the injury will suffer from severe swelling, pain, and a burning sensation where the injury occurred.
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It was believed that tendons could not undergo matrix turnover and that tenocytes were not capable of repair. However, it has since been shown that, throughout the lifetime of a person, tenocytes in the tendon actively synthesize matrix components as well as enzymes such as matrix metalloproteinases (MMPs) can degrade the matrix.[33] Tendons are capable of healing and recovering from injuries in a process that is controlled by the tenocytes and their surrounding extracellular matrix.
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The three main stages of tendon healing are inflammation, repair or proliferation, and remodeling, which can be further divided into consolidation and maturation. These stages can overlap with each other. In the first stage, inflammatory cells such as neutrophils are recruited to the injury site, along with erythrocytes. Monocytes and macrophages are recruited within the first 24 hours, and phagocytosis of necrotic materials at the injury site occurs. After the release of vasoactive and chemotactic factors, angiogenesis and the proliferation of tenocytes are initiated. Tenocytes then move into the site and start to synthesize collagen III.[34][35] After a few days, the repair or proliferation stage begins. In this stage, the tenocytes are involved in the synthesis of large amounts of collagen and proteoglycans at the site of injury, and the levels of GAG and water are high.[36] After about six weeks, the remodeling stage begins. The first part of this stage is consolidation, which lasts from about six to ten weeks after the injury. During this time, the synthesis of collagen and GAGs is decreased, and the cellularity is also decreased as the tissue becomes more fibrous as a result of increased production of collagen I and the fibrils become aligned in the direction of mechanical stress.[35] The final maturation stage occurs after ten weeks, and during this time there is an increase in crosslinking of the collagen fibrils, which causes the tissue to become stiffer. Gradually, over about one year, the tissue will turn from fibrous to scar-like.[36]
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Matrix metalloproteinases (MMPs) have a very important role in the degradation and remodeling of the ECM during the healing process after a tendon injury. Certain MMPs including MMP-1, MMP-2, MMP-8, MMP-13, and MMP-14 have collagenase activity, meaning that, unlike many other enzymes, they are capable of degrading collagen I fibrils. The degradation of the collagen fibrils by MMP-1 along with the presence of denatured collagen are factors that are believed to cause weakening of the tendon ECM and an increase in the potential for another rupture to occur.[37] In response to repeated mechanical loading or injury, cytokines may be released by tenocytes and can induce the release of MMPs, causing degradation of the ECM and leading to recurring injury and chronic tendinopathies.[35]
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A variety of other molecules are involved in tendon repair and regeneration. There are five growth factors that have been shown to be significantly upregulated and active during tendon healing: insulin-like growth factor 1 (IGF-I), platelet-derived growth factor (PDGF), vascular endothelial growth factor (VEGF), basic fibroblast growth factor (bFGF), and transforming growth factor beta (TGF-β).[36] These growth factors all have different roles during the healing process. IGF-1 increases collagen and proteoglycan production during the first stage of inflammation, and PDGF is also present during the early stages after injury and promotes the synthesis of other growth factors along with the synthesis of DNA and the proliferation of tendon cells.[36] The three isoforms of TGF-β (TGF-β1, TGF-β2, TGF-β3) are known to play a role in wound healing and scar formation.[38] VEGF is well known to promote angiogenesis and to induce endothelial cell proliferation and migration, and VEGF mRNA has been shown to be expressed at the site of tendon injuries along with collagen I mRNA.[39] Bone morphogenetic proteins (BMPs) are a subgroup of TGF-β superfamily that can induce bone and cartilage formation as well as tissue differentiation, and BMP-12 specifically has been shown to influence formation and differentiation of tendon tissue and to promote fibrogenesis.
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In animal models, extensive studies have been conducted to investigate the effects of mechanical strain in the form of activity level on tendon injury and healing. While stretching can disrupt healing during the initial inflammatory phase, it has been shown that controlled movement of the tendons after about one week following an acute injury can help to promote the synthesis of collagen by the tenocytes, leading to increased tensile strength and diameter of the healed tendons and fewer adhesions than tendons that are immobilized. In chronic tendon injuries, mechanical loading has also been shown to stimulate fibroblast proliferation and collagen synthesis along with collagen realignment, all of which promote repair and remodeling.[36] To further support the theory that movement and activity assist in tendon healing, it has been shown that immobilization of the tendons after injury often has a negative effect on healing. In rabbits, collagen fascicles that are immobilized have shown decreased tensile strength, and immobilization also results in lower amounts of water, proteoglycans, and collagen crosslinks in the tendons.[34]
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Several mechanotransduction mechanisms have been proposed as reasons for the response of tenocytes to mechanical force that enable them to alter their gene expression, protein synthesis, and cell phenotype, and eventually cause changes in tendon structure. A major factor is mechanical deformation of the extracellular matrix, which can affect the actin cytoskeleton and therefore affect cell shape, motility, and function. Mechanical forces can be transmitted by focal adhesion sites, integrins, and cell-cell junctions. Changes in the actin cytoskeleton can activate integrins, which mediate “outside-in” and “inside-out” signaling between the cell and the matrix. G-proteins, which induce intracellular signaling cascades, may also be important, and ion channels are activated by stretching to allow ions such as calcium, sodium, or potassium to enter the cell.[36]
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Sinew was widely used throughout pre-industrial eras as a tough, durable fiber. Some specific uses include using sinew as thread for sewing, attaching feathers to arrows (see fletch), lashing tool blades to shafts, etc. It is also recommended in survival guides as a material from which strong cordage can be made for items like traps or living structures. Tendon must be treated in specific ways to function usefully for these purposes. Inuit and other circumpolar people utilized sinew as the only cordage for all domestic purposes due to the lack of other suitable fiber sources in their ecological habitats. The elastic properties of particular sinews were also used in composite recurved bows favoured by the steppe nomads of Eurasia, and Native Americans. The first stone throwing artillery also used the elastic properties of sinew.
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Sinew makes for an excellent cordage material for three reasons: It is extremely strong, it contains natural glues, and it shrinks as it dries, doing away with the need for knots.
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Tendon (in particular, beef tendon) is used as a food in some Asian cuisines (often served at yum cha or dim sum restaurants). One popular dish is suan bao niu jin, in which the tendon is marinated in garlic. It is also sometimes found in the Vietnamese noodle dish phở.
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Tendons are subject to many types of injuries. There are various forms of tendinopathies or tendon injuries due to overuse. These types of injuries generally result in inflammation and degeneration or weakening of the tendons, which may eventually lead to tendon rupture.[34] Tendinopathies can be caused by a number of factors relating to the tendon extracellular matrix (ECM), and their classification has been difficult because their symptoms and histopathology often are similar.
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The first category of tendinopathy is paratenonitis, which refers to inflammation of the paratenon, or paratendinous sheet located between the tendon and its sheath. Tendinosis refers to non-inflammatory injury to the tendon at the cellular level. The degradation is caused by damage to collagen, cells, and the vascular components of the tendon, and is known to lead to rupture.[40] Observations of tendons that have undergone spontaneous rupture have shown the presence of collagen fibrils that are not in the correct parallel orientation or are not uniform in length or diameter, along with rounded tenocytes, other cell abnormalities, and the ingrowth of blood vessels.[34] Other forms of tendinosis that have not led to rupture have also shown the degeneration, disorientation, and thinning of the collagen fibrils, along with an increase in the amount of glycosaminoglycans between the fibrils.[35] The third is paratenonitis with tendinosis, in which combinations of paratenon inflammation and tendon degeneration are both present. The last is tendinitis, which refers to degeneration with inflammation of the tendon as well as vascular disruption.[5]
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Tendinopathies may be caused by several intrinsic factors including age, body weight, and nutrition. The extrinsic factors are often related to sports and include excessive forces or loading, poor training techniques, and environmental conditions.[33]
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In some organisms, notable ones being birds[41] and ornithischian dinosaurs,[42] portions of the tendon can become ossified. In this process, osteocytes infiltrate the tendon and lay down bone as they would in sesamoid bone such as the patella. In birds, tendon ossification primarily occurs in the hindlimb, while in ornithischian dinosaurs, ossified axial muscle tendons form a latticework along the neural and haemal spines on the tail, presumably for support.
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1 |
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Tennessee (/ˌtɛnəˈsiː/ (listen),[7][8] locally /ˈtɛnəsi/[9]), officially the State of Tennessee, is a state in the southeastern United States. Tennessee is the 36th largest by area and the 16th most populous of the 50 states. It is bordered by eight states, with Kentucky to the north, Virginia to the northeast, North Carolina to the east, Georgia, Alabama, and Mississippi to the south, Arkansas to the west, and Missouri to the northwest. The Appalachian Mountains dominate the eastern part of the state, and the Mississippi River forms its western border. Nashville is the state's capital and largest city, with a 2019 population of 670,820 and a 2019 metro population of 1,934,317. Tennessee's second largest city is Memphis, which had a population of 651,073 and metro population of 1,346,045 in 2019.[10]
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The state of Tennessee is rooted in the Watauga Association, a 1772 frontier pact generally regarded as the first constitutional government west of the Appalachians.[11] What is now Tennessee was initially part of North Carolina, and later part of the Southwest Territory. Tennessee was admitted to the Union as the 16th state on June 1, 1796. Tennessee was the last state to leave the Union and join the Confederacy at the outbreak of the American Civil War in 1861. Occupied by Union forces from 1862, it was the first state to be readmitted to the Union at the end of the war.[12]
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Tennessee furnished more soldiers for the Confederate Army than any other state besides Virginia, and more soldiers for the Union Army than the rest of the Confederacy combined.[12] Beginning during Reconstruction, it had competitive party politics, but a Democratic takeover in the late 1880s resulted in passage of disenfranchisement laws that excluded most blacks and many poor whites from voting. This sharply reduced competition in politics in the state until after passage of civil rights legislation in the mid-20th century.[13] In the 20th century, Tennessee transitioned from an agrarian economy to a more diversified economy, aided by massive federal investment in the Tennessee Valley Authority and, in the early 1940s, the city of Oak Ridge. This city was established just outside of Knoxville to house the Manhattan Project's uranium enrichment facilities, helping to build the world's first atomic bombs, two of which were dropped on Imperial Japan near the end of World War II. After the war, the Oak Ridge National Laboratory became and remains a key center for scientific research. In 2016, the element tennessine was named for the state, largely in recognition of the roles played by Oak Ridge, Vanderbilt University, and the University of Tennessee in the element’s discovery.[14]
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Tennessee's major industries include agriculture, manufacturing, and tourism. Poultry, soybeans, and cattle are the state's primary agricultural products,[15] and major manufacturing exports include chemicals, transportation equipment, and electrical equipment.[16] The Great Smoky Mountains National Park, the nation's most visited national park, is located in the eastern part of the state, and a section of the Appalachian Trail roughly follows the Tennessee–North Carolina border.[17] Other major tourist attractions include the Tennessee Aquarium and Chattanooga Choo-Choo Hotel in Chattanooga; Dollywood in Pigeon Forge; Ripley's Aquarium of the Smokies and Ober Gatlinburg in Gatlinburg; the Parthenon, the Country Music Hall of Fame and Museum, and Ryman Auditorium in Nashville; the Jack Daniel's Distillery in Lynchburg; Elvis Presley's Graceland residence and tomb, the Memphis Zoo, the National Civil Rights Museum in Memphis; and Bristol Motor Speedway in Bristol.
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The earliest variant of the name that became Tennessee was recorded by Captain Juan Pardo, the Spanish explorer, when he and his men passed through an American Indian village named "Tanasqui" in 1567 while traveling inland from South Carolina. In the early 18th century, British traders encountered a Cherokee town named Tanasi (or "Tanase") in present-day Monroe County, Tennessee. The town was located on a river of the same name (now known as the Little Tennessee River), and appears on maps as early as 1725. It is not known whether this was the same town as the one encountered by Juan Pardo, although recent research suggests that Pardo's "Tanasqui" was located at the confluence of the Pigeon River and the French Broad River, near modern Newport.[18]
|
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The meaning and origin of the word are uncertain. Some accounts suggest it is a Cherokee modification of an earlier Yuchi word. It has been said to mean "meeting place", "winding river", or "river of the great bend".[19][20] According to ethnographer James Mooney, the name "can not be analyzed" and its meaning is lost.[21]
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The modern spelling, Tennessee, is attributed to James Glen, the governor of South Carolina, who used this spelling in his official correspondence during the 1750s. The spelling was popularized by the publication of Henry Timberlake's "Draught of the Cherokee Country" in 1765. In 1788, North Carolina created "Tennessee County", the third county to be established in what is now Middle Tennessee. (Tennessee County was the predecessor to present-day Montgomery and Robertson counties.) When a constitutional convention met in 1796 to organize a new state out of the Southwest Territory, it adopted "Tennessee" as the name of the state.
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Tennessee is known as The Volunteer State, a nickname some claimed was earned during the War of 1812 because of the prominent role played by volunteer soldiers from Tennessee, especially during the Battle of New Orleans.[22] Other sources differ on the origin of the state nickname; according to The Columbia Encyclopedia,[23] the name refers to volunteers for the Mexican–American War from 1846 to 1848. This explanation is more likely, because President Polk's call for 2,600 nationwide volunteers at the beginning of the Mexican–American War resulted in 30,000 volunteers from Tennessee alone, largely in response to the death of Davy Crockett and appeals by former Tennessee Governor and then Texas politician, Sam Houston.[24]
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Tennessee borders eight other states: Kentucky and Virginia to the north; North Carolina to the east; Georgia, Alabama, and Mississippi on the south; Arkansas and Missouri on the Mississippi River to the west. Tennessee is tied with Missouri as the state bordering the most other states. The state is trisected by the Tennessee River.
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The highest point in the state is Clingmans Dome at 6,643 feet (2,025 m).[25] Clingmans Dome, which lies on Tennessee's eastern border, is the highest point on the Appalachian Trail, and is the third highest peak in the United States east of the Mississippi River. The state line between Tennessee and North Carolina crosses the summit. The state's lowest point is the Mississippi River at the Mississippi state line: 178 feet (54 m). The geographical center of the state is located in Murfreesboro.
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The state of Tennessee is geographically, culturally, economically, and legally divided into three Grand Divisions: East Tennessee, Middle Tennessee, and West Tennessee. The state constitution allows no more than two justices of the five-member Tennessee Supreme Court to be from one Grand Division and a similar rule applies to certain commissions and boards.
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Tennessee features six principal physiographic regions: the Blue Ridge, the Appalachian Ridge and Valley Region, the Cumberland Plateau, the Highland Rim, the Nashville Basin, and the Gulf Coastal Plain. Tennessee is home to the most caves in the United States, with more than ten thousand documented.[26]
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About half the state is in the Tennessee Valley drainage basin of the Tennessee River.[27] Approximately the northern half of Middle Tennessee, including Nashville and Clarksville, and a small portion of East Tennessee is in the Cumberland River basin.[28] A small part of north-central Tennessee in Sumner, Macon, and Clay counties is in the Green River watershed.[29] All three of these basins are tributaries of the Ohio River watershed. Most of West Tennessee is in the Lower Mississippi River watershed.[30] The entirety of the state is in the Mississippi River watershed, except for a small area in Bradley and Polk counties traversed by the Conasauga River, which is part of the Mobile Bay watershed.[31]
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The Blue Ridge area lies on the eastern edge of Tennessee, which borders North Carolina. This region of Tennessee is characterized by the high mountains and rugged terrain of the western Blue Ridge Mountains, which are subdivided into several subranges, namely the Great Smoky Mountains, the Bald Mountains, the Unicoi Mountains, the Unaka Mountains and Roan Highlands, and the Iron Mountains.
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The average elevation of the Blue Ridge area is 5,000 feet (1,500 m) above sea level. Clingmans Dome, the state's highest point, is located in this region. The Blue Ridge area was never more than sparsely populated, and today much of it is protected by the Cherokee National Forest, the Great Smoky Mountains National Park, and several federal wilderness areas and state parks.
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Stretching west from the Blue Ridge for approximately 55 miles (89 km) is the Ridge and Valley region, in which numerous tributaries join to form the Tennessee River in the Tennessee Valley. This area of Tennessee is covered by fertile valleys separated by wooded ridges, such as Bays Mountain and Clinch Mountain. The western section of the Tennessee Valley, where the depressions become broader and the ridges become lower, is called the Great Valley. In this valley are numerous towns and two of the region's three urban areas, Knoxville, the third largest city in the state, and Chattanooga, the fourth largest city in the state. The third urban area, the Tri-Cities, comprising Bristol, Johnson City, and Kingsport and their environs, is located to the northeast of Knoxville.
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The Cumberland Plateau rises to the west of the Tennessee Valley; this area is covered with flat-topped mountains separated by sharp valleys. The elevation of the Cumberland Plateau ranges from 1,500 to about 2,000 feet (460 to about 610 m) above sea level.
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East Tennessee has several important transportation links with Middle and West Tennessee, as well as the rest of the nation and the world, including several major airports and interstates. Knoxville's McGhee Tyson Airport (TYS) and Chattanooga's Chattanooga Metropolitan Airport (CHA), as well as the Tri-Cities' Tri-Cities Regional Airport (TRI), provide air service to numerous destinations. I-24, I-81, I-40, I-75, and I-26 along with numerous state highways and other important roads, traverse the Grand Division and connect Chattanooga, Knoxville, and the Tri-Cities, along with other cities and towns such as Cleveland, Athens, and Sevierville.
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West of the Cumberland Plateau is the Highland Rim, an elevated plain that surrounds the Nashville Basin. The northern section of the Highland Rim, known for its high tobacco production, is sometimes called the Pennyroyal Plateau; it is located primarily in Southwestern Kentucky. The Nashville Basin is characterized by rich, fertile farm country and great diversity of natural wildlife.
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Middle Tennessee was a common destination of settlers crossing the Appalachians from Virginia in the late 18th century and early 19th century. An important trading route called the Natchez Trace, created and used for many generations by American Indians, connected Middle Tennessee to the lower Mississippi River town of Natchez. The route of the Natchez Trace was used as the basis for a scenic highway called the Natchez Trace Parkway.
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Some of the last remaining large American chestnut trees grow in this region. They are being used to help breed blight-resistant trees.
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Middle Tennessee is one of the primary state population and transportation centers along with the heart of state government. Nashville (the capital), Clarksville, and Murfreesboro are its largest cities.[32] Interstates 24, 40, 65, and 840 service the Division, with the first three meeting in Nashville.
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West of the Highland Rim and Nashville Basin is the Gulf Coastal Plain, which includes the Mississippi embayment. The Gulf Coastal Plain is, in terms of area, the predominant land region in Tennessee. It is part of the large geographic land area that begins at the Gulf of Mexico and extends north into southern Illinois. In Tennessee, the Gulf Coastal Plain is divided into three sections that extend from the Tennessee River in the east to the Mississippi River in the west.
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The easternmost section, about 10 miles (16 km) in width, consists of hilly land that runs along the western bank of the Tennessee River. To the west of this narrow strip of land is a wide area of rolling hills and streams that stretches all the way to the Mississippi River; this area is called the Tennessee Bottoms or bottom land. In Memphis, the Tennessee Bottoms end in steep bluffs overlooking the river. To the west of the Tennessee Bottoms is the Mississippi Alluvial Plain, less than 300 feet (91 m) above sea level. This area of lowlands, flood plains, and swamp land is sometimes referred to as the Delta region. Memphis is the economic center of West Tennessee.
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Most of West Tennessee remained Indian land until the Chickasaw Cession of 1818, when the Chickasaw ceded their land between the Tennessee River and the Mississippi River. The portion of the Chickasaw Cession that lies in Kentucky is known today as the Jackson Purchase.
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Areas under the control and management of the National Park Service include the following:
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Fifty-four state parks, covering some 132,000 acres (530 km2) as well as parts of the Great Smoky Mountains National Park and Cherokee National Forest, and Cumberland Gap National Historical Park are in Tennessee. Sportsmen and visitors are attracted to Reelfoot Lake, originally formed by the 1811–12 New Madrid earthquakes; stumps and other remains of a once dense forest, together with the lotus bed covering the shallow waters, give the lake an eerie beauty.
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Most of the state has a humid subtropical climate, with the exception of some of the higher elevations in the Appalachians, which are classified as having a mountain temperate or humid continental climate due to cooler temperatures.[33] The Gulf of Mexico is the dominant factor in the climate of Tennessee, with winds from the south being responsible for most of the state's annual precipitation. Generally, the state has hot summers and mild to cool winters with generous precipitation throughout the year, with highest average monthly precipitation generally in the winter and spring months, between December and April. The driest months, on average, are August to October. On average the state receives 50 inches (130 cm) of precipitation annually. Snowfall ranges from 5 inches (13 cm) in West Tennessee to over 80 inches (200 cm) in the highest mountains in East Tennessee.[34][35]
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Summers in the state are generally hot and humid, with most of the state averaging a high of around 90 °F (32 °C) during the summer months. Winters tend to be mild to cool, increasing in coolness at higher elevations. Generally, for areas outside the highest mountains, the average overnight lows are near freezing for most of the state. The highest recorded temperature is 113 °F (45 °C) at Perryville on August 9, 1930, while the lowest recorded temperature is −32 °F (−36 °C) at Mountain City on December 30, 1917.
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While the state is far enough from the coast to avoid any direct impact from a hurricane, the location of the state makes it likely to be impacted from the remnants of tropical cyclones which weaken over land and can cause significant rainfall, such as Tropical Storm Chris in 1982 and Hurricane Opal in 1995.[36] The state averages about fifty days of thunderstorms per year, some of which can be severe with large hail and damaging winds. Tornadoes are possible throughout the state, with West and Middle Tennessee the most vulnerable. Occasionally, strong or violent tornadoes occur, such as the devastating April 2011 tornadoes that killed twenty people in North Georgia and Southeast Tennessee.[37] On average, the state has 15 tornadoes per year.[38] Tornadoes in Tennessee can be severe, and Tennessee leads the nation in the percentage of total tornadoes which have fatalities.[39] Winter storms are an occasional problem, such as the infamous Blizzard of 1993, although ice storms are a more likely occurrence. Fog is a persistent problem in parts of the state, especially in East Tennessee.
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The capital and largest city is Nashville, though Knoxville, Kingston, and Murfreesboro have all served as state capitals in the past. Nashville's 13-county metropolitan area has been the state's largest since c. 1990. Memphis was the largest city in the state until 2018 when it was surpassed by Nashville. Chattanooga and Knoxville, both in the east near the Great Smoky Mountains, are about a third that size. The city of Clarksville is a fifth significant population center, 45 miles (72 km) northwest of Nashville. Murfreesboro is the sixth-largest city in Tennessee, consisting of 146,900 residents.
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The area now known as Tennessee was first inhabited by Paleo-Indians nearly 12,000 years ago.[42] The names of the cultural groups who inhabited the area between first settlement and the time of European contact are unknown, but several distinct cultural phases have been named by archaeologists, including Archaic (8000–1000 BC), Woodland (1000 BC – 1000 AD), and Mississippian (1000–1600 AD), whose chiefdoms were the cultural predecessors of the Muscogee people who inhabited the Tennessee River Valley before Cherokee migration into the river's headwaters.
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The first recorded European excursions into what is now called Tennessee were three expeditions led by Spanish explorers, namely Hernando de Soto in 1540, Tristan de Luna in 1559, and Juan Pardo in 1567. Pardo recorded the name "Tanasqui" from a local Indian village, which evolved to the state's current name. At that time, Tennessee was inhabited by tribes of Muscogee and Yuchi people. Possibly because of European diseases devastating the Indian tribes, which would have left a population vacuum, and also from expanding European settlement in the north, the Cherokee moved south from the area now called Virginia. As European colonists spread into the area, the Indian populations were forcibly displaced to the south and west, including all Muscogee and Yuchi peoples, the Chickasaw and Choctaw, and ultimately, the Cherokee in 1838.
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The first British settlement in what is now Tennessee was built in 1756 by settlers from the colony of South Carolina at Fort Loudoun, near present-day Vonore. Fort Loudoun became the westernmost British outpost to that date. The fort was designed by John William Gerard de Brahm and constructed by forces under British Captain Raymond Demeré. After its completion, Captain Raymond Demeré relinquished command on August 14, 1757, to his brother, Captain Paul Demeré. Hostilities erupted between the British and the neighboring Overhill Cherokees, and a siege of Fort Loudoun ended with its surrender on August 7, 1760. The following morning, Captain Paul Demeré and a number of his men were killed in an ambush nearby, and most of the rest of the garrison was taken prisoner.[43]
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In the 1760s, long hunters from Virginia explored much of East and Middle Tennessee, and the first permanent European settlers began arriving late in the decade. The majority of 18th century settlers were English or of primarily English descent but nearly 20% of them were also Scotch-Irish.[44] These settlers formed the Watauga Association, a community built on lands leased from the Cherokee peoples.
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During the American Revolutionary War, Fort Watauga at Sycamore Shoals (in present-day Elizabethton) was attacked (1776) by Dragging Canoe and his warring faction of Cherokee who were aligned with the British Loyalists. These renegade Cherokee were referred to by settlers as the Chickamauga. They opposed North Carolina's annexation of the Washington District and the concurrent settling of the Transylvania Colony further north and west. The lives of many settlers were spared from the initial warrior attacks through the warnings of Dragging Canoe's cousin, Nancy Ward. The frontier fort on the banks of the Watauga River later served as a 1780 staging area for the Overmountain Men in preparation to trek over the Appalachian Mountains, to engage, and to later defeat the British Army at the Battle of Kings Mountain in South Carolina.
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Three counties of the Washington District (now part of Tennessee) broke off from North Carolina in 1784 and formed the State of Franklin. Efforts to obtain admission to the Union failed, and the counties (now numbering eight) had re-joined North Carolina by 1789. North Carolina ceded the area to the federal government in 1790, after which it was organized into the Southwest Territory. In an effort to encourage settlers to move west into the new territory, in 1787 the mother state of North Carolina ordered a road to be cut to take settlers into the Cumberland Settlements—from the south end of Clinch Mountain (in East Tennessee) to French Lick (Nashville). The Trace was called the "North Carolina Road" or "Avery's Trace", and sometimes "The Wilderness Road" (although it should not be confused with Daniel Boone's "Wilderness Road" through the Cumberland Gap).
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Tennessee was admitted to the Union on June 1, 1796 as the 16th state. It was the first state created from territory under the jurisdiction of the United States federal government. Apart from the former Thirteen Colonies only Vermont and Kentucky predate Tennessee's statehood, and neither was ever a federal territory.[45] The Constitution of the State of Tennessee, Article I, Section 31, states that the beginning point for identifying the boundary is the extreme height of the Stone Mountain, at the place where the line of Virginia intersects it, and basically runs the extreme heights of mountain chains through the Appalachian Mountains separating North Carolina from Tennessee past the Native American towns of Cowee and Old Chota, thence along the main ridge of the said mountain (Unicoi Mountain) to the southern boundary of the state; all the territory, lands and waters lying west of said line are included in the boundaries and limits of the newly formed state of Tennessee. Part of the provision also stated that the limits and jurisdiction of the state would include future land acquisition, referencing possible land trade with other states, or the acquisition of territory from west of the Mississippi River.
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As more white settlers moved into Tennessee, they came into increasing conflict with Native American tribes. During the administration of U.S. President Martin Van Buren, nearly 17,000 Cherokees, along with approximately 2,000 black slaves owned by Cherokees, were uprooted from their homes between 1838 and 1839 and were forced by the U.S. military to march from "emigration depots" in Eastern Tennessee, such as Fort Cass, toward the more distant Indian Territory west of Arkansas, now the state of Oklahoma.[46] During this relocation an estimated 4,000 Cherokees died along the way west.[47] In the Cherokee language, the event is called Nunna daul Isunyi,"the Trail Where We Cried". The Cherokees were not the only Native Americans forced to emigrate as a result of the Indian removal efforts of the United States, and so the phrase "Trail of Tears" is sometimes used to refer to similar events endured by other American Indian peoples, especially among the "Five Civilized Tribes". The phrase originated as a description of the earlier emigration of the Choctaw nation.
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In February 1861, secessionists in Tennessee's state government—led by Governor Isham Harris—sought voter approval for a convention to sever ties with the United States, but Tennessee voters rejected the referendum by a 54–46% margin. The strongest opposition to secession came from East Tennessee, which later tried to form a separate Union-aligned state. Following the Confederate attack upon Fort Sumter in April and Lincoln's call for troops from Tennessee and other states in response, Governor Isham Harris began military mobilization, submitted an ordinance of secession to the General Assembly, and made direct overtures to the Confederate government. The Tennessee legislature ratified an agreement to enter a military league with the Confederate States on May 7, 1861. On June 8, 1861, with people in Middle Tennessee having significantly changed their position, voters approved a second referendum calling for secession, becoming the last state to do so. But the Union-backing State of Scott was also established at this time, and remained a de facto enclave of the United States throughout the war.
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Many major battles of the American Civil War were fought in Tennessee—most of them Union victories. Ulysses S. Grant and the U.S. Navy captured control of the Cumberland and Tennessee rivers in February 1862. They held off the Confederate counterattack at Shiloh in April. Memphis fell to the Union in June, following a naval battle on the Mississippi River in front of the city. The Capture of Memphis and Nashville gave the Union control of the western and middle sections. This control was confirmed at the Battle of Murfreesboro in early January 1863 and by the subsequent Tullahoma Campaign.
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Despite the strength of Unionist sentiment in East Tennessee (with the exception of Sullivan County, which was heavily pro-Confederate), Confederates held the area. The Confederates, led by General James Longstreet, did attack General Burnside's Fort Sanders at Knoxville and lost. It was a big blow to East Tennessee Confederate momentum, but Longstreet won the Battle of Bean's Station a few weeks later. The Confederates besieged Chattanooga during the Chattanooga Campaign in early fall 1863, but were driven off by Grant in November. Many of the Confederate defeats can be attributed to the poor strategic vision of General Braxton Bragg, who led the Army of Tennessee from Perryville, Kentucky, to another Confederate defeat at Chattanooga.
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The last major battles came when the Confederates invaded Middle Tennessee in November 1864 and were checked at Franklin, then completely dispersed by George Thomas at Nashville in December. Meanwhile, President Abraham Lincoln appointed displaced senator and native Tennessean Andrew Johnson military governor of the state.
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When the Emancipation Proclamation was announced, Tennessee was largely held by Union forces. Thus it was not among the states enumerated in the Proclamation, and the Proclamation did not free any slaves there. Nonetheless, enslaved African Americans escaped to Union lines to gain freedom without waiting for official action. Old and young, men, women, and children camped near Union troops. Thousands of former slaves ended up fighting on the Union side, nearly 200,000 in total across the South.
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Tennessee's legislature approved an amendment to the state constitution prohibiting slavery on February 22, 1865.[48] Voters in the state approved the amendment in March.[49] It also ratified the Thirteenth Amendment to the United States Constitution (abolishing slavery in every state) on April 7, 1865.
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In 1864, Andrew Johnson, a War Democrat from Tennessee, was elected Vice President under Abraham Lincoln. He became President after Lincoln's assassination in 1865. Under Johnson's lenient re-admission policy, Tennessee was the first of the seceding states to have its elected members readmitted to the U.S. Congress, on July 24, 1866. Because Tennessee had ratified the Fourteenth Amendment, it was the only one of the formerly secessionist states that did not have a military governor during the Reconstruction period.
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After the formal end of Reconstruction, the struggle over power in Southern society continued. Through violence and intimidation against freedmen and their allies, White Democrats regained political power in Tennessee and other states across the South in the late 1870s and 1880s. Over the next decade, the state legislature passed increasingly restrictive laws to control African Americans. In 1889 the General Assembly passed four laws described as electoral reform, with the cumulative effect of essentially disfranchising most African Americans in rural areas and small towns, as well as many poor Whites. Legislation included implementation of a poll tax, timing of registration, and recording requirements. Tens of thousands of taxpaying citizens were without representation for decades into the 20th century.[13] Disfranchising legislation accompanied Jim Crow laws passed in the late 19th century, which imposed segregation in the state. In 1900, African Americans made up nearly 24% of the state's population, and numbered 480,430 citizens who lived mostly in the central and western parts of the state.[50]
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In 1897, Tennessee celebrated its centennial of statehood (though one year late of the 1896 anniversary) with a great exposition in Nashville. A full-scale replica of the Parthenon was constructed for the celebration, located in what is now Nashville's Centennial Park.
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On August 18, 1920, Tennessee became the thirty-sixth and final state necessary to ratify the Nineteenth Amendment to the United States Constitution, which provided women the right to vote. Disenfranchising voter registration requirements continued to keep most African Americans and many poor whites, both men and women, off the voter rolls.
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The need to create work for the unemployed during the Great Depression, a desire for rural electrification, the need to control annual spring flooding and improve shipping capacity on the Tennessee River were all factors that drove the federal creation of the Tennessee Valley Authority (TVA) in 1933. Through the power of the TVA projects, Tennessee quickly became the nation's largest public utility supplier.
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During World War II, the availability of abundant TVA electrical power led the Manhattan Project to locate one of the principal sites for production and isolation of weapons-grade fissile material in East Tennessee. The planned community of Oak Ridge was built from scratch to provide accommodations for the facilities and workers. These sites are now Oak Ridge National Laboratory, the Y-12 National Security Complex, and the East Tennessee Technology Park.
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Despite recognized effects of limiting voting by poor whites, successive legislatures expanded the reach of the disfranchising laws until they covered the state. Political scientist V. O. Key, Jr. argued in 1949:
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... the size of the poll tax did not inhibit voting as much as the inconvenience of paying it. County officers regulated the vote by providing opportunities to pay the tax (as they did in Knoxville), or conversely by making payment as difficult as possible. Such manipulation of the tax, and therefore the vote, created an opportunity for the rise of urban bosses and political machines. Urban politicians bought large blocks of poll tax receipts and distributed them to blacks and whites, who then voted as instructed.[13]
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In 1953 state legislators amended the state constitution, removing the poll tax. In many areas both blacks and poor whites still faced subjectively applied barriers to voter registration that did not end until after passage of national civil rights legislation, including the Voting Rights Act of 1965.[13]
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Tennessee celebrated its bicentennial in 1996. With a yearlong statewide celebration entitled "Tennessee 200", it opened a new state park (Bicentennial Mall) at the foot of Capitol Hill in Nashville.
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The state has had major disasters, such as the Great Train Wreck of 1918, one of the worst train accidents in U.S. history,[51] and the Sultana explosion on the Mississippi River near Memphis, the deadliest maritime disaster in U.S. history.[52]
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In 2002, businessman Phil Bredesen was elected to become the 48th governor in January 2003. Also in 2002, Tennessee amended the state constitution to allow for the establishment of a lottery. Tennessee's Bob Corker was the only freshman Republican elected to the United States Senate in the 2006 midterm elections. The state constitution was amended to reject same-sex marriage. In January 2007, Ron Ramsey became the first Republican elected as Speaker of the State Senate since Reconstruction,[citation needed] as a result of the realignment of the Democratic and Republican parties in the South since the late 20th century, with Republicans now elected by conservative voters, who previously had supported Democrats.
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In 2010, during the 2010 midterm elections, Bill Haslam was elected to succeed Bredesen, who was term-limited, to become the 49th Governor of Tennessee in January 2011. In April and May 2010, flooding in Middle Tennessee devastated Nashville and other parts of Middle Tennessee. In 2011, parts of East Tennessee, including Hamilton and Bradley counties, were devastated by the April 2011 tornado outbreak.
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In 2018, Republican businessman Bill Lee was elected to succeed term-limited Haslam, and he became the 50th Governor of Tennessee.
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The United States Census Bureau estimates that the population of Tennessee was 6,829,174 on July 1, 2019, an increase of 483,069 people since the 2010 United States Census, or 7.61%.[4] This includes a natural increase since the last census of 124,385 (584,236 births minus 459,851 deaths), and an increase from net migration of 244,537 people into the state. Immigration from outside the United States resulted in a net increase of 66,412, and migration within the country produced a net increase of 178,125.[54]
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Twenty percent of Tennesseans were born outside the South in 2008, compared to a figure of 13.5% in 1990.[55] In recent years, Tennessee has received an influx of people relocating from California, Florida, New York, New Jersey and the New England States for the low cost of living, and the booming healthcare and automobile industries. Metropolitan Nashville is one of the fastest-growing areas in the country due in part to these factors.[56]
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The center of population of Tennessee is located in Rutherford County, in the city of Murfreesboro.[57]
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In 2010, 4.6% of the total population was of Hispanic or Latino origin (they may be of any race), up from 2.2% in 2000. Between 2000 and 2010, the Hispanic population in Tennessee grew by 134.2%, the third highest of any state.[61] That same year Non-Hispanic whites were 75.6% of the population, compared to 63.7% of the population nationwide.[62]
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In 2010, approximately 4.4% of Tennessee's population was foreign-born, an increase of about 118.5% since 2000. Of the foreign-born population, approximately 31.0% were naturalized citizens and 69.0% non-citizens. The foreign-born population was approximately 49.9% from Latin America, 27.1% from Asia, 11.9% from Europe, 7.7% from Africa, 2.7% from Northern America, and 0.6% from Oceania.[63]
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In 2010, the five most common self-reported ethnic groups in the state were: American (26.5%), English (8.2%), Irish (6.6%), German (5.5%), and Scotch-Irish (2.7%).[64] Most Tennesseans who self-identify as having American ancestry are of English and Scotch-Irish ancestry. An estimated 21–24% of Tennesseans are of predominantly English ancestry.[65][66] In the 1980 census 1,435,147 Tennesseans claimed "English" or "mostly English" ancestry out of a state population of 3,221,354 making them 45% of the state at the time.[67]
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According to the 2010 census, 6.4% of Tennessee's population were reported as under age 5, 23.6% under 18, and 13.4% were 65 or older.[68]
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On June 19, 2010, the Tennessee Commission of Indian Affairs granted state recognition to six Native American tribes, which was later repealed by the state's Attorney General because the action by the commission was illegal. The tribes were as follows:
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As of 2011, 36.3% of Tennessee's population younger than age 1 were minorities.[70]
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Note: Births in table do not add up, because Hispanics are counted both by their ethnicity and by their race, giving a higher overall number.
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Religious affiliation of Tennessee population as of 2014[77]
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The religious affiliations of the people of Tennessee as of 2014:[77]
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The largest denominations by number of adherents in 2010 were the Southern Baptist Convention with 1,483,356; the United Methodist Church with 375,693; the Roman Catholic Church with 222,343; and the Churches of Christ with 214,118.[78]
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As of January 1, 2009, The Church of Jesus Christ of Latter-day Saints (LDS Church) reported 43,179 members, 10 stakes, 92 Congregations (68 wards and 24 branches), two missions, and two temples in Tennessee.[79]
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Tennessee is home to several Protestant denominations, such as the National Baptist Convention (headquartered in Nashville); the Church of God in Christ and the Cumberland Presbyterian Church (both headquartered in Memphis); the Church of God and The Church of God of Prophecy (both headquartered in Cleveland). The Free Will Baptist denomination is headquartered in Antioch; its main Bible college is in Nashville. The Southern Baptist Convention maintains its general headquarters in Nashville. Publishing houses of several denominations are located in Nashville.
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According to the U.S. Bureau of Economic Analysis, in 2011 Tennessee's real gross state product was $233.997 billion.
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In 2003, the per capita personal income was $28,641, 36th in the nation, and 91% of the national per capita personal income of $31,472. In 2004, the median household income was $38,550, 41st in the nation, and 87% of the national median of $44,472.
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For 2012, the state held an asset surplus of $533 million, one of only eight states in the nation to report a surplus.[81]
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Major outputs for the state include textiles, cotton, cattle, and electrical power. Tennessee has more than 82,000 farms, roughly 59 percent of which accommodate beef cattle.[82] Although cotton was an early crop in Tennessee, large-scale cultivation of the fiber did not begin until the 1820s with the opening of the land between the Tennessee and Mississippi Rivers. The upper wedge of the Mississippi Delta extends into southwestern Tennessee, and it was in this fertile section that cotton took hold. Soybeans are also heavily planted in West Tennessee, focusing on the northwest corner of the state.[83]
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Large corporations with headquarters in Tennessee include FedEx, AutoZone and International Paper, all based in Memphis; Pilot Corporation and Regal Entertainment Group, based in Knoxville; Eastman Chemical Company, based in Kingsport; Hospital Corporation of America and Caterpillar Financial, based in Nashville; and Unum, based in Chattanooga.
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Tennessee is also a major hub for the automotive industry.[84] Three automobile manufacturers have factories in Tennessee: Nissan in Smyrna, General Motors in Spring Hill, and Volkswagen in Chattanooga.[85] Nissan moved its North American corporate headquarters from California to Franklin, Tennessee in 2005,[86] and Mitsubishi Motors did the same in 2019.[84]
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Other major manufacturers include a $2 billion polysilicon production facility owned by Wacker Chemie in Bradley County and a $1.2 billion polysilicon production facility owned by Hemlock Semiconductor in Clarksville.
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Tennessee is a right to work state, as are most of its Southern neighbors. Unionization has historically been low and continues to decline as in most of the U.S. generally. As of August 2019, the state has an unemployment rate of 3.5%, which is ranked 28th in the country.[87] As of 2015, 16.7% of the population of Tennessee lives below the poverty line, which is higher than the national average of 14.7%.[88]
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Tennessee has a reputation as low-tax state and is usually ranked as one of the five states with the lowest tax burden on residents.[89] It is one of nine states that do not have a general income tax; the sales tax is the primary means of funding the government.[90] The Hall income tax is a tax imposed on most dividends and interest. The tax rate was 6% from 1937 to 2016, but is in the process of being completely phased out by January 1, 2021.[91] The first $1,250 of individual income and $2,500 of joint income is exempt from this tax.[92]
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The state's sales and use tax rate for most items is 7%, the second-highest in the nation, along with Mississippi, Rhode Island, New Jersey, and Indiana. Food is taxed at a lower rate of 4%, but candy, dietary supplements and prepared food are taxed at 7%.[93] Local sales taxes are collected in most jurisdictions at rates varying from 1.5% to 2.75%, bringing the total sales tax to between 8.5% and 9.75%, with an average rate of about 9.5%, the nation's highest average sales tax.[94] Intangible property tax is assessed on the shares of stock of stockholders of any loan, investment, insurance, or for-profit cemetery companies. The assessment ratio is 40% of the value times the jurisdiction's tax rate.[95] Since January 1, 2016, Tennessee has had no inheritance tax.[96]
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While sales tax remains the main source of state government funding, property taxes are the primary source of revenue for local governments.[95]
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Tourism contributes billions of dollars every year to the state's economy, and Tennessee is ranked among the Top 10 destinations in the nation.[97] In 2014 a record 100 million people visited the state resulting in $17.7 billion in tourism related spending within the state, an increase of 6.3% over 2013; tax revenue from tourism equaled $1.5 billion. Each county in Tennessee saw at least $1 million from tourism while 19 counties received at least $100 million (Davidson, Shelby, and Sevier counties were the top three). Tourism-generated jobs for the state reached 152,900, a 2.8% increase.[97] International travelers to Tennessee accounted for $533 million in spending.[98]
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In 2013 tourism within the state from local citizens accounted for 39.9% of tourists, the second highest originating location for tourists to Tennessee is the state of Georgia, accounting for 8.4% of tourists.[99]:17 Forty-four percent of stays in the state were "day trips", 25% stayed one night, 15% stayed two nights, and 11% stayed four or more nights. The average stay was 2.16 nights, compared to 2.03 nights for the U.S. as a whole.[99]:40 The average person spent $118 per day: 29% on transportation, 24% on food, 17% on accommodation, and 28% on shopping and entertainment.[99]:44
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Tennessee is home to the Great Smoky Mountains National Park, the most visited national park in the United States. The park anchors a large tourism industry based primarily in nearby Gatlinburg and Pigeon Forge, which consists of such attractions as Dollywood, the most visited ticketed attraction in Tennessee, Ober Gatlinburg, and Ripley's Aquarium of the Smokies.[100] Major attractions in Memphis include Graceland, the home of Elvis Presley, Beale Street, the National Civil Rights Museum, the Memphis Zoo, and the Stax Museum of American Soul Music.[101] Nashville contains many attractions related to its musical heritage, including Lower Broadway, the Country Music Hall of Fame, the Ryman Auditorium, Grand Ole Opry, and the Gaylord Opryland Resort. Other major attractions in Nashville include the Tennessee State Museum, Parthenon, and the Belle Meade Plantation.[102] Major attractions in Chattanooga include Lookout Mountain, the Chattanooga Choo-Choo Hotel, Ruby Falls, and the Tennessee Aquarium, the largest freshwater aquarium in the United States.[100] Other major attractions include the American Museum of Science and Energy in Oak Ridge, the Bristol Motor Speedway in Bristol, Jack Daniel's Distillery in Lynchburg, and the Hiwassee and Ocoee rivers in Polk County.[100]
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Four Civil War battlefields in Tennessee are preserved by the National Park Service: Chickamauga and Chattanooga National Military Park, Stones River National Battlefield, Shiloh National Military Park, and Fort Donelson National Battlefield.[103] Big South Fork National River and Recreation Area is within the Cumberland Mountains in northeastern Tennessee. Other major attractions preserved by the National Park Service include Cumberland Gap National Historical Park and Overmountain Victory National Historic Trail.[104]
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Tennessee's electric utilities are regulated monopolies, as in many other states.[105][106] As of 2019, the Tennessee Valley Authority owned over 90% of generating capacity.[107] Nuclear power is the largest source of electricity generation in Tennessee, producing about 43.7% of its power in 2019. The same year, 22.9% of the power was produced by coal, 20.3% from natural gas, 12.1% from hydroelectric power, and 1.6% from other renewables. About 56% of the electricity generated in Tennessee produces no greenhouse gas emissions.[108] Tennessee is a net consumer of electricity, receiving power from other TVA facilities in neighboring states such as the Browns Ferry Nuclear Plant in northern Alabama.[109]
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Tennessee is home to the two newest civilian nuclear power reactors in the United States, at Watts Bar Nuclear Plant in Rhea County. Unit 1 began operation in 1996 and Unit 2 began operation in 2016, making it the first and only new nuclear power reactor to begin operation in the United States in the 21st century.[110] Tennessee was also an early leader in hydroelectric power, first with the now defunct Chattanooga and Tennessee Electric Power Company; later, the United States Army Corps of Engineers and the TVA constructed several hydroelectric dams on Tennessee rivers.[111] Today, Tennessee is the third-largest hydroelectric power-producing state east of the Rocky Mountains.[112]
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Tennessee has very little petroleum and natural gas reserves, but is home to one oil refinery, in Memphis.[112] Bituminous coal is mined in small quantities in the Cumberland Plateau and Cumberland Mountains.[113] There are sizable reserves of lignite coal in West Tennessee that remain untapped.[113] Coal production in Tennessee peaked in 1972, and today less than 0.1% of coal production in the United States comes from Tennessee mines.[112]
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Tennessee is the leading producer of ball clay in the United States.[113] Other major mineral products produced in Tennessee include sand, gravel, crushed stone, Portland cement, marble, sandstone, common clay, lime, and zinc.[113][114] The Copper Basin, in Tennessee's southeastern corner in Polk County, was one of the most productive copper mining districts in the United States between the 1840s and 1980s.[115] Mines in the basin supplied about 90% of the copper used by the Confederacy during the Civil War,[116] and also marketed chemical byproducts of the mining, including sulfuric acid.[117] Mining activities in the basin resulted in a major environmental disaster, which left the landscape in the basin barren for more than a century.[118] Iron ore was another major mineral mined in Tennessee until the early 20th century.[119] Tennessee was also a top producer of phosphate until the early 1990s.[120]
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Tennessee has played a critical role in the development of many forms of American popular music, including rock and roll, blues, country, and rockabilly. Beale Street in Memphis is considered by many to be the birthplace of the blues, with musicians such as W. C. Handy performing in its clubs as early as 1909.[121] Memphis is also home to Sun Records, where musicians such as Elvis Presley, Johnny Cash, Carl Perkins, Jerry Lee Lewis, Roy Orbison, and Charlie Rich began their recording careers, and where rock and roll took shape in the 1950s.[122] The 1927 Victor recording sessions in Bristol generally mark the beginning of the country music genre and the rise of the Grand Ole Opry in the 1930s helped make Nashville the center of the country music recording industry.[123][124] Three brick-and-mortar museums recognize Tennessee's role in nurturing various forms of popular music: the Memphis Rock N' Soul Museum, the Country Music Hall of Fame and Museum in Nashville, and the International Rock-A-Billy Museum in Jackson. Moreover, the Rockabilly Hall of Fame, an online site recognizing the development of rockabilly in which Tennessee played a crucial role, is based in Nashville.
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Notable writers with ties to Tennessee include Cormac McCarthy, Peter Taylor, James Agee, Francis Hodgson Burnett, Thomas S. Stribling, Ida B. Wells, Nikki Giovanni, Shelby Foote, Ann Patchett, Ishmael Reed, and Randall Jarrell.
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Tennessee is home to four major professional sports franchises: the Tennessee Titans have played in the National Football League since 1997, the Memphis Grizzlies have played in the National Basketball Association since 2001, the Nashville Predators have played in the National Hockey League since 1998, and Nashville SC has played in Major League Soccer since 2020.
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The state is also home to 13 teams playing in minor leagues. Nine Minor League Baseball teams call the state their home. The Memphis Redbirds and Nashville Sounds, each of the Pacific Coast League, compete at the Triple-A level, the highest before Major League Baseball. The Chattanooga Lookouts, Jackson Generals, and Tennessee Smokies play in the Double-A Southern League. The Elizabethton Twins, Greeneville Reds, Johnson City Cardinals, and Kingsport Mets are Rookie League teams of the Appalachian League.
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The Knoxville Ice Bears are a minor league ice hockey team of the Southern Professional Hockey League. Memphis 901 FC, a soccer team, joined the USL Championship in 2019.[125] Chattanooga Red Wolves SC became an inaugural member of the third-level USL League One in 2019. Chattanooga FC began play in the National Independent Soccer Association in 2020.
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In Knoxville, the Tennessee Volunteers college team has played in the Southeastern Conference (SEC) of the National Collegiate Athletic Association since the conference was formed in 1932. The football team has won 13 SEC championships and 28 bowls, including four Sugar Bowls, three Cotton Bowls, an Orange Bowl and a Fiesta Bowl. Meanwhile, the men's basketball team has won four SEC championships and reached the NCAA Elite Eight in 2010. In addition, the women's basketball team has won a host of SEC regular-season and tournament titles along with eight national titles.
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In Nashville, the Vanderbilt Commodores are also charter members of the SEC. The Tennessee–Vanderbilt football rivalry began in 1892, having since played more than a hundred times. In June 2014, the Vanderbilt Commodores baseball team won its first men's national championship by winning the 2014 College World Series.
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The state is home to 10 other NCAA Division I programs. Two of these participate in the top level of college football, the Football Bowl Subdivision. The Memphis Tigers are members of the American Athletic Conference, and the Middle Tennessee Blue Raiders from Murfreesboro play in Conference USA. In addition to the Commodores, Nashville is also home to the Belmont Bruins and Tennessee State Tigers, both members of the Ohio Valley Conference (OVC), and the Lipscomb Bisons, members of the Atlantic Sun Conference. Tennessee State plays football in Division I's second level, the Football Championship Subdivision (FCS), while Belmont and Lipscomb do not have football teams. Belmont and Lipscomb have an intense rivalry in men's and women's basketball known as the Battle of the Boulevard, with both schools' men's and women's teams playing two games each season against each other (a rare feature among non-conference rivalries). The OVC also includes the Austin Peay Governors from Clarksville, the UT Martin Skyhawks from Martin, and the Tennessee Tech Golden Eagles from Cookeville. These three schools, along with fellow OVC member Tennessee State, play each season in football for the Sgt. York Trophy. The Chattanooga Mocs and Johnson City's East Tennessee State Buccaneers are full members, including football, of the Southern Conference.
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Tennessee is also home to Bristol Motor Speedway which features NASCAR Cup Series racing two weekends a year, routinely selling out more than 160,000 seats on each date; it also was the home of the Nashville Superspeedway, which held Nationwide and IndyCar races until it was shut down in 2012. Tennessee's only graded stakes horse race, the Iroquois Steeplechase, is also held in Nashville each May.
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The FedEx St. Jude Classic is a PGA Tour golf tournament held at Memphis since 1958. The U.S. National Indoor Tennis Championships has been held at Memphis since 1976 (men's) and 2002 (women's).
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Interstate 40 crosses the state in an east–west orientation. Its branch interstate highways include I-240 in Memphis; I-440 in Nashville; I-840 in Nashville; I-140 from Knoxville to Alcoa; and I-640 in Knoxville. I-26, although technically an east–west interstate, runs from the North Carolina border below Johnson City to its terminus at Kingsport. I-24 is an east–west interstate that runs from Chattanooga to Clarksville. In a north–south orientation are highways I-55, I-65, I-75, and I-81. Interstate 65 crosses the state through Nashville, while Interstate 75 serves Chattanooga and Knoxville and Interstate 55 serves Memphis. Interstate 81 enters the state at Bristol and terminates at its junction with I-40 near Dandridge. I-155 is a branch highway from I-55. The only spur highway of I-75 in Tennessee is I-275, which is in Knoxville.[126] An extension of I-69 is proposed to travel through the western part of the state, from South Fulton to Memphis.[127] A branch interstate, I-269 also exists from Millington to Collierville.[126]
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Major airports within the state include Memphis International Airport (MEM), Nashville International Airport (BNA), McGhee Tyson Airport (TYS) outside of Knoxville in Blount County, Chattanooga Metropolitan Airport (CHA), Tri-Cities Regional Airport (TRI), and McKellar-Sipes Regional Airport (MKL), in Jackson. Because Memphis International Airport is the major hub for FedEx Corporation, it is the world's largest air cargo operation.
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For passenger rail service, Memphis and Newbern, are served by the Amtrak City of New Orleans line on its run between Chicago, Illinois, and New Orleans, Louisiana. Nashville is served by the Music City Star commuter rail service.
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Cargo services in Tennessee are primarily served by CSX Transportation, which has a hump yard in Nashville called Radnor Yard. Norfolk Southern Railway operates lines in East Tennessee, through cities including Knoxville and Chattanooga, and operates a classification yard near Knoxville, the John Sevier Yard. BNSF operates a major intermodal facility in Memphis.
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Similar to the United States Federal Government, Tennessee's government has three parts. The Executive Branch is led by Tennessee's governor, who holds office for a four-year term and may serve a maximum of two consecutive terms. The governor is the only official who is elected statewide. Unlike most states, the state does not elect the lieutenant governor directly; the Tennessee Senate elects its Speaker, who serves as lieutenant governor. The governor is supported by 22 cabinet-level departments, most headed by a commissioner who serves at the pleasure of the governor:
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The Executive Branch also includes several agencies, boards and commissions, some of which are under the auspices of one of the cabinet-level departments.[128]
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The bicameral Legislative Branch, known as the Tennessee General Assembly, consists of the 33-member Senate and the 99-member House of Representatives. Senators serve four-year terms, and House members serve two-year terms. Each chamber chooses its own speaker. The speaker of the state Senate also holds the title of lieutenant governor. Constitutional officials in the legislative branch are elected by a joint session of the legislature.
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The highest court in Tennessee is the state Supreme Court. It has a chief justice and four associate justices. No more than two justices can be from the same Grand Division. The Supreme Court of Tennessee also appoints the Attorney General, a practice not found in any of the other 49 states. Both the Court of Appeals and the Court of Criminal Appeals have 12 judges.[129] A number of local, circuit, and federal courts provide judicial services.
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Tennessee's current state constitution was adopted in 1870. The state had two earlier constitutions. The first was adopted in 1796, the year Tennessee joined the union, and the second was adopted in 1834. The 1870 Constitution outlaws martial law within its jurisdiction. This may be a result of the experience of Tennessee residents and other Southerners during the period of military control by Union (Northern) forces of the U.S. government after the American Civil War.
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Tennessee politics, like that of most U.S. states, are dominated by the Republican and the Democratic parties. Historian Dewey W. Grantham traces divisions in the state to the period of the American Civil War; for decades afterward, the eastern third of the state was heavily Republican and the western two thirds mostly voted Democratic.[130] This division was related to the state's pattern of farming, plantations and slaveholding. The eastern section was made up of yeoman farmers, but Middle and West Tennessee farmers cultivated crops such as tobacco and cotton which were dependent on the use of slave labor. These areas became defined as Democratic after the war.
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During Reconstruction, freedmen and former free people of color were granted the right to vote; most joined the Republican Party. Numerous African Americans were elected to local offices, and some to state office. Following Reconstruction, Tennessee continued to have competitive party politics, but in the 1880s, the white-dominated state government passed four laws, the last of which imposed a poll tax requirement for voter registration. These served to disenfranchise most African Americans, and their power in state and local politics was markedly reduced. In 1900 African Americans comprised 23.8 percent of the state's population, concentrated in Middle and West Tennessee.[50] In the early 1900s, the state legislature approved a form of commission government for cities based on at-large voting for a few positions on a Board of Commission; several cities adopted this as another means to limit African-American political participation. In 1913 the state legislature enacted a bill enabling cities to adopt this structure without legislative approval.[131]
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After disenfranchisement of blacks, the Republican Party in Tennessee became a primarily white sectional party supported only in the eastern part of the state. In the 20th century, except for two nationwide Republican landslides of the 1920s (in 1920, when Tennessee narrowly supported Warren G. Harding over Ohio Governor James Cox, and in 1928, when it more decisively voted for Herbert Hoover over New York Governor Al Smith), the state was part of the Democratic Solid South until the 1950s. In that postwar decade, it twice voted for Republican Dwight D. Eisenhower, former Allied Commander of the Armed Forces during World War II. Since then, more of the state's voters have shifted to supporting Republicans, and Democratic presidential candidates have carried Tennessee only four times.
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By 1960 African Americans comprised 16.45% of the state's population. It was not until after the mid-1960s and passage of the Voting Rights Act of 1965 that they were able to vote in full again, but new devices, such as at-large commission city governments, had been adopted in several jurisdictions to limit their political participation. Former Gov. Winfield Dunn and former U.S. Sen. Bill Brock wins in 1970 helped make the Republican Party competitive among whites for the statewide victory. Tennessee has selected governors from different parties since 1970.
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In the early 21st century, Republican voters control most of the state, especially in the more rural and suburban areas outside of the cities; Democratic strength is mostly confined to the urban cores of the four major cities, and is particularly strong in the cities of Nashville and Memphis. The latter area includes a large African-American population.[132] Historically, Republicans had their greatest strength in East Tennessee before the 1960s. Tennessee's 1st and 2nd congressional districts, based in the Tri-Cities and Knoxville, respectively, are among the few historically Republican districts in the South. Those districts' residents supported the Union over the Confederacy during the Civil War; they identified with the GOP after the war and have stayed with that party ever since. The first has been in Republican hands continuously since 1881, and Republicans (or their antecedents) have held it for all but four years since 1859. The second has been held continuously by Republicans or their antecedents since 1859.
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In the 2000 presidential election, Vice President Al Gore, a Democratic U.S. Senator from Tennessee, failed to carry his home state, an unusual occurrence but indicative of strengthening Republican support. Republican George W. Bush received increased support in 2004, with his margin of victory in the state increasing from 4% in 2000 to 14% in 2004.[133] Democratic presidential nominees from Southern states, such as Lyndon B. Johnson, Jimmy Carter, and Bill Clinton, usually fare better than their Northern counterparts do in Tennessee, especially among split-ticket voters outside the metropolitan areas.
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Tennessee sends nine members to the U.S. House of Representatives, of whom there are seven Republicans and two Democrats. Lieutenant Governor Ron Ramsey is the first Republican speaker of the state Senate in 140 years.[citation needed] In the 2008 elections, the Republican party gained control of both houses of the Tennessee state legislature for the first time since Reconstruction. In 2008, some 30% of the state's electorate identified as independents.[134]
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The Baker v. Carr (1962) decision of the U.S. Supreme Court established the principle of "one man, one vote", requiring state legislatures to redistrict to bring Congressional apportionment in line with decennial censuses. It also required both houses of state legislatures to be based on population for representation and not geographic districts such as counties. This case arose out of a lawsuit challenging the longstanding rural bias of apportionment of seats in the Tennessee legislature.[135][136][137] After decades in which urban populations had been underrepresented in many state legislatures, this significant ruling led to an increased (and proportional) prominence in state politics by urban and, eventually, suburban, legislators and statewide officeholders in relation to their population within the state. The ruling also applied to numerous other states long controlled by rural minorities, such as Alabama, Vermont, and Montana.
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The state of Tennessee maintains four dedicated law enforcement entities: the Tennessee Highway Patrol, the Tennessee Wildlife Resources Agency (TWRA), the Tennessee Bureau of Investigation (TBI), and the Tennessee Department of Environment and Conservation (TDEC).
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The Highway Patrol is the primary law enforcement entity that concentrates on highway safety regulations and general non-wildlife state law enforcement and is under the jurisdiction of the Tennessee Department of Safety. The TWRA is an independent agency tasked with enforcing all wildlife, boating, and fisheries regulations outside of state parks. The TBI maintains state-of-the-art investigative facilities and is the primary state-level criminal investigative department. Tennessee State Park Rangers are responsible for all activities and law enforcement inside the Tennessee State Parks system.
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Local law enforcement is divided between County Sheriff's Offices and Municipal Police Departments. Tennessee's Constitution requires that each County have an elected Sheriff. In 94 of the 95 counties the Sheriff is the chief law enforcement officer in the county and has jurisdiction over the county as a whole. Each Sheriff's Office is responsible for warrant service, court security, jail operations and primary law enforcement in the unincorporated areas of a county as well as providing support to the municipal police departments. Incorporated municipalities are required to maintain a police department to provide police services within their corporate limits.
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The three counties in Tennessee to adopt metropolitan governments have taken different approaches to resolving the conflict that a Metro government presents to the requirement to have an elected Sheriff.
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Gun laws in Tennessee regulate the sale, possession, and use of firearms and ammunition. Concealed carry and open-carry of a handgun is permitted with a Tennessee handgun carry permit or an equivalent permit from a reciprocating state. As of July 1, 2014, a permit is no longer required to possess a loaded handgun in a motor vehicle.[138]
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Capital punishment has existed in Tennessee at various times since statehood. Before 1913, the method of execution was hanging. From 1913 to 1915, there was a hiatus on executions but they were reinstated in 1916 when electrocution became the new method. From 1972 to 1978, after the Supreme Court ruled (Furman v. Georgia) capital punishment unconstitutional, there were no further executions. Capital punishment was restarted in 1978, although those prisoners awaiting execution between 1960 and 1978 had their sentences mostly commuted to life in prison.[139] From 1916 to 1960 the state executed 125 inmates.[140] For a variety of reasons there were no further executions until 2000. Since 2000, Tennessee has executed seven prisoners. Tennessee has 59 prisoners on death row (as of October 2018).[141]
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Lethal injection was approved by the legislature in 1998, though those who were sentenced to death before January 1, 1999, may request electrocution.[139] In May 2014, the Tennessee General Assembly passed a law allowing the use of the electric chair for death row executions when lethal injection drugs are not available.[142][143]
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The Mississippi Band of Choctaw Indians is the only federally recognized Native American Indian tribe in the state. It owns 79 acres (32 ha) in Henning, which was placed into federal trust by the tribe in 2012. This is governed directly by the tribe.[144][145]
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Tennessee has a rich variety of public, private, charter, and specialized education facilities ranging from pre-school through university.
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Public higher education is under the oversight of the Tennessee Higher Education Commission which provides guidance to two public university systems—the University of Tennessee system and the Tennessee Board of Regents. In addition a number of private colleges and universities are located throughout the state.
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Public primary and secondary education systems are operated by county, city, or special school districts to provide education at the local level. These school districts operate under the direction of the Tennessee Department of Education. Private schools are found in many counties.
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Over 250 public and private museums are located in Tennessee, including renowned museums such as the Brooks Museum, Pink Palace Museum and Planetarium, Country Music Hall of Fame, and Discovery Park of America.
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State symbols, found in Tennessee Code Annotated; Title 4, Chapter 1, Part 3, include:
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Coordinates: 36°N 86°W / 36°N 86°W / 36; -86
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Table tennis, also known as ping-pong and whiff-whaff, is a sport in which two or four players hit a lightweight ball, also known as the ping-pong ball, back and forth across a table using small rackets. The game takes place on a hard table divided by a net. Except for the initial serve, the rules are generally as follows: players must allow a ball played toward them to bounce one time on their side of the table, and must return it so that it bounces on the opposite side at least once. A point is scored when a player fails to return the ball within the rules. Play is fast and demands quick reactions. Spinning the ball alters its trajectory and limits an opponent's options, giving the hitter a great advantage.
|
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Table tennis is governed by the worldwide organization International Table Tennis Federation (ITTF), founded in 1926. ITTF currently includes 226 member associations.[3] The table tennis official rules are specified in the ITTF handbook.[4] Table tennis has been an Olympic sport since 1988,[5] with several event categories. From 1988 until 2004, these were men's singles, women's singles, men's doubles and women's doubles. Since 2008, a team event has been played instead of the doubles.
|
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The sport originated in Victorian England, where it was played among the upper-class as an after-dinner parlour game.[1][2] It has been suggested that makeshift versions of the game were developed by British military officers in India around the 1860s or 1870s, who brought it back with them.[6] A row of books stood up along the center of the table as a net, two more books served as rackets and were used to continuously hit a golf-ball.[7][8]
|
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+
|
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+
The name "ping-pong" was in wide use before British manufacturer J. Jaques & Son Ltd trademarked it in 1901. The name "ping-pong" then came to describe the game played using the rather expensive Jaques's equipment, with other manufacturers calling it table tennis. A similar situation arose in the United States, where Jaques sold the rights to the "ping-pong" name to Parker Brothers. Parker Brothers then enforced its trademark for the term in the 1920s making the various associations change their names to "table tennis" instead of the more common, but trademarked, term.[9]
|
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|
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The next major innovation was by James W. Gibb, a British enthusiast of table tennis, who discovered novelty celluloid balls on a trip to the US in 1901 and found them to be ideal for the game. This was followed by E.C. Goode who, in 1901, invented the modern version of the racket by fixing a sheet of pimpled, or stippled, rubber to the wooden blade. Table tennis was growing in popularity by 1901 to the extent that tournaments were being organized, books being written on the subject,[7] and an unofficial world championship was held in 1902.
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In 1921, the Table Tennis Association was founded, and in 1926 renamed the English Table Tennis Association.[10] The International Table Tennis Federation (ITTF) followed in 1926.[1][11] London hosted the first official World Championships in 1926. In 1933, the United States Table Tennis Association, now called USA Table Tennis, was formed.[1][12]
|
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In the 1930s, Edgar Snow commented in Red Star Over China that the Communist forces in the Chinese Civil War had a "passion for the English game of table tennis" which he found "bizarre".[13] On the other hand, the popularity of the sport waned in 1930s Soviet Union, partly because of the promotion of team and military sports, and partly because of a theory that the game had adverse health effects.[14]
|
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+
In the 1950s, paddles that used a rubber sheet combined with an underlying sponge layer changed the game dramatically,[1] introducing greater spin and speed.[15] These were introduced to Britain by sports goods manufacturer S.W. Hancock Ltd. The use of speed glue beginning in the mid 1980s increased the spin and speed even further, resulting in changes to the equipment to "slow the game down". Table tennis was introduced as an Olympic sport at the Olympics in 1988.[16]
|
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+
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+
After the 2000 Olympics in Sydney, the ITTF instituted several rule changes that were aimed at making table tennis more viable as a televised spectator sport.[17][18] First, the older 38 mm (1.50 in) balls were officially replaced by 40 mm (1.57 in) balls in October 2000.[7][19] This increased the ball's air resistance and effectively slowed down the game. By that time, players had begun increasing the thickness of the fast sponge layer on their paddles, which made the game excessively fast and difficult to watch on television. A few months later, the ITTF changed from a 21-point to an 11-point scoring system (and the serve rotation was reduced from five points to two), effective in September 2001.[7] This was intended to make games more fast-paced and exciting. The ITTF also changed the rules on service to prevent a player from hiding the ball during service, in order to increase the average length of rallies and to reduce the server's advantage, effective in 2002.[20]
|
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+
For the opponent to have time to realize a serve is taking place, the ball must be tossed a minimum of 16 centimetres (6.3 in) in the air. The ITTF states that all events after July 2014 are played with a new poly material ball.[21]
|
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+
[22]
|
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+
|
23 |
+
The international rules specify that the game is played with a sphere having a mass of 2.7 grams (0.095 oz) and a diameter of 40 millimetres (1.57 in).[23] The rules say that the ball shall bounce up 24–26 cm (9.4–10.2 in) when dropped from a height of 30.5 cm (12.0 in) onto a standard steel block thereby having a coefficient of restitution of 0.89 to 0.92. Balls are now made of a polymer instead of celluloid as of 2015, colored white or orange, with a matte finish. The choice of ball color is made according to the table color and its surroundings. For example, a white ball is easier to see on a green or blue table than it is on a grey table. Manufacturers often indicate the quality of the ball with a star rating system, usually from one to three, three being the highest grade. As this system is not standard across manufacturers, the only way a ball may be used in official competition is upon ITTF approval[23] (the ITTF approval can be seen printed on the ball).
|
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+
|
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+
The 40 mm ball was introduced after the end of the 2000 Summer Olympics; previously a 38 mm ball was standard.[19] This created some controversies. Then World No 1 table tennis professional Vladimir Samsonov threatened to pull out of the World Cup, which was scheduled to debut the new regulation ball on October 12, 2000.[24]
|
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+
|
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+
The table is 2.74 m (9.0 ft) long, 1.525 m (5.0 ft) wide, and 76 cm (2.5 ft) high with any continuous material so long as the table yields a uniform bounce of about 23 cm (9.1 in) when a standard ball is dropped onto it from a height of 30 cm (11.8 in), or about 77%.[25][26] The table or playing surface is uniformly dark coloured and matte, divided into two halves by a net at 15.25 cm (6.0 in) in height. The ITTF approves only wooden tables or their derivates. Concrete tables with a steel net or a solid concrete partition are sometimes available in outside public spaces, such as parks.[27]
|
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+
|
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+
Players are equipped with a laminated wooden racket covered with rubber on one or two sides depending on the grip of the player. The ITTF uses the term "racket",[28] though "bat" is common in Britain, and "paddle" in the U.S. and Canada.
|
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+
|
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+
The wooden portion of the racket, often referred to as the "blade", commonly features anywhere between one and seven plies of wood, though cork, glass fiber, carbon fiber, aluminum fiber, and Kevlar are sometimes used. According to the ITTF regulations, at least 85% of the blade by thickness shall be of natural wood.[29] Common wood types include balsa, limba, and cypress or "hinoki", which is popular in Japan. The average size of the blade is about 17 centimetres (6.7 in) long and 15 centimetres (5.9 in) wide, although the official restrictions only focus on the flatness and rigidity of the blade itself, these dimensions are optimal for most play styles.
|
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+
|
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+
Table tennis regulations allow different rubber surfaces on each side of the racket.[30] Various types of surfaces provide various levels of spin or speed, and in some cases they nullify spin. For example, a player may have a rubber that provides much spin on one side of their racket, and one that provides no spin on the other. By flipping the racket in play, different types of returns are possible. To help a player distinguish between the rubber used by his opposing player, international rules specify that one side must be red while the other side must be black.[29] The player has the right to inspect their opponent's racket before a match to see the type of rubber used and what colour it is. Despite high speed play and rapid exchanges, a player can see clearly what side of the racket was used to hit the ball. Current rules state that, unless damaged in play, the racket cannot be exchanged for another racket at any time during a match.[31]
|
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+
|
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+
According to ITTF rule 2.13.1, the first service is decided by lot,[32] normally a coin toss.[33] It is also common for one player (or the umpire/scorer) to hide the ball in one or the other hand, usually hidden under the table, allowing the other player to guess which hand the ball is in. The correct or incorrect guess gives the "winner" the option to choose to serve, receive, or to choose which side of the table to use. (A common but non-sanctioned method is for the players to play the ball back and forth three times and then play out the point. This is commonly referred to as "serve to play", "rally to serve", "play for serve", or "volley for serve".)
|
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+
|
37 |
+
In game play, the player serving the ball commences a play.[34] The server first stands with the ball held on the open palm of the hand not carrying the paddle, called the freehand, and tosses the ball directly upward without spin, at least 16 cm (6.3 in) high.[35] The server strikes the ball with the racket on the ball's descent so that it touches first his court and then touches directly the receiver's court without touching the net assembly. In casual games, many players do not toss the ball upward; however, this is technically illegal and can give the serving player an unfair advantage.
|
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+
|
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+
The ball must remain behind the endline and above the upper surface of the table, known as the playing surface, at all times during the service. The server cannot use his/her body or clothing to obstruct sight of the ball; the opponent and the umpire must have a clear view of the ball at all times. If the umpire is doubtful of the legality of a service they may first interrupt play and give a warning to the server. If the serve is a clear failure or is doubted again by the umpire after the warning, the receiver scores a point.
|
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+
|
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+
If the service is "good", then the receiver must make a "good" return by hitting the ball back before it bounces a second time on receiver's side of the table so that the ball passes the net and touches the opponent's court, either directly or after touching the net assembly.[36] Thereafter, the server and receiver must alternately make a return until the rally is over. Returning the serve is one of the most difficult parts of the game, as the server's first move is often the least predictable and thus most advantageous shot due to the numerous spin and speed choices at his or her disposal.
|
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+
|
43 |
+
A Let is a rally of which the result is not scored, and is called in the following circumstances:[37]
|
44 |
+
|
45 |
+
A let is also called foul service, if the ball hits the server's side of the table, if the ball does not pass further than the edge and if the ball hits the table edge and hits the net.
|
46 |
+
|
47 |
+
A point is scored by the player for any of several results of the rally:[38]
|
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+
|
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+
A game shall be won by the player first scoring 11 points unless both players score 10 points, when the game shall be won by the first player subsequently gaining a lead of 2 points. A match shall consist of the best of any odd number of games.[41] In competition play, matches are typically best of five or seven games.
|
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+
|
51 |
+
Service alternates between opponents every two points (regardless of winner of the rally) until the end of the game, unless both players score ten points or the expedite system is operated, when the sequences of serving and receiving stay the same but each player serves for only one point in turn (Deuce).[42] The player serving first in a game receives first in the next game of the match.
|
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+
|
53 |
+
After each game, players switch sides of the table. In the last possible game of a match, for example the seventh game in a best of seven matches, players change ends when the first player scores five points, regardless of whose turn it is to serve. If the sequence of serving and receiving is out of turn or the ends are not changed, points scored in the wrong situation are still calculated and the game shall be resumed with the order at the score that has been reached.
|
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+
|
55 |
+
In addition to games between individual players, pairs may also play table tennis. Singles and doubles are both played in international competition, including the Olympic Games since 1988 and the Commonwealth Games since 2002.[43] In 2005, the ITTF announced that doubles table tennis only was featured as a part of team events in the 2008 Olympics.
|
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In doubles, all the rules of single play are applied except for the following.
|
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Service
|
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Order of play, serving and receiving
|
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Men's doubles. Brothers Dmitry Mazunov and Andrey Mazunov in 1989.
|
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Women's doubles finals, 2013 World Table Tennis Championships.
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Mixed doubles finals, 2013 World Table Tennis Championships.
|
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If a game is unfinished after 10 minutes' play and fewer than 18 points have been scored, the expedite system is initiated.[39] The umpire interrupts the game, and the game resumes with players serving for one point in turn. If the expedite system is introduced while the ball is not in play, the previous receiver shall serve first. Under the expedite system, the server must win the point before the opponent makes 13 consecutive returns or the point goes to the opponent. The system can also be initiated at any time at the request of both players or pairs. Once introduced, the expedite system remains in force until the end of the match. A rule to shorten the time of a match, it is mainly seen in defensive players' games.
|
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|
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Though table tennis players grip their rackets in various ways, their grips can be classified into two major families of styles, penhold and shakehand.[45] The rules of table tennis do not prescribe the manner in which one must grip the racket, and numerous grips are employed.
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The penhold grip is so-named because one grips the racket similarly to the way one holds a writing instrument.[46] The style of play among penhold players can vary greatly from player to player. The most popular style, usually referred to as the Chinese penhold style, involves curling the middle, ring, and fourth finger on the back of the blade with the three fingers always touching one another.[46] Chinese penholders favour a round racket head, for a more over-the-table style of play. In contrast, another style, sometimes referred to as the Japanese/Korean penhold grip, involves splaying those three fingers out across the back of the racket, usually with all three fingers touching the back of the racket, rather than stacked upon one another.[46] Sometimes a combination of the two styles occurs, wherein the middle, ring and fourth fingers are straight, but still stacked, or where all fingers may be touching the back of the racket, but are also in contact with one another. Japanese and Korean penholders will often use a square-headed racket for an away-from-the-table style of play. Traditionally these square-headed rackets feature a block of cork on top of the handle, as well as a thin layer of cork on the back of the racket, for increased grip and comfort. Penhold styles are popular among players originating from East Asian countries such as China, Japan, South Korea, and Taiwan.
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Traditionally, penhold players use only one side of the racket to hit the ball during normal play, and the side which is in contact with the last three fingers is generally not used. This configuration is sometimes referred to as "traditional penhold" and is more commonly found in square-headed racket styles. However, the Chinese developed a technique in the 1990s in which a penholder uses both sides of the racket to hit the ball, where the player produces a backhand stroke (most often topspin) known as a reverse penhold backhand by turning the traditional side of the racket to face one's self, and striking the ball with the opposite side of the racket. This stroke has greatly improved and strengthened the penhold style both physically and psychologically, as it eliminates the strategic weakness of the traditional penhold backhand.
|
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+
|
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+
The shakehand grip is so-named because the racket is grasped as if one is performing a handshake.[47] Though it is sometimes referred to as the "tennis" or "Western" grip, it bears no relation to the Western tennis grip, which was popularized on the West Coast of the United States in which the racket is rotated 90°, and played with the wrist turned so that on impact the knuckles face the target. In table tennis, "Western" refers to Western nations, for this is the grip that players native to Europe and the Americas have almost exclusively employed.
|
78 |
+
|
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+
The shakehand grip's simplicity and versatility, coupled with the acceptance among top-level Chinese trainers that the European style of play should be emulated and trained against, has established it as a common grip even in China.[48] Many world class European and East Asian players currently use the shakehand grip, and it is generally accepted that shakehands is easier to learn than penholder, allowing a broader range of playing styles both offensive and defensive.[49]
|
80 |
+
|
81 |
+
The Seemiller grip is named after the American table tennis champion Danny Seemiller, who used it. It is achieved by placing the thumb and index finger on either side of the bottom of the racquet head and holding the handle with the rest of the fingers. Since only one side of the racquet is used to hit the ball, two contrasting rubber types can be applied to the blade, offering the advantage of "twiddling" the racket to fool the opponent. Seemiller paired inverted rubber with anti-spin rubber. Many players today combine inverted and long-pipped rubber. The grip is considered exceptional for blocking, especially on the backhand side, and for forehand loops of backspin balls.[50]
|
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+
The Seemiller grip's popularity reached its apex in 1985 when four (Danny Seemiller, Ricky Seemiller, Eric Boggan and Brian Masters) of the United States' five participants in the World Championships used it.[50]
|
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|
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Shakehand grip (Vladimir Samsonov)
|
85 |
+
|
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+
Chinese penhold (Ma Lin)
|
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|
88 |
+
Traditional penhold (Ryu Seung-min)
|
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|
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+
'A good ready position will enable you to move quickly into position and to stay balanced whilst playing powerful strokes.'[51]
|
91 |
+
|
92 |
+
The stance in table tennis is also known as the 'ready position'. It is the position every player initially adopts when receiving and returns to after playing a shot in order to be prepared to make the next shot. It involves the feet being spaced wider than shoulder width and a partial crouch being adopted; the crouch is an efficient posture for moving quickly from and also preloads the muscles enabling a more dynamic movement. The upper torso is positioned slightly forward and the player is looking forwards. The racket is held at the ready with a bent arm. The position should feel balanced and provide a solid base for striking and quick lateral movement. Players may tailor their stance based upon their personal preferences, and alter it during the game based upon the specific circumstances.[52]
|
93 |
+
|
94 |
+
Table tennis strokes generally break down into offensive and defensive categories.
|
95 |
+
|
96 |
+
Also known as speed drive, a direct hit on the ball propelling it forward back to the opponent. This stroke differs from speed drives in other racket sports like tennis because the racket is primarily perpendicular to the direction of the stroke and most of the energy applied to the ball results in speed rather than spin, creating a shot that does not arc much, but is fast enough that it can be difficult to return. A speed drive is used mostly for keeping the ball in play, applying pressure on the opponent, and potentially opening up an opportunity for a more powerful attack.
|
97 |
+
|
98 |
+
Perfected during the 1960s,[1][53] the loop is essentially the reverse of the chop. The racket is parallel to the direction of the stroke ("closed") and the racket thus grazes the ball, resulting in a large amount of topspin. A good loop drive will arc quite a bit, and once striking the opponent's side of the table will jump forward, much like a kick serve in tennis. Most professional players nowadays, such as Ding Ning, Timo Boll and Zhang Jike, primarily use loop for offense.
|
99 |
+
|
100 |
+
The counter-hit is usually a counterattack against drives, normally high loop drives. The racket is held closed and near to the ball, which is hit with a short movement "off the bounce" (immediately after hitting the table) so that the ball travels faster to the other side. Kenta Matsudaira is known for primarily using counter-hit for offense.
|
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|
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+
When a player tries to attack a ball that has not bounced beyond the edge of the table, the player does not have the room to wind up in a backswing. The ball may still be attacked, however, and the resulting shot is called a flip because the backswing is compressed into a quick wrist action. A flip is not a single stroke and can resemble either a loop drive or a loop in its characteristics. What identifies the stroke is that the backswing is compressed into a short wrist flick.
|
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+
|
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A player will typically execute a smash when the opponent has returned a ball that bounces too high or too close to the net. It is nearly always done with a forehand stroke. Smashing use rapid acceleration to impart as much speed on the ball as possible so that the opponent cannot react in time. The racket is generally perpendicular to the direction of the stroke. Because the speed is the main aim of this shot, the spin on the ball is often minimal, although it can be applied as well. An offensive table tennis player will think of a rally as a build-up to a winning smash. Smash is used more often with penhold grip.
|
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|
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+
The push (or "slice" in Asia) is usually used for keeping the point alive and creating offensive opportunities. A push resembles a tennis slice: the racket cuts underneath the ball, imparting backspin and causing the ball to float slowly to the other side of the table. A push can be difficult to attack because the backspin on the ball causes it to drop toward the table upon striking the opponent's racket. In order to attack a push, a player must usually loop (if the push is long) or flip (if the push is short) the ball back over the net. Often, the best option for beginners is to simply push the ball back again, resulting in pushing rallies. Against good players, it may be the worst option because the opponent will counter with a loop, putting the first player in a defensive position. Pushing can have advantages in some circumstances, such as when the opponent makes easy mistakes.
|
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+
|
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+
A chop is the defensive, backspin counterpart to the offensive loop drive.[54] A chop is essentially a bigger, heavier push, taken well back from the table. The racket face points primarily horizontally, perhaps a little bit upward, and the direction of the stroke is straight down. The object of a defensive chop is to match the topspin of the opponent's shot with backspin. A good chop will float nearly horizontally back to the table, in some cases having so much backspin that the ball actually rises. Such a chop can be extremely difficult to return due to its enormous amount of backspin. Some defensive players can also impart no-spin or sidespin variations of the chop. Some famous choppers include Joo Sae-hyuk and Wu Yang.
|
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+
|
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+
A block is executed by simply placing the racket in front of the ball right after the ball bounces; thus, the ball rebounds back toward the opponent with nearly as much energy as it came in with. This requires precision, since the ball's spin, speed, and location all influence the correct angle of a block. It is very possible for an opponent to execute a perfect loop, drive, or smash, only to have the blocked shot come back just as fast. Due to the power involved in offensive strokes, often an opponent simply cannot recover quickly enough to return the blocked shot, especially if the block is aimed at an unexpected side of the table. Blocks almost always produce the same spin as was received, many times topspin.
|
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+
|
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+
The defensive lob propels the ball about five metres in height, only to land on the opponent's side of the table with great amounts of spin.[55] The stroke itself consists of lifting the ball to an enormous height before it falls back to the opponent's side of the table. A lob can have nearly any kind of spin. Though the opponent may smash the ball hard and fast, a good defensive lob could be more difficult to return due to the unpredictability and heavy amounts of the spin on the ball.[55] Thus, though backed off the table by tens of feet and running to reach the ball, a good defensive player can still win the point using good lobs. Lob is used less frequently by professional players. A notable exception is Michael Maze.
|
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+
|
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+
Adding spin onto the ball causes major changes in table tennis gameplay. Although nearly every stroke or serve creates some kind of spin, understanding the individual types of spin allows players to defend against and use different spins effectively.[56]
|
115 |
+
|
116 |
+
Backspin is where the bottom half of the ball is rotating away from the player, and is imparted by striking the base of the ball with a downward movement.[56] At the professional level, backspin is usually used defensively in order to keep the ball low.[57] Backspin is commonly employed in service because it is harder to produce an offensive return, though at the professional level most people serve sidespin with either backspin or topspin. Due to the initial lift of the ball, there is a limit on how much speed with which one can hit the ball without missing the opponent's side of the table. However, backspin also makes it harder for the opponent to return the ball with great speed because of the required angular precision of the return. Alterations are frequently made to regulations regarding equipment in an effort to maintain a balance between defensive and offensive spin choices.[citation needed] It is actually possible to smash with backspin offensively, but only on high balls that are close to the net.
|
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+
|
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The topspin stroke has a smaller influence on the first part of the ball-curve. Like the backspin stroke, however, the axis of spin remains roughly perpendicular to the trajectory of the ball thus allowing for the Magnus effect to dictate the subsequent curvature. After the apex of the curve, the ball dips downwards as it approaches the opposing side, before bouncing. On the bounce, the topspin will accelerate the ball, much in the same way that a wheel which is already spinning would accelerate upon making contact with the ground. When the opponent attempts to return the ball, the topspin causes the ball to jump upwards and the opponent is forced to compensate for the topspin by adjusting the angle of his or her racket. This is known as "closing the racket".
|
119 |
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|
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The speed limitation of the topspin stroke is minor compared to the backspin stroke. This stroke is the predominant technique used in professional competition because it gives the opponent less time to respond. In table tennis topspin is regarded as an offensive technique due to increased ball speed, lower bio-mechanical efficiency and the pressure that it puts on the opponent by reducing reaction time. (It is possible to play defensive topspin-lobs from far behind the table, but only highly skilled players use this stroke with any tactical efficiency.) Topspin is the least common type of spin to be found in service at the professional level, simply because it is much easier to attack a top-spin ball that is not moving at high speed.
|
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|
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This type of spin is predominantly employed during service, wherein the contact angle of the racket can be more easily varied. Unlike the two aforementioned techniques, sidespin causes the ball to spin on an axis which is vertical, rather than horizontal. The axis of rotation is still roughly perpendicular to the trajectory of the ball. In this circumstance, the Magnus effect will still dictate the curvature of the ball to some degree. Another difference is that unlike backspin and topspin, sidespin will have relatively very little effect on the bounce of the ball, much in the same way that a spinning top would not travel left or right if its axis of rotation were exactly vertical. This makes sidespin a useful weapon in service, because it is less easily recognized when bouncing, and the ball "loses" less spin on the bounce. Sidespin can also be employed in offensive rally strokes, often from a greater distance, as an adjunct to topspin or backspin. This stroke is sometimes referred to as a "hook". The hook can even be used in some extreme cases to circumvent the net when away from the table.
|
123 |
+
|
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+
Players employ this type of spin almost exclusively when serving, but at the professional level, it is also used from time to time in the lob. Unlike any of the techniques mentioned above, corkspin (or "drill-spin") has the axis of spin relatively parallel to the ball's trajectory, so that the Magnus effect has little or no effect on the trajectory of a cork-spun ball: upon bouncing, the ball will dart right or left (according to the direction of the spin), severely complicating the return. In theory this type of spin produces the most obnoxious effects, but it is less strategically practical than sidespin or backspin, because of the limitations that it imposes upon the opponent during their return. Aside from the initial direction change when bouncing, unless it goes out of reach, the opponent can counter with either topspin or backspin. A backspin stroke is similar in the fact that the corkspin stroke has a lower maximum velocity, simply due to the contact angle of the racket when producing the stroke. To impart a spin on the ball which is parallel to its trajectory, the racket must be swung more or less perpendicular to the trajectory of the ball, greatly limiting the forward momentum that the racket transfers to the ball. Corkspin is almost always mixed with another variety of spin, since alone, it is not only less effective but also harder to produce.
|
125 |
+
|
126 |
+
Competitive table tennis is popular in East Asia and Europe, and has been[vague] gaining attention in the United States.[62] The most important international competitions are the World Table Tennis Championships, the Table Tennis World Cup, the Olympics and the ITTF World Tour. Continental competitions include the following:
|
127 |
+
|
128 |
+
Chinese players have won 60% of the men's World Championships since 1959;[63] in the women's competition for the Corbillin Cup, Chinese players have won all but three of the World Championships since 1971.[64] Other strong teams come from East Asia and Europe, including countries such as Austria, Belarus, Germany, Hong Kong, Portugal, Japan, South Korea, Singapore, Sweden, and Taiwan.[65]
|
129 |
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|
130 |
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There are professional competitions at the clubs level; the respective leagues of Austria, Belgium, China (China Table Tennis Super League), Japan (T.League), France, Germany (Bundesliga), and Russia are examples of the highest level. There are also some important international club teams competitions such as the European Champions League and its former competitor,[vague] the European Club Cup, where the top club teams from European countries compete.
|
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+
|
132 |
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According to the New York Times, 31% of the table tennis players at the 2016 Summer Olympics were naturalized. The rate was twice as high as the next sport, basketball, which featured 15% of naturalized players.[66]
|
133 |
+
|
134 |
+
In particular, Chinese-born players representing Singapore have won three Olympic medals, more than native Singaporeans have ever won in all sports. However, these successes have been very controversial in Singapore.[67] In 2014, Singapore Table Tennis Association's president Lee Bee Wah quit over this issue;[68] however, her successor Ellen Lee has basically continued on this path.[69]
|
135 |
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|
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+
The rate of naturalization accelerated after the ITTF's 2009 decision (one year after China won every possible Olympic medal in the sport) to reduce the number of entries per association in both the Olympics and the World Table Tennis Championships.[70]
|
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+
|
138 |
+
In 2019, the ITTF adopted new regulations which state that players who acquired a new nationality may not represent their new association before:[71]
|
139 |
+
|
140 |
+
An official hall of fame exists at the ITTF Museum.[72] A Grand Slam is earned by a player who wins singles crowns at the Olympic Games, World Championships, and World Cup.[73] Jan-Ove Waldner of Sweden first completed the grand slam at 1992 Olympic Games. Deng Yaping of China is the first female recorded at the inaugural Women's World Cup in 1996. The following table presents an exhaustive list of all players to have completed a grand slam.
|
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|
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+
Jean-Philippe Gatien (France) and Wang Hao (China) won both the World Championships and the World Cup, but lost in the gold medal matches at the Olympics. Jörgen Persson (Sweden) also won the titles except the Olympic Games. Persson is one of the three table tennis players to have competed at seven Olympic Games. Ma Lin (China) won both the Olympic gold and the World Cup, but lost (three times, in 1999, 2005, and 2007) in the finals of the World Championships.
|
143 |
+
|
144 |
+
Founded in 1926, the International Table Tennis Federation (ITTF) is the worldwide governing body for table tennis, which maintains an international ranking system in addition to organizing events like the World Table Tennis Championships.[12] In 2007, the governance for table tennis for persons with a disability was transferred from the International Paralympic Committee to the ITTF.[83]
|
145 |
+
|
146 |
+
On many continents, there is a governing body responsible for table tennis on that continent. For example, the European Table Tennis Union (ETTU) is the governing body responsible for table tennis in Europe.[84] There are also national bodies and other local authorities responsible for the sport, such as USA Table Tennis (USATT), which is the national governing body for table tennis in the United States.[12]
|
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+
|
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+
Hardbat table tennis uses rackets with short outward "pips" and no sponge, resulting in decreased speeds and reduced spin. World Championship of Ping Pong uses old-fashioned wooden paddles covered with sandpaper.
|
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Tennis is a racket sport that can be played individually against a single opponent (singles) or between two teams of two players each (doubles). Each player uses a tennis racket that is strung with cord to strike a hollow rubber ball covered with felt over or around a net and into the opponent's court. The object of the game is to maneuver the ball in such a way that the opponent is not able to play a valid return. The player who is unable to return the ball will not gain a point, while the opposite player will.
|
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|
7 |
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Tennis is an Olympic sport and is played at all levels of society and at all ages. The sport can be played by anyone who can hold a racket, including wheelchair users. The modern game of tennis originated in Birmingham, England, in the late 19th century as lawn tennis.[1] It had close connections both to various field (lawn) games such as croquet and bowls as well as to the older racket sport today called real tennis. During most of the 19th century, in fact, the term tennis referred to real tennis, not lawn tennis.
|
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|
9 |
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The rules of modern tennis have changed little since the 1890s. Two exceptions are that from 1908 to 1961 the server had to keep one foot on the ground at all times, and the adoption of the tiebreak in the 1970s. A recent addition to professional tennis has been the adoption of electronic review technology coupled with a point-challenge system, which allows a player to contest the line call of a point, a system known as Hawk-Eye.
|
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|
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Tennis is played by millions of recreational players and is also a popular worldwide spectator sport. The four Grand Slam tournaments (also referred to as the Majors) are especially popular: the Australian Open played on hard courts, the French Open played on red clay courts, Wimbledon played on grass courts, and the US Open also played on hard courts.
|
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|
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Historians believe that the game's ancient origin lay in 12th century northern France, where a ball was struck with the palm of the hand.[2] Louis X of France was a keen player of jeu de paume ("game of the palm"), which evolved into real tennis, and became notable as the first person to construct indoor tennis courts in the modern style. Louis was unhappy with playing tennis outdoors and accordingly had indoor, enclosed courts made in Paris "around the end of the 13th century".[3] In due course this design spread across royal palaces all over Europe.[3] In June 1316 at Vincennes, Val-de-Marne and following a particularly exhausting game, Louis drank a large quantity of cooled wine and subsequently died of either pneumonia or pleurisy, although there was also suspicion of poisoning.[4] Because of the contemporary accounts of his death, Louis X is history's first tennis player known by name.[4] Another of the early enthusiasts of the game was King Charles V of France, who had a court set up at the Louvre Palace.[5]
|
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|
15 |
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It was not until the 16th century that rackets came into use and the game began to be called "tennis", from the French term tenez, which can be translated as "hold!", "receive!" or "take!", an interjection used as a call from the server to his opponent.[6] It was popular in England and France, although the game was only played indoors where the ball could be hit off the wall. Henry VIII of England was a big fan of this game, which is now known as real tennis.[7] During the 18th and early 19th centuries, as real tennis declined, new racket sports emerged in England.[8]
|
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|
17 |
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The invention of the first lawn mower in 1830, in Britain, is believed to have been a catalyst, for the preparation of modern-style grass courts, sporting ovals, playing fields, pitches, greens, etc. This in turn led to the codification of modern rules for many sports, including lawn tennis, most football codes, lawn bowls and others.[9]
|
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|
19 |
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Between 1859 and 1865 Harry Gem, a solicitor and his friend Augurio Perera developed a game that combined elements of racquets and the Basque ball game pelota, which they played on Perera's croquet lawn in Birmingham in England.[10][11] In 1872, along with two local doctors, they founded the world's first tennis club on Avenue Road, Leamington Spa.[12] This is where "lawn tennis" was used as a name of activity by a club for the first time. After Leamington, the second club to take up the game of lawn tennis appears to have been the Edgbaston Archery and Croquet Society, also in Birmingham.
|
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|
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In Tennis: A Cultural History, Heiner Gillmeister reveals that on December 8, 1874, British army officer Walter Clopton Wingfield wrote to Harry Gem, commenting that he (Wingfield) had been experimenting with his version of lawn tennis “for a year and a half”.[13] In December 1873, Wingfield designed and patented a game which he called sphairistikè (Greek: σφαιριστική, meaning "ball-playing"), and was soon known simply as "sticky" – for the amusement of guests at a garden party on his friend's estate of Nantclwyd Hall, in Llanelidan, Wales.[14] According to R. D. C. Evans, turfgrass agronomist, "Sports historians all agree that [Wingfield] deserves much of the credit for the development of modern tennis."[8][15] According to Honor Godfrey, museum curator at Wimbledon, Wingfield "popularized this game enormously. He produced a boxed set which included a net, poles, rackets, balls for playing the game – and most importantly you had his rules. He was absolutely terrific at marketing and he sent his game all over the world. He had very good connections with the clergy, the law profession, and the aristocracy and he sent thousands of sets out in the first year or so, in 1874."[16] The world's oldest annual tennis tournament took place at Leamington Lawn Tennis Club in Birmingham in 1874.[17] This was three years before the All England Lawn Tennis and Croquet Club would hold its first championships at Wimbledon, in 1877. The first Championships culminated in a significant debate on how to standardise the rules.[16]
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In the U.S. in 1874 Mary Ewing Outerbridge, a young socialite, returned from Bermuda with a sphairistikè set. She became fascinated by the game of tennis after watching British army officers play.[18] She laid out a tennis court at the Staten Island Cricket Club at Camp Washington, Tompkinsville, Staten Island, New York. The first American National championship was played there in September 1880. An Englishman named O.E. Woodhouse won the singles title, and a silver cup worth $100, by defeating Canadian I. F. Hellmuth.[19] There was also a doubles match which was won by a local pair. There were different rules at each club. The ball in Boston was larger than the one normally used in New York.
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On 21 May 1881, the oldest nationwide tennis organization in the world[20] was formed, the United States National Lawn Tennis Association (now the United States Tennis Association) in order to standardize the rules and organize competitions.[21] The U.S. National Men's Singles Championship, now the US Open, was first held in 1881 at the Newport Casino, Newport, Rhode Island.[22] The U.S. National Women's Singles Championships were first held in 1887 in Philadelphia.[23]
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Tennis also became popular in France, where the French Championships dates to 1891 although until 1925 it was open only to tennis players who were members of French clubs.[24] Thus, Wimbledon, the US Open, the French Open, and the Australian Open (dating to 1905) became and have remained the most prestigious events in tennis.[25][26] Together these four events are called the Majors or Slams (a term borrowed from bridge rather than baseball).[27]
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In 1913, the International Lawn Tennis Federation (ILTF), now the International Tennis Federation (ITF), was founded and established three official tournaments as the major championships of the day. The World Grass Court Championships were awarded to Great Britain. The World Hard Court Championships were awarded to France; the term "hard court" was used for clay courts at the time. Some tournaments were held in Belgium instead. And the World Covered Court Championships for indoor courts was awarded annually; Sweden, France, Great Britain, Denmark, Switzerland and Spain each hosted the tournament.[28] At a meeting held on 16 March 1923 in Paris, the title 'World Championship' was dropped and a new category of Official Championship was created for events in Great Britain, France, the United States, and Australia – today's Grand Slam events.[28][29] The impact on the four recipient nations to replace the ‘world championships’ with ‘official
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championships’ was simple in a general sense: each became a major nation of the federation with enhanced voting power and each now operated a major event.[28]
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The comprehensive rules promulgated in 1924 by the ILTF, have remained largely stable in the ensuing eighty years, the one major change being the addition of the tiebreak system designed by Jimmy Van Alen.[30] That same year, tennis withdrew from the Olympics after the 1924 Games but returned 60 years later as a 21-and-under demonstration event in 1984. This reinstatement was credited by the efforts by the then ITF President Philippe Chatrier, ITF General Secretary David Gray and ITF Vice President Pablo Llorens, and support from IOC President Juan Antonio Samaranch. The success of the event was overwhelming and the IOC decided to reintroduce tennis as a full medal sport at Seoul in 1988.[31][32]
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The Davis Cup, an annual competition between men's national teams, dates to 1900.[33] The analogous competition for women's national teams, the Fed Cup, was founded as the Federation Cup in 1963 to celebrate the 50th anniversary of the founding of the ITF.[34]
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In 1926, promoter C. C. Pyle established the first professional tennis tour with a group of American and French tennis players playing exhibition matches to paying audiences.[26][35] The most notable of these early professionals were the American Vinnie Richards and the Frenchwoman Suzanne Lenglen.[26][36] Once a player turned pro he or she was no longer permitted to compete in the major (amateur) tournaments.[26]
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In 1968, commercial pressures and rumors of some amateurs taking money under the table led to the abandonment of this distinction, inaugurating the Open Era, in which all players could compete in all tournaments, and top players were able to make their living from tennis. With the beginning of the Open Era, the establishment of an international professional tennis circuit, and revenues from the sale of television rights, tennis's popularity has spread worldwide, and the sport has shed its middle-class English-speaking image[37] (although it is acknowledged that this stereotype still exists).[37][38]
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In 1954, Van Alen founded the International Tennis Hall of Fame, a non-profit museum in Newport, Rhode Island.[39] The building contains a large collection of tennis memorabilia as well as a hall of fame honouring prominent members and tennis players from all over the world. Each year, a grass court tournament and an induction ceremony honoring new Hall of Fame members are hosted on its grounds.
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Part of the appeal of tennis stems from the simplicity of equipment required for play. Beginners need only a racket and balls.
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The components of a tennis racket include a handle, known as the grip, connected to a neck which joins a roughly elliptical frame that holds a matrix of tightly pulled strings. For the first 100 years of the modern game, rackets were made of wood and of standard size, and strings were of animal gut. Laminated wood construction yielded more strength in rackets used through most of the 20th century until first metal and then composites of carbon graphite, ceramics, and lighter metals such as titanium were introduced. These stronger materials enabled the production of oversized rackets that yielded yet more power. Meanwhile, technology led to the use of synthetic strings that match the feel of gut yet with added durability.
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Under modern rules of tennis, the rackets must adhere to the following guidelines;[40]
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The rules regarding rackets have changed over time, as material and engineering advances have been made. For example, the maximum length of the frame had been 32 inches (81 cm) until 1997, when it was shortened to 29 inches (74 cm).[41]
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Many companies manufacture and distribute tennis rackets. Wilson, Head and Babolat are some of the more commonly used brands; however, many more companies exist.[example needed] The same companies sponsor players to use these rackets in the hopes that the company name will become more well known by the public.
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Tennis balls were originally made of cloth strips stitched together with thread and stuffed with feathers.[42] Modern tennis balls are made of hollow vulcanized rubber with a felt coating. Traditionally white, the predominant colour was gradually changed to optic yellow in the latter part of the 20th century to allow for improved visibility. Tennis balls must conform to certain criteria for size, weight, deformation, and bounce to be approved for regulation play. The International Tennis Federation (ITF) defines the official diameter as 65.41–68.58 mm (2.575–2.700 in). Balls must weigh between 56.0 and 59.4 g (1.98 and 2.10 oz).[43] Tennis balls were traditionally manufactured in the United States and Europe. Although the process of producing the balls has remained virtually unchanged for the past 100 years, the majority of manufacturing now takes place in the Far East. The relocation is due to cheaper labour costs and materials in the region.[44] Tournaments that are played under the ITF Rules of Tennis must use balls that are approved by the International Tennis Federation (ITF) and be named on the official ITF list of approved tennis balls.[45]
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Advanced players improve their performance through a number of accoutrements. Vibration dampeners may be interlaced in the proximal part of the string array for improved feel. Racket handles may be customized with absorbent or rubber-like materials to improve the players' grip. Players often use sweat bands on their wrists to keep their hands dry and head bands or bandanas to keep the sweat out of their eyes as well. Finally, although the game can be played in a variety of shoes, specialized tennis shoes have wide, flat soles for stability and a built-up front structure to avoid excess wear.
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Tennis is played on a rectangular, flat surface. The court is 78 feet (23.77 m) long, and 27 feet (8.2 m) wide for singles matches and 36 ft (11 m) for doubles matches.[46] Additional clear space around the court is required in order for players to reach overrun balls. A net is stretched across the full width of the court, parallel with the baselines, dividing it into two equal ends. It is held up by either a cord or metal cable of diameter no greater than 0.8 cm (1⁄3 in).[47] The net is 3 feet 6 inches (1.07 m) high at the posts and 3 feet (0.91 m) high in the center.[46] The net posts are 3 feet (0.91 m) outside the doubles court on each side or, for a singles net, 3 feet (0.91 m) outside the singles court on each side.
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The modern tennis court owes its design to Major Walter Clopton Wingfield. In 1873, Wingfield patented a court much the same as the current one for his stické tennis (sphairistike). This template was modified in 1875 to the court design that exists today, with markings similar to Wingfield's version, but with the hourglass shape of his court changed to a rectangle.[48]
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Tennis is unusual in that it is played on a variety of surfaces.[49] Grass, clay, and hardcourts of concrete or asphalt topped with acrylic are the most common. Occasionally carpet is used for indoor play, with hardwood flooring having been historically used. Artificial turf courts can also be found.
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The lines that delineate the width of the court are called the baseline (farthest back) and the service line (middle of the court). The short mark in the center of each baseline is referred to as either the hash mark or the center mark. The outermost lines that make up the length are called the doubles sidelines; they are the boundaries for doubles matches. The lines to the inside of the doubles sidelines are the singles sidelines, and are the boundaries in singles play. The area between a doubles sideline and the nearest singles sideline is called the doubles alley, playable in doubles play. The line that runs across the center of a player's side of the court is called the service line because the serve must be delivered into the area between the service line and the net on the receiving side. Despite its name, this is not where a player legally stands when making a serve.[50]
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The line dividing the service line in two is called the center line or center service line. The boxes this center line creates are called the service boxes; depending on a player's position, they have to hit the ball into one of these when serving.[51] A ball is out only if none of it has hit the area inside the lines, or the line, upon its first bounce. All lines are required to be between 1 and 2 inches (25 and 51 mm) in width, with the exception of the baseline which can be up to 4 inches (100 mm) wide, although in practice it is often the same width as the others.[50]
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The players or teams start on opposite sides of the net. One player is designated the server, and the opposing player is the receiver. The choice to be server or receiver in the first game and the choice of ends is decided by a coin toss before the warm-up starts. Service alternates game by game between the two players or teams. For each point, the server starts behind the baseline, between the center mark and the sideline. The receiver may start anywhere on their side of the net. When the receiver is ready, the server will serve, although the receiver must play to the pace of the server.
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For a service to be legal, the ball must travel over the net without touching it into the diagonally opposite service box. If the ball hits the net but lands in the service box, this is a let or net service, which is void, and the server retakes that serve. The player can serve any number of let services in a point and they are always treated as voids and not as faults. A fault is a serve that falls long or wide of the service box, or does not clear the net. There is also a "foot fault" when a player's foot touches the baseline or an extension of the center mark before the ball is hit. If the second service, after a fault, is also a fault, the server double faults, and the receiver wins the point. However, if the serve is in, it is considered a legal service.
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A legal service starts a rally, in which the players alternate hitting the ball across the net. A legal return consists of a player hitting the ball so that it falls in the server's court, before it has bounced twice or hit any fixtures except the net. A player or team cannot hit the ball twice in a row. The ball must travel over the net into the other players' court. A ball that hits the net during a rally is considered a legal return as long as it crosses into the opposite side of the court. The first player or team to fail to make a legal return loses the point. The server then moves to the other side of the service line at the start of a new point.[52]
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A game consists of a sequence of points played with the same player serving. A game is won by the first player to have won at least four points in total and at least two points more than the opponent. The running score of each game is described in a manner peculiar to tennis: scores from zero to three points are described as "love", "15", "30", and "40", respectively. If at least three points have been scored by each player, making the player's scores equal at 40 apiece, the score is not called out as "40–40", but rather as "deuce". If at least three points have been scored by each side and a player has one more point than his opponent, the score of the game is "advantage" for the player in the lead. During informal games, "advantage" can also be called "ad in" or "van in" when the serving player is ahead, and "ad out" or "van out" when the receiving player is ahead; alternatively, either player may simply call out "my ad" or "your ad" during informal play.
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The score of a tennis game during play is always read with the serving player's score first. In tournament play, the chair umpire calls the point count (e.g., "15-love") after each point. At the end of a game, the chair umpire also announces the winner of the game and the overall score.
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A set consists of a sequence of games played with service alternating between games, ending when the count of games won meets certain criteria. Typically, a player wins a set by winning at least six games and at least two games more than the opponent. If one player has won six games and the opponent five, an additional game is played. If the leading player wins that game, the player wins the set 7–5. If the trailing player wins the game (tying the set 6–6) a tie-break is played. A tie-break, played under a separate set of rules, allows one player to win one more game and thus the set, to give a final set score of 7–6. A "love" set means that the loser of the set won zero games, colloquially termed a 'jam donut' in the US.[53] In tournament play, the chair umpire announces the winner of the set and the overall score. The final score in sets is always read with the winning player's score first, e.g. "6–2, 4–6, 6–0, 7–5".
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A match consists of a sequence of sets. The outcome is determined through a best of three or five sets system. On the professional circuit, men play best-of-five-set matches at all four Grand Slam tournaments, Davis Cup, and the final of the Olympic Games and best-of-three-set matches at all other tournaments, while women play best-of-three-set matches at all tournaments. The first player to win two sets in a best-of-three, or three sets in a best-of-five, wins the match.[54] Only in the final sets of matches at the French Open, the Olympic Games, and Fed Cup are tie-breaks not played. In these cases, sets are played indefinitely until one player has a two-game lead, occasionally leading to some remarkably long matches.
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In tournament play, the chair umpire announces the end of the match with the well-known phrase "Game, set, match" followed by the winning person's or team's name.
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A game point occurs in tennis whenever the player who is in the lead in the game needs only one more point to win the game. The terminology is extended to sets (set point), matches (match point), and even championships (championship point). For example, if the player who is serving has a score of 40-love, the player has a triple game point (triple set point, etc.) as the player has three consecutive chances to win the game. Game points, set points, and match points are not part of official scoring and are not announced by the chair umpire in tournament play.
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A break point occurs if the receiver, not the server, has a chance to win the game with the next point. Break points are of particular importance because serving is generally considered advantageous, with servers being expected to win games in which they are serving. A receiver who has one (score of 30–40 or advantage), two (score of 15–40) or three (score of love-40) consecutive chances to win the game has break point, double break point or triple break point, respectively. If the receiver does, in fact, win their break point, the game is awarded to the receiver, and the receiver is said to have converted their break point. If the receiver fails to win their break point it is called a failure to convert. Winning break points, and thus the game, is also referred to as breaking serve, as the receiver has disrupted, or broken the natural advantage of the server. If in the following game the previous server also wins a break point it is referred to as breaking back. Except where tie-breaks apply, at least one break of serve is required to win a set (otherwise a two-game lead would never occur).
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Another, however informal, tennis format is called Canadian doubles. This involves three players, with one person playing against a doubles team. The single player gets to utilize the alleys normally reserved only for a doubles team. Conversely, the doubles team does not use the alleys when executing a shot. The scoring is the same as for a regular game. This format is not sanctioned by any official body.
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"Australian doubles", another informal and unsanctioned form of tennis, is played with similar rules to the Canadian doubles style, only in this version, players rotate court position after each game, each player taking a turn at playing alone against the other two. As such, each player plays doubles and singles over the course of a match, with the singles player always serving. Scoring styles vary, but one popular method is to assign a value of 2 points to each game, with the server taking both points if he or she holds serve and the doubles team each taking one if they break serve.
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Wheelchair tennis can be played by able-bodied players as well as people who require a wheelchair for mobility. An extra bounce is permitted. This rule makes it possible to have mixed wheelchair and able-bodied matches. It is possible for a doubles team to consist of a wheelchair player and an able-bodied player (referred to as "one-up, one-down"), or for a wheelchair player to play against an able-bodied player. In such cases, the extra bounce is permitted for the wheelchair users only.
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In most professional play and some amateur competition, there is an officiating head judge or chair umpire (usually referred to simply as the umpire), who sits in a raised chair to one side of the court. The umpire has absolute authority to make factual determinations. The umpire may be assisted by line judges, who determine whether the ball has landed within the required part of the court and who also call foot faults. There also may be a net judge who determines whether the ball has touched the net during service. The umpire has the right to overrule a line judge or a net judge if the umpire is sure that a clear mistake has been made.[57]
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In past tournaments, line judges tasked with calling the serve were sometimes assisted by electronic sensors that beeped to indicate an out-of-bounds serve; one such system was called "Cyclops".[58] Cyclops has since largely been replaced by the Hawk-Eye system.[59][60] In professional tournaments using this system, players are allowed three unsuccessful appeals per set, plus one additional appeal in the tie-break to challenge close line calls by means of an electronic review. The US Open, Miami Masters, US Open Series, and World Team Tennis started using this challenge system in 2006 and the Australian Open and Wimbledon introduced the system in 2007.[61] In clay-court matches, such as at the French Open, a call may be questioned by reference to the mark left by the ball's impact on the court surface.
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The referee, who is usually located off the court, is the final authority about tennis rules. When called to the court by a player or team captain, the referee may overrule the umpire's decision if the tennis rules were violated (question of law) but may not change the umpire's decision on a question of fact. If, however, the referee is on the court during play, the referee may overrule the umpire's decision. (This would only happen in Davis Cup or Fed Cup matches, not at the World Group level, when a chair umpire from a non-neutral country is in the chair).[57]
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Ball boys and girls may be employed to retrieve balls, pass them to the players, and hand players their towels. They have no adjudicative role. In rare events (e.g., if they are hurt or if they have caused a hindrance), the umpire may ask them for a statement of what actually happened. The umpire may consider their statements when making a decision. In some leagues, especially junior leagues, players make their own calls, trusting each other to be honest. This is the case for many school and university level matches. The referee or referee's assistant, however, can be called on court at a player's request, and the referee or assistant may change a player's call. In unofficiated matches, a ball is out only if the player entitled to make the call is sure that the ball is out.
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In tennis, a junior is a player under 18 who is still legally protected by a parent or guardian. Players on the main adult tour who are under 18 must have documents signed by a parent or guardian. These players, however, are still eligible to play in junior tournaments.
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The International Tennis Federation (ITF) conducts a junior tour that allows juniors to establish a world ranking and an Association of Tennis Professionals (ATP) or Women's Tennis Association (WTA) ranking. Most juniors who enter the international circuit do so by progressing through ITF, Satellite, Future, and Challenger tournaments before entering the main circuit. The latter three circuits also have adults competing in them. Some juniors, however, such as Australian Lleyton Hewitt and Frenchman Gaël Monfils, have catapulted directly from the junior tour to the ATP tour by dominating the junior scene or by taking advantage of opportunities given to them to participate in professional tournaments.
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In 2004, the ITF implemented a new rankings scheme to encourage greater participation in doubles, by combining two rankings (singles and doubles) into one combined tally.[62] Junior tournaments do not offer prize money except for the Grand Slam tournaments, which are the most prestigious junior events. Juniors may earn income from tennis by participating in the Future, Satellite, or Challenger tours. Tournaments are broken up into different tiers offering different amounts of ranking points, culminating with Grade A.
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Leading juniors are allowed to participate for their nation in the Junior Fed Cup and Davis Cup competitions. To succeed in tennis often means having to begin playing at a young age. To facilitate and nurture a junior's growth in tennis, almost all tennis playing nations have developed a junior development system. Juniors develop their play through a range of tournaments on all surfaces, accommodating all different standards of play. Talented juniors may also receive sponsorships from governing bodies or private institutions.
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A tennis match is intended to be continuous.[63] Because stamina is a relevant factor, arbitrary delays are not permitted. In most cases, service is required to occur no more than 20 seconds after the end of the previous point.[63] This is increased to 90 seconds when the players change ends (after every odd-numbered game), and a 2-minute break is permitted between sets.[63] Other than this, breaks are permitted only when forced by events beyond the players' control, such as rain, damaged footwear, damaged racket, or the need to retrieve an errant ball. Should a player be deemed to be stalling repeatedly, the chair umpire may initially give a warning followed by subsequent penalties of "point", "game", and default of the match for the player who is consistently taking longer than the allowed time limit.[64]
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In the event of a rain delay, darkness or other external conditions halting play, the match is resumed at a later time, with the same score as at the time of the delay, and each player at the same end of the court as when rain halted play, or as close to the same relative compass point if play is resumed on a different court.
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Balls wear out quickly in serious play and, therefore, in ATP and WTA tournaments, they are changed after every nine games with the first change occurring after only seven games, because the first set of balls is also used for the pre-match warm-up.[43] In ITF tournaments like Fed Cup, the balls are changed after every eleven games (rather than nine) with the first change occurring after only nine games (instead of seven). An exception is that a ball change may not take place at the beginning of a tiebreaker, in which case the ball change is delayed until the beginning of the second game of the next set.[47] As a courtesy to the receiver, the server will often signal to the receiver before the first serve of the game in which new balls are used as a reminder that they are using new balls. Continuity of the balls' condition is considered part of the game, so if a re-warm-up is required after an extended break in play (usually due to rain), then the re-warm-up is done using a separate set of balls, and use of the match balls is resumed only when play resumes.
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A recent rule change is to allow coaching on court on a limited basis during a match.[65][66][67][68] This has been introduced in women's tennis for WTA Tour events in 2009 and allows the player to request her coach once per set.[69]
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Stance refers to the way a player prepares themselves in order to best be able to return a shot. Essentially, it enables them to move quickly in order to achieve a particular stroke. There are four main stances in modern tennis: open, semi-open, closed, and neutral. All four stances involve the player crouching in some manner: as well as being a more efficient striking posture, it allows them to isometrically preload their muscles in order to play the stroke more dynamically. What stance is selected is strongly influenced by shot selection. A player may quickly alter their stance depending on the circumstances and the type of shot they intend to play. Any given stance also alters dramatically based upon the actual playing of the shot with dynamic movements and shifts of body weight occurring.[70][71]
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This is the most common stance in tennis. The player’s feet are placed parallel to the net. They may be pointing sideways, directly at the net or diagonally towards it. This stance allows for a high degree of torso rotation which can add significant power to the stroke. This process is sometimes likened to the coiling and uncoiling of a spring. i.e the torso is rotated as a means of preloading the muscular system in preparation for playing the stroke: this is the coiling phase. When the stroke is played the torso rotates to face forwards again, called uncoiling, and adds significant power to the stroke. A disadvantage of this stance is that it does not always allow ‘for proper weight transfer and maintenance of balance’[70] when making powerful strokes. It is commonly used for forehand strokes; double-handed backhands can also be made effectively from it.
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This stance is somewhere between open and closed and is a very flexible stance. The feet are aligned diagonally towards the net. It allows for a lot of shoulder rotation and the torso can be coiled, before being uncoiled into the shot in order to increase the power of the shot. It is commonly used in modern tennis especially by ‘top professional players on the forehand’.[72] Two-handed backhands can also be employed from this stance.
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The closed stance is the least commonly used of the three main stances. One foot is placed further towards the net with the other foot further from it; there is a diagonal alignment between the feet. It allows for effective torso rotation in order to increase the power of the shot. It is usually used to play backhand shots and it is rare to see forehand shots played from it. A stroke from this stance may entail the rear foot coming completely off the floor with bodyweight being transferred entirely to the front foot.[70]
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[71]
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This is sometimes also referred to as the square stance. One foot is positioned closer to the net and ahead of the other which is behind and in line with it. Both feet are aligned at a 90 degree angle to the net. The neutral stance is often taught early because ‘It allows beginners to learn about shifting weight and rotation of the body.’[71] Forehands and backhands may be made from it.[73]
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A competent tennis player has eight basic shots in his or her repertoire: the serve, forehand, backhand, volley, half-volley, overhead smash, drop shot, and lob.
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A grip is a way of holding the racket in order to hit shots during a match. The grip affects the angle of the racket face when it hits the ball and influences the pace, spin, and placement of the shot. Players use various grips during play, including the Continental (The "Handshake Grip"), Eastern (Can be either semi-eastern or full eastern. Usually used for backhands.), and Western (semi-western or full western, usually for forehand grips) grips. Most players change grips during a match depending on what shot they are hitting; for example, slice shots and serves call for a Continental grip.[74]
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A serve (or, more formally, a "service") in tennis is a shot to start a point. The serve is initiated by tossing the ball into the air and hitting it (usually near the apex of its trajectory) into the diagonally opposite service box without touching the net. The serve may be hit under- or overhand although underhand serving remains a rarity.[75] If the ball hits the net on the first serve and bounces over into the correct diagonal box then it is called a "let" and the server gets two more additional serves to get it in. There can also be a let if the server serves the ball and the receiver isn't prepared.[47] If the server misses his or her first serve and gets a let on the second serve, then they get one more try to get the serve in the box.
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Experienced players strive to master the conventional overhand serve to maximize its power and placement. The server may employ different types of serve including flat serve, topspin serve, slice serve, and kick (American twist) serve. A reverse type of spin serve is hit in a manner that spins the ball opposite the natural spin of the server, the spin direction depending upon right- or left-handedness. If the ball is spinning counterclockwise, it will curve right from the hitter's point of view and curve left if spinning clockwise.[76]
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Some servers are content to use the serve simply to initiate the point; however, advanced players often try to hit a winning shot with their serve. A winning serve that is not touched by the opponent is called an "ace".
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For a right-handed player, the forehand is a stroke that begins on the right side of the body, continues across the body as contact is made with the ball, and ends on the left side of the body. There are various grips for executing the forehand, and their popularity has fluctuated over the years. The most important ones are the continental, the eastern, the semi-western, and the western. For a number of years, the small, frail 1920s player Bill Johnston was considered by many to have had the best forehand of all time, a stroke that he hit shoulder-high using a western grip. Few top players used the western grip after the 1920s, but in the latter part of the 20th century, as shot-making techniques and equipment changed radically, the western forehand made a strong comeback and is now used by many modern players. No matter which grip is used, most forehands are generally executed with one hand holding the racket, but there have been fine players with two-handed forehands. In the 1940s and 50s, the Ecuadorian/American player Pancho Segura used a two-handed forehand to achieve a devastating effect against larger, more powerful players. Players such as Monica Seles or France's Fabrice Santoro and Marion Bartoli are also notable players known for their two-handed forehands.[77]
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For right-handed players, the backhand is a stroke that begins on the left side of their body, continues across their body as contact is made with the ball, and ends on the right side of their body. It can be executed with either one hand or with both and is generally considered more difficult to master than the forehand. For most of the 20th century, the backhand was performed with one hand, using either an eastern or a continental grip. The first notable players to use two hands were the 1930s Australians Vivian McGrath and John Bromwich, but they were lonely exceptions. The two-handed grip gained popularity in the 1970s as Björn Borg, Chris Evert, Jimmy Connors, and later Mats Wilander and Marat Safin used it to great effect, and it is now used by a large number of the world's best players, including Rafael Nadal and Serena Williams.[78]
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Two hands give the player more control, while one hand can generate a slice shot, applying backspin on the ball to produce a low trajectory bounce. Reach is also limited with the two-handed shot. The player long considered to have had the best backhand of all time, Don Budge, had a powerful one-handed stroke in the 1930s and 1940s that imparted topspin onto the ball. Ken Rosewall, another player noted for his one-handed backhand, used a very accurate slice backhand through the 1950s and 1960s. A small number of players, notably Monica Seles, use two hands on both the backhand and forehand sides.
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A volley is a shot returned to the opponent in mid-air before the ball bounces, generally performed near the net, and is usually made with a stiff-wristed punching motion to hit the ball into an open area of the opponent's court. The half volley is made by hitting the ball on the rise just after it has bounced, also generally in the vicinity of the net, and played with the racket close to the ground.[79] The swinging volley is hit out of the air as the player approaches the net. It is an offensive shot used to take preparation time away from the opponent, as it returns the ball into the opponent's court much faster than a standard volley.
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From a poor defensive position on the baseline, the lob can be used as either an offensive or defensive weapon, hitting the ball high and deep into the opponent's court to either enable the lobber to get into better defensive position or to win the point outright by hitting it over the opponent's head. If the lob is not hit deeply enough into the other court, however, an opponent near the net may then hit an overhead smash, a hard, serve-like shot, to try to end the point.
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A difficult shot in tennis is the return of an attempted lob over the backhand side of a player. When the contact point is higher than the reach of a two-handed backhand, most players will try to execute a high slice (under the ball or sideways). Fewer players attempt the backhand sky-hook or smash. Rarely, a player will go for a high topspin backhand, while themselves in the air. A successful execution of any of these alternatives requires balance and timing, with less margin of error than the lower contact point backhands, since this shot is a break in the regular pattern of play.
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If an opponent is deep in his court, a player may suddenly employ an unexpected drop shot, by softly tapping the ball just over the net so that the opponent is unable to run in fast enough to retrieve it. Advanced players will often apply back spin to a drop shot, causing the ball to "skid" upon landing and bounce sideways, with less forward momentum toward their opponent, or even backwards towards the net, thus making it even more difficult to return.
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Muscle strain is one of the most common injuries in tennis.[80] When an isolated large-energy appears during the muscle contraction and at the same time body weight apply huge amount of pressure to the lengthened muscle, muscle strain can occur.[81] Inflammation and bleeding are triggered when muscle strain occurs, which can result in redness, pain and swelling.[81] Overuse is also common in tennis players of all levels. Muscle, cartilage, nerves, bursae, ligaments and tendons may be damaged from overuse. The repetitive use of a particular muscle without time for repair and recovery is the most common cause of injury.[81]
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Tournaments are often organized by gender and number of players. Common tournament configurations include men's singles, women's singles, and doubles, where two players play on each side of the net. Tournaments may be organized for specific age groups, with upper age limits for youth and lower age limits for senior players. Example of this include the Orange Bowl and Les Petits As junior tournaments. There are also tournaments for players with disabilities, such as wheelchair tennis and deaf tennis.[82] In the four Grand Slam tournaments, the singles draws are limited to 128 players for each gender.
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Most large tournaments seed players, but players may also be matched by their skill level. According to how well a person does in sanctioned play, a player is given a rating that is adjusted periodically to maintain competitive matches. For example, the United States Tennis Association administers the National Tennis Rating Program (NTRP), which rates players between 1.0 and 7.0 in 1/2 point increments. Average club players under this system would rate 3.0–4.5 while world class players would be 7.0 on this scale.
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The four Grand Slam tournaments are considered to be the most prestigious tennis events in the world. They are held annually and comprise, in chronological order, the Australian Open, the French Open, Wimbledon, and the US Open. Apart from the Olympic Games, Davis Cup, Fed Cup, and Hopman Cup, they are the only tournaments regulated by the International Tennis Federation (ITF).[83] The ITF's national associations, Tennis Australia (Australian Open), the Fédération Française de Tennis (French Open), the Lawn Tennis Association (Wimbledon) and the United States Tennis Association (US Open) are delegated the responsibility to organize these events.[83]
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Aside from the historical significance of these events, they also carry larger prize funds than any other tour event and are worth double the number of ranking points to the champion than in the next echelon of tournaments, the Masters 1000 (men) and Premier events (women).[84][85] Another distinguishing feature is the number of players in the singles draw. There are 128, more than any other professional tennis tournament. This draw is composed of 32 seeded players, other players ranked in the world's top 100, qualifiers, and players who receive invitations through wild cards. Grand Slam men's tournaments have best-of-five set matches while the women play best-of-three. Grand Slam tournaments are among the small number of events that last two weeks, the others being the Indian Wells Masters and the Miami Masters.
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Currently, the Grand Slam tournaments are the only tour events that have mixed doubles contests. Grand Slam tournaments are held in conjunction with wheelchair tennis tournaments and junior tennis competitions. These tournaments also contain their own idiosyncrasies. For example, players at Wimbledon are required to wear predominantly white. Andre Agassi chose to skip Wimbledon from 1988 through 1990 citing the event's traditionalism, particularly its "predominantly white" dress code.[86] Wimbledon has its own particular methods for disseminating tickets, often leading tennis fans to follow complex procedures to obtain tickets.[87]
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* The international tournament began in 1925
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The ATP World Tour Masters 1000 is a group of nine tournaments that form the second-highest echelon in men's tennis. Each event is held annually, and a win at one of these events is worth 1000 ranking points. When the ATP, led by Hamilton Jordan, began running the men's tour in 1990, the directors designated the top nine tournaments, outside of the Grand Slam events, as "Super 9" events.[88] In 2000 this became the Tennis Masters Series and in 2004 the ATP Masters Series. In November at the end of the tennis year, the world's top eight players compete in the ATP World Tour Finals, a tournament with a rotating locale. It is currently held in London, England.[89]
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In August 2007 the ATP announced major changes to the tour that were introduced in 2009. The Masters Series was renamed to the "Masters 1000", the addition of the number 1000 referring to the number of ranking points earned by the winner of each tournament. Contrary to earlier plans, the number of tournaments was not reduced from nine to eight and the Monte Carlo Masters remains part of the series although, unlike the other events, it does not have a mandatory player commitment. The Hamburg Masters has been downgraded to a 500-point event. The Madrid Masters moved to May and onto clay courts, and a new tournament in Shanghai took over Madrid's former indoor October slot. As of 2011 six of the nine "1000" level tournaments are combined ATP and WTA events.[90]
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The third and fourth tier of men's tennis tournaments are formed by the ATP World Tour 500 series, consisting of 11 tournaments, and the ATP World Tour 250 series with 40 tournaments.[91] Like the ATP World Tour Masters 1000, these events offer various amounts of prize money and the numbers refer to the amount of ranking points earned by the winner of a tournament.[84] The Dubai Tennis Championships offer the largest financial incentive to players, with total prize money of US$2,313,975 (2012).[92] These series have various draws of 28, 32, 48 and 56 for singles and 16 and 24 for doubles. It is mandatory for leading players to enter at least four 500 events, including at least one after the US Open.
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The Challenger Tour for men is the lowest level of tournament administered by the ATP. It is composed of about 150 events and, as a result, features a more diverse range of countries hosting events.[93] The majority of players use the Challenger Series at the beginning of their career to work their way up the rankings. Andre Agassi, between winning Grand Slam tournaments, plummeted to World No. 141 and used Challenger Series events for match experience and to progress back up the rankings.[94] The Challenger Series offers prize funds of between US$25,000 and US$150,000.
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Below the Challenger Tour are the Futures tournaments, events on the ITF Men's Circuit. These tournaments also contribute towards a player's ATP rankings points. Futures Tournaments offer prize funds of between US$10,000 and US$15,000.[95] Approximately 530 Futures Tournaments are played each year.
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Premier events for women form the most prestigious level of events on the Women's Tennis Association Tour after the Grand Slam tournaments. These events offer the largest rewards in terms of points and prize money. Within the Premier category are Premier Mandatory, Premier 5, and Premier tournaments. The Premier events were introduced in 2009 replacing the previous Tier I and II tournament categories. Currently four tournaments are Premier Mandatory, five tournaments are Premier 5, and twelve tournaments are Premier. The first tiering system in women's tennis was introduced in 1988. At the time of its creation, only two tournaments, the Lipton International Players Championships in Florida and the German Open in Berlin, comprised the Tier I category.
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International tournaments are the second main tier of the WTA tour and consist of 31 tournaments, with a prize money for every event at U.S.$220,000, except for the year-ending Commonwealth Bank Tournament of Champions in Bali, which has prize money of U.S.$600,000.
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Professional tennis players enjoy the same relative perks as most top sports personalities: clothing, equipment and endorsements. Like players of other individual sports such as golf, they are not salaried, but must play and finish highly in tournaments to obtain prize money.
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In recent years, some controversy has surrounded the involuntary or deliberate noise caused by players' grunting.[citation needed]
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While players are gradually less competitive in singles by their late 20s and early 30s, they can still continue competitively in doubles (as instanced by Martina Navratilova and John McEnroe, who won doubles titles in their 40s).
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In the Open Era, several female players such as Martina Navratilova, Margaret Court, Martina Hingis, Serena Williams, and Venus Williams (the latter two sisters playing together) have been prolific at both singles and doubles events throughout their careers. John McEnroe is one of the very few professional male players to be top ranked in both singles and doubles at the same time,[96][97][98] and Yevgeny Kafelnikov is the most recent male player to win multiple Grand Slams in both singles and doubles during the same period of his career.
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In terms of public attention and earnings (see below), singles champions have far surpassed their doubles counterparts. The Open Era, particularly the men's side, has seen many top-ranked singles players that only sparingly compete in doubles, while having "doubles specialists" who are typically being eliminated early in the singles draw but do well in the doubles portion of a tournament. Notable doubles pairings include The Woodies (Todd Woodbridge and Mark Woodforde) and the Bryan Brothers (identical twin brothers Robert Charles "Bob" Bryan and Michael Carl "Mike" Bryan). Woodbridge has disliked the term "doubles ‘specialists’", saying that he and Woodforde "set a singles schedule and doubles fitted in around that", although later in Woodbridge's career he focused exclusively on doubles as his singles ranking fell too low that it was no longer financially viable to recover at that age. Woodbridge noted that while top singles players earn enough that they don't need to nor want to play doubles, he suggested that lower-ranked singles players outside the Top Ten should play doubles to earn more playing time and money.[99][100]
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The Olympics doubles tennis tournament necessitates that both members of a doubles pairing be from the same country, hence several top professional pairs such as Jamie Murray and Bruno Soares cannot compete in the Olympics. Top-ranked singles players that are usually rivals on the professional circuit, such as Boris Becker and Michael Stich, and Roger Federer and Stan Wawrinka have formed a rare doubles partnership for the Olympics. Unlike professional tennis tournaments (see below) where singles players receive much more prize money than doubles players, an Olympic medal for both singles and doubles has similar prestige. The Olympics is more of a priority for doubles champions while singles champions often skip the tournament.[99][100] While the ATP has voted for Olympic results to count towards player ranking points, WTA players voted against it.[101]
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For the 2000 Olympics, Lisa Raymond was passed over for Team USA in favor of Serena Williams by captain Billie Jean King, even though Raymond was the top-ranked doubles player in the world at the time, and Raymond unsuccessfully challenged the selection.[101]
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In professional tennis tournaments such as Wimbledon, the singles competition receives the most prize money and coverage, followed by doubles, and then mixed doubles usually receive the lowest monetary awards.[102] For instance in the US Open as of 2018, the men's and women's singles prize money (US$40,912,000) accounts for 80.9 percent of total player base compensation, while men's and women's doubles (US$6,140,840), men's and women's singles qualifying (US$3,008,000), and mixed doubles (US$505,000) account for 12.1 percent, 5.9 percent, and 1.0 percent, respectively. The singles winner receives US$3,800,000, while the doubles winning pair receives $700,000 and the mixed doubles winning pair receives US$155,000.[103]
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The following players have won at least five singles titles at Grand Slam tournaments:
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A frequent topic of discussion among tennis fans and commentators is who was the greatest male singles player of all time. By a large margin, an Associated Press poll in 1950 named Bill Tilden as the greatest player of the first half of the 20th century.[104] From 1920 to 1930, Tilden won singles titles at Wimbledon three times and the U.S. Championships seven times. In 1938, however, Donald Budge became the first person to win all four major singles titles during the same calendar year, the Grand Slam, and won six consecutive major titles in 1937 and 1938. Tilden called Budge "the finest player 365 days a year that ever lived."[105] In his 1979 autobiography, Jack Kramer said that, based on consistent play, Budge was the greatest player ever.[106] Some observers, however, also felt that Kramer deserved consideration for the title. Kramer was among the few who dominated amateur and professional tennis during the late 1940s and early 1950s. Tony Trabert has said that of the players he saw before the start of the Open Era, Kramer was the best male champion.[107]
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By the 1960s, Budge and others had added Pancho Gonzales and Lew Hoad to the list of contenders. Budge reportedly believed that Gonzales was the greatest player ever.[108] Gonzales said about Hoad, "When Lew's game was at its peak nobody could touch him. ... I think his game was the best game ever. Better than mine. He was capable of making more shots than anybody. His two volleys were great. His overhead was enormous. He had the most natural tennis mind with the most natural tennis physique."[109]
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Before and during the Open Era, Rod Laver remains the only male player in history to have won the calendar year Grand Slam twice in 1962 and 1969 [110] and also the calendar year Professional Grand Slam in 1967.[111] More recently Björn Borg and Pete Sampras were regarded by many of their contemporaries as among the greatest ever. Andre Agassi, the first of two male players in history to have achieved a Career Golden Slam in singles tennis (followed by Rafael Nadal), has been called the best service returner in the history of the game.[112][113][114][115] He is the first man to win grand slams on all modern surfaces (previous holders of all grand slam tournaments played in an era of grass and clay only), and is regarded by a number of critics and fellow players to be among the greatest players of all time.[112][116][117] Both Rod Laver and Ken Rosewall also won major Pro Slam tournaments on all three surfaces (grass, clay, hard court) Rosewall in 1963 and Laver in 1967.[118]
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By the early twenty-first century, Roger Federer is considered by many observers to have the most "complete" game in modern tennis. He has won 20 grand slam titles and 6 World Tour Finals, the most for any male player. Many experts of tennis, former tennis players and his own tennis peers believe Federer is the greatest player in the history of the game.[119][120][121][122][123][124] Federer's biggest rival Rafael Nadal is regarded as the greatest competitor in tennis history by some former players and is regarded to have the potential to be the greatest of all time.[125][126] Nadal is regarded as the greatest clay court player of all time.[127] Novak Djokovic, a rival of both Nadal and Federer, is also considered to be one of the greatest tennis players of all time.[128]
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As with the men there are frequent discussions about who is the greatest female singles player of all time with Steffi Graf, Martina Navratilova and Serena Williams being the three players most often nominated.
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In March 2012 the TennisChannel published a combined list of the 100 greatest men and women tennis players of all time.[129] It ranked Steffi Graf as the greatest female player (in 3rd place overall), followed by Martina Navratilova (4th place) and Margaret Court (8th place). The rankings were determined by an international panel.
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Sportswriter John Wertheim of Sports Illustrated stated in an article in July 2010 that Serena Williams is the greatest female tennis player ever with the argument that "Head-to-head, on a neutral surface (i.e. hard courts), everyone at their best, I can't help feeling that she crushes the other legends.".[130] In a reaction to this article Yahoo sports blog Busted Racket published a list of the top-10 women's tennis players of all time placing Martina Navratilova in first spot.[131] This top-10 list was similar to the one published in June 2008 by the Bleacher Report who also ranked Martina Navratilova as the top female player of all time.[132]
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Steffi Graf is considered by some to be the greatest female player. Billie Jean King said in 1999, "Steffi is definitely the greatest women's tennis player of all time."[133] Martina Navratilova has included Graf on her list of great players.[133] In December 1999, Graf was named the greatest female tennis player of the 20th century by a panel of experts assembled by the Associated Press.[134] Tennis writer Steve Flink, in his book The Greatest Tennis Matches of the Twentieth Century, named her as the best female player of the 20th century, directly followed by Martina Navratilova.[135]
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Tennis magazine selected Martina Navratilova as the greatest female tennis player for the years 1965 through 2005.[136][137] Tennis historian and journalist Bud Collins has called Navratilova "arguably, the greatest player of all time."[138] Billie Jean King said about Navratilova in 2006, "She's the greatest singles, doubles and mixed doubles player who's ever lived."[139]
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Baseball is a bat-and-ball game played between two opposing teams who take turns batting and fielding. The game proceeds when a player on the fielding team, called the pitcher, throws a ball which a player on the batting team tries to hit with a bat. The objective of the offensive team (batting team) is to hit the ball into the field of play, allowing its players to run the bases, having them advance counter-clockwise around four bases to score what are called "runs". The objective of the defensive team (fielding team) is to prevent batters from becoming runners, and to prevent runners' advance around the bases.[2] A run is scored when a runner legally advances around the bases in order and touches home plate (the place where the player started as a batter). The team that scores the most runs by the end of the game is the winner.
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The first objective of the batting team is to have a player reach first base safely. A player on the batting team who reaches first base without being called "out" can attempt to advance to subsequent bases as a runner, either immediately or during teammates' turns batting. The fielding team tries to prevent runs by getting batters or runners "out", which forces them out of the field of play. Both the pitcher and fielders have methods of getting the batting team's players out. The opposing teams switch back and forth between batting and fielding; the batting team's turn to bat is over once the fielding team records three outs. One turn batting for each team constitutes an inning. A game is usually composed of nine innings, and the team with the greater number of runs at the end of the game wins. If scores are tied at the end of nine innings, extra innings are usually played. Baseball has no game clock, although most games end in the ninth inning.
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Baseball evolved from older bat-and-ball games already being played in England by the mid-18th century. This game was brought by immigrants to North America, where the modern version developed. By the late 19th century, baseball was widely recognized as the national sport of the United States. Baseball is popular in North America and parts of Central and South America, the Caribbean, and East Asia, particularly in Japan, South Korea and Taiwan.
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In the United States and Canada, professional Major League Baseball (MLB) teams are divided into the National League (NL) and American League (AL), each with three divisions: East, West, and Central. The MLB champion is determined by playoffs that culminate in the World Series. The top level of play is similarly split in Japan between the Central and Pacific Leagues and in Cuba between the West League and East League. The World Baseball Classic, organized by the World Baseball Softball Confederation, is the major international competition of the sport and attracts the top national teams from around the world.
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A baseball game is played between two teams, each composed of nine players, that take turns playing offense (batting and baserunning) and defense (pitching and fielding). A pair of turns, one at bat and one in the field, by each team constitutes an inning. A game consists of nine innings (seven innings at the high school level and in doubleheaders in college and minor leagues, and six innings at the Little League level).[3] One team—customarily the visiting team—bats in the top, or first half, of every inning. The other team—customarily the home team—bats in the bottom, or second half, of every inning. The goal of the game is to score more points (runs) than the other team. The players on the team at bat attempt to score runs by touching all four bases, in order, set at the corners of the square-shaped baseball diamond. A player bats at home plate and must attempt to safely reach a base before proceeding, counterclockwise, from first base, to second base, third base, and back home to score a run. The team in the field attempts to prevent runs from scoring by recording outs, which remove opposing players from offensive action, until their next turn at bat comes up again. When three outs are recorded, the teams switch roles for the next half-inning. If the score of the game is tied after nine innings, extra innings are played to resolve the contest. Many amateur games, particularly unorganized ones, involve different numbers of players and innings.[4]
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The game is played on a field whose primary boundaries, the foul lines, extend forward from home plate at 45-degree angles. The 90-degree area within the foul lines is referred to as fair territory; the 270-degree area outside them is foul territory. The part of the field enclosed by the bases and several yards beyond them is the infield; the area farther beyond the infield is the outfield. In the middle of the infield is a raised pitcher's mound, with a rectangular rubber plate (the rubber) at its center. The outer boundary of the outfield is typically demarcated by a raised fence, which may be of any material and height. The fair territory between home plate and the outfield boundary is baseball's field of play, though significant events can take place in foul territory, as well.[5]
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There are three basic tools of baseball: the ball, the bat, and the glove or mitt:
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Protective helmets are also standard equipment for all batters.[9]
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At the beginning of each half-inning, the nine players of the fielding team arrange themselves around the field. One of them, the pitcher, stands on the pitcher's mound. The pitcher begins the pitching delivery with one foot on the rubber, pushing off it to gain velocity when throwing toward home plate. Another fielding team player, the catcher, squats on the far side of home plate, facing the pitcher. The rest of the fielding team faces home plate, typically arranged as four infielders—who set up along or within a few yards outside the imaginary lines (basepaths) between first, second, and third base—and three outfielders. In the standard arrangement, there is a first baseman positioned several steps to the left of first base, a second baseman to the right of second base, a shortstop to the left of second base, and a third baseman to the right of third base. The basic outfield positions are left fielder, center fielder, and right fielder. With the exception of the catcher, all fielders are required to be in fair territory when the pitch is delivered. A neutral umpire sets up behind the catcher.[10] Other umpires will be distributed around the field as well.[11]
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Play starts with a member of the batting team, the batter, standing at home plate, holding a bat.[12] The batter waits for the pitcher to throw a pitch (the ball) toward home plate, and attempts to hit the ball[13] with the bat.[12] The catcher catches pitches that the batter does not hit—as a result of either electing not to swing or failing to connect—and returns them to the pitcher. A batter who hits the ball into the field of play must drop the bat and begin running toward first base, at which point the player is referred to as a runner (or, until the play is over, a batter-runner). A batter-runner who reaches first base without being put out is said to be safe and is on base. A batter-runner may choose to remain at first base or attempt to advance to second base or even beyond—however far the player believes can be reached safely. A player who reaches base despite proper play by the fielders has recorded a hit. A player who reaches first base safely on a hit is credited with a single. If a player makes it to second base safely as a direct result of a hit, it is a double; third base, a triple. If the ball is hit in the air within the foul lines over the entire outfield (and outfield fence, if there is one), or otherwise safely circles all the bases, it is a home run: the batter and any runners on base may all freely circle the bases, each scoring a run. This is the most desirable result for the batter. A player who reaches base due to a fielding mistake is not credited with a hit—instead, the responsible fielder is charged with an error.[12]
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Any runners already on base may attempt to advance on batted balls that land, or contact the ground, in fair territory, before or after the ball lands. A runner on first base must attempt to advance if a ball lands in play. If a ball hit into play rolls foul before passing through the infield, it becomes dead and any runners must return to the base they occupied when the play began. If the ball is hit in the air and caught before it lands, the batter has flied out and any runners on base may attempt to advance only if they tag up (contact the base they occupied when the play began, as or after the ball is caught). Runners may also attempt to advance to the next base while the pitcher is in the process of delivering the ball to home plate; a successful effort is a stolen base.[14]
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A pitch that is not hit into the field of play is called either a strike or a ball. A batter against whom three strikes are recorded strikes out. A batter against whom four balls are recorded is awarded a base on balls or walk, a free advance to first base. (A batter may also freely advance to first base if the batter's body or uniform is struck by a pitch outside the strike zone, provided the batter does not swing and attempts to avoid being hit.)[15] Crucial to determining balls and strikes is the umpire's judgment as to whether a pitch has passed through the strike zone, a conceptual area above home plate extending from the midpoint between the batter's shoulders and belt down to the hollow of the knee.[16]
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While the team at bat is trying to score runs, the team in the field is attempting to record outs. In addition to the strikeout, common ways a member of the batting team may be put out include the flyout, ground out, force out, and tag out. It is possible to record two outs in the course of the same play. This is called a double play. Three outs in one play, a triple play, is possible, though rare. Players put out or retired must leave the field, returning to their team's dugout or bench. A runner may be stranded on base when a third out is recorded against another player on the team. Stranded runners do not benefit the team in its next turn at bat as every half-inning begins with the bases empty.[17]
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An individual player's turn batting or plate appearance is complete when the player reaches base, hits a home run, makes an out, or hits a ball that results in the team's third out, even if it is recorded against a teammate. On rare occasions, a batter may be at the plate when, without the batter's hitting the ball, a third out is recorded against a teammate—for instance, a runner getting caught stealing (tagged out attempting to steal a base). A batter with this sort of incomplete plate appearance starts off the team's next turn batting; any balls or strikes recorded against the batter the previous inning are erased. A runner may circle the bases only once per plate appearance and thus can score at most a single run per batting turn. Once a player has completed a plate appearance, that player may not bat again until the eight other members of the player's team have all taken their turn at bat in the batting order. The batting order is set before the game begins, and may not be altered except for substitutions. Once a player has been removed for a substitute, that player may not reenter the game. Children's games often have more lenient rules, such as Little League rules, which allow players to be substituted back into the same game.[18][3]
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If the designated hitter (DH) rule is in effect, each team has a tenth player whose sole responsibility is to bat (and run). The DH takes the place of another player—almost invariably the pitcher—in the batting order, but does not field. Thus, even with the DH, each team still has a batting order of nine players and a fielding arrangement of nine players.[19]
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The number of players on a baseball roster, or squad, varies by league and by the level of organized play. A Major League Baseball (MLB) team has a roster of 25 players with specific roles. A typical roster features the following players:[20]
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Most baseball leagues worldwide have the DH rule, including MLB's American League, Japan's Pacific League, and Caribbean professional leagues, along with major American amateur organizations.[21] The Central League in Japan and the National League do not have the rule, and high-level minor league clubs connected to National League teams are not required to field a DH.[22] In leagues that apply the designated hitter rule, a typical team has nine offensive regulars (including the DH), five starting pitchers,[23] seven or eight relievers, a backup catcher, and two or three other reserve players.[24][25]
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The manager, or head coach, oversees the team's major strategic decisions, such as establishing the starting rotation, setting the lineup, or batting order, before each game, and making substitutions during games—in particular, bringing in relief pitchers. Managers are typically assisted by two or more coaches; they may have specialized responsibilities, such as working with players on hitting, fielding, pitching, or strength and conditioning. At most levels of organized play, two coaches are stationed on the field when the team is at bat: the first base coach and third base coach, whom occupy designated coaches' boxes, just outside the foul lines. These coaches assist in the direction of baserunners, when the ball is in play, and relay tactical signals from the manager to batters and runners, during pauses in play.[26] In contrast to many other team sports, baseball managers and coaches generally wear their team's uniforms; coaches must be in uniform to be allowed on the field to confer with players during a game.[27]
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Any baseball game involves one or more umpires, who make rulings on the outcome of each play. At a minimum, one umpire will stand behind the catcher, to have a good view of the strike zone, and call balls and strikes. Additional umpires may be stationed near the other bases, thus making it easier to judge plays such as attempted force outs and tag outs. In MLB, four umpires are used for each game, one near each base. In the playoffs, six umpires are used: one at each base and two in the outfield along the foul lines.[28]
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Many of the pre-game and in-game strategic decisions in baseball revolve around a fundamental fact: in general, right-handed batters tend to be more successful against left-handed pitchers and, to an even greater degree, left-handed batters tend to be more successful against right-handed pitchers.[29] A manager with several left-handed batters in the regular lineup, who knows the team will be facing a left-handed starting pitcher, may respond by starting one or more of the right-handed backups on the team's roster. During the late innings of a game, as relief pitchers and pinch hitters are brought in, the opposing managers will often go back and forth trying to create favorable matchups with their substitutions. The manager of the fielding team trying to arrange same-handed pitcher-batter matchups and the manager of the batting team trying to arrange opposite-handed matchups. With a team that has the lead in the late innings, a manager may remove a starting position player—especially one whose turn at bat is not likely to come up again—for a more skillful fielder (known as a defensive substitution).[30]
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The tactical decision that precedes almost every play in a baseball game involves pitch selection.[31] By gripping and then releasing the baseball in a certain manner, and by throwing it at a certain speed, pitchers can cause the baseball to break to either side, or downward, as it approaches the batter; thus creating differing pitches that can be selected.[32] Among the resulting wide variety of pitches that may be thrown, the four basic types are the fastball, the changeup (or off-speed pitch), and two breaking balls—the curveball and the slider.[33] Pitchers have different repertoires of pitches they are skillful at throwing. Conventionally, before each pitch, the catcher signals the pitcher what type of pitch to throw, as well as its general vertical and/or horizontal location.[34] If there is disagreement on the selection, the pitcher may shake off the sign and the catcher will call for a different pitch.
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With a runner on base and taking a lead, the pitcher may attempt a pickoff, a quick throw to a fielder covering the base to keep the runner's lead in check or, optimally, effect a tag out.[35] Pickoff attempts, however, are subject to rules that severely restrict the pitcher's movements before and during the pickoff attempt. Violation of any one of these rules could result in the umpire calling a balk against the pitcher, which permits any runners on base to advance one base with impunity.[36] If an attempted stolen base is anticipated, the catcher may call for a pitchout, a ball thrown deliberately off the plate, allowing the catcher to catch it while standing and throw quickly to a base.[37] Facing a batter with a strong tendency to hit to one side of the field, the fielding team may employ a shift, with most or all of the fielders moving to the left or right of their usual positions. With a runner on third base, the infielders may play in, moving closer to home plate to improve the odds of throwing out the runner on a ground ball, though a sharply hit grounder is more likely to carry through a drawn-in infield.[38]
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Several basic offensive tactics come into play with a runner on first base, including the fundamental choice of whether to attempt a steal of second base. The hit and run is sometimes employed, with a skillful contact hitter, the runner takes off with the pitch, drawing the shortstop or second baseman over to second base, creating a gap in the infield for the batter to poke the ball through.[39] The sacrifice bunt, calls for the batter to focus on making soft contact with the ball, so that it rolls a short distance into the infield, allowing the runner to advance into scoring position as the batter is thrown out at first. A batter, particularly one who is a fast runner, may also attempt to bunt for a hit. A sacrifice bunt employed with a runner on third base, aimed at bringing that runner home, is known as a squeeze play.[40] With a runner on third and fewer than two outs, a batter may instead concentrate on hitting a fly ball that, even if it is caught, will be deep enough to allow the runner to tag up and score—a successful batter, in this case, gets credit for a sacrifice fly.[38] The manager will sometimes signal a batter who is ahead in the count (i.e., has more balls than strikes) to take, or not swing at, the next pitch.[41]
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The evolution of baseball from older bat-and-ball games is difficult to trace with precision. Consensus once held that today's baseball is a North American development from the older game rounders, popular among children in Great Britain and Ireland.[42][43][44] Baseball Before We Knew It: A Search for the Roots of the Game (2005), by American baseball historian David Block, suggests that the game originated in England; recently uncovered historical evidence supports this position. Block argues that rounders and early baseball were actually regional variants of each other, and that the game's most direct antecedents are the English games of stoolball and "tut-ball".[42] The earliest known reference to baseball is in a 1744 British publication, A Little Pretty Pocket-Book, by John Newbery.[45] Block discovered that the first recorded game of "Bass-Ball" took place in 1749 in Surrey, and featured the Prince of Wales as a player.[46] This early form of the game was apparently brought to Canada by English immigrants.[47]
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By the early 1830s, there were reports of a variety of uncodified bat-and-ball games recognizable as early forms of baseball being played around North America.[48] The first officially recorded baseball game on this continent was played in Beachville, Ontario, Canada, on June 4, 1838.[49] In 1845, Alexander Cartwright, a member of New York City's Knickerbocker Club, led the codification of the so-called Knickerbocker Rules,[50] which in turn were based on rules developed in 1837 by William R. Wheaton of the Gotham Club.[51] While there are reports that the New York Knickerbockers played games in 1845, the contest long recognized as the first officially recorded baseball game in U.S. history took place on June 19, 1846, in Hoboken, New Jersey: the "New York Nine" defeated the Knickerbockers, 23–1, in four innings.[52] With the Knickerbocker code as the basis, the rules of modern baseball continued to evolve over the next half-century.[53]
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In the mid-1850s, a baseball craze hit the New York metropolitan area,[54] and by 1856, local journals were referring to baseball as the "national pastime" or "national game".[55] A year later, the sport's first governing body, the National Association of Base Ball Players, was formed. In 1867, it barred participation by African Americans.[56] The more formally structured National League was founded in 1876.[57] Professional Negro leagues formed, but quickly folded.[58] In 1887, softball, under the name of indoor baseball or indoor-outdoor, was invented as a winter version of the parent game.[59] The National League's first successful counterpart, the American League, which evolved from the minor Western League, was established in 1893, and virtually all of the modern baseball rules were in place by then.[60][61]
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The National Agreement of 1903 formalized relations both between the two major leagues and between them and the National Association of Professional Base Ball Leagues, representing most of the country's minor professional leagues.[62] The World Series, pitting the two major league champions against each other, was inaugurated that fall.[63] The Black Sox Scandal of the 1919 World Series led to the formation of a new National Commission of baseball that drew the two major leagues closer together.[64] The first major league baseball commissioner, Kenesaw Mountain Landis, was elected in 1920. That year also saw the founding of the Negro National League; the first significant Negro league, it would operate until 1931. For part of the 1920s, it was joined by the Eastern Colored League.[65]
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Compared with the present, professional baseball in the early 20th century was lower-scoring, and pitchers were more dominant.[66] The so-called dead-ball era ended in the early 1920s with several changes in rule and circumstance that were advantageous to hitters. Strict new regulations governed the ball's size, shape and composition, along with a new rule officially banning the spitball and other pitches that depended on the ball being treated or roughed-up with foreign substances, resulted in a ball that traveled farther when hit.[67] The rise of the legendary player Babe Ruth, the first great power hitter of the new era, helped permanently alter the nature of the game.[68] In the late 1920s and early 1930s, St. Louis Cardinals general manager Branch Rickey invested in several minor league clubs and developed the first modern farm system.[69] A new Negro National League was organized in 1933; four years later, it was joined by the Negro American League. The first elections to the National Baseball Hall of Fame took place in 1936. In 1939, Little League Baseball was founded in Pennsylvania.[70]
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A large number of minor league teams disbanded when World War II led to a player shortage. Chicago Cubs owner Philip K. Wrigley led the formation of the All-American Girls Professional Baseball League to help keep the game in the public eye.[71] The first crack in the unwritten agreement barring blacks from white-controlled professional ball occurred in 1945: Jackie Robinson was signed by the National League's Brooklyn Dodgers and began playing for their minor league team in Montreal.[72] In 1947, Robinson broke the major leagues' color barrier when he debuted with the Dodgers.[73] Latin American players, largely overlooked before, also started entering the majors in greater numbers. In 1951, two Chicago White Sox, Venezuelan-born Chico Carrasquel and black Cuban-born Minnie Miñoso, became the first Hispanic All-Stars.[74][75] Integration proceeded slowly: by 1953, only six of the 16 major league teams had a black player on the roster.[74]
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In 1975, the union's power—and players' salaries—began to increase greatly when the reserve clause was effectively struck down, leading to the free agency system.[76] Significant work stoppages occurred in 1981 and 1994, the latter forcing the cancellation of the World Series for the first time in 90 years.[77] Attendance had been growing steadily since the mid-1970s and in 1994, before the stoppage, the majors were setting their all-time record for per-game attendance.[78][79] After play resumed in 1995, non-division-winning wild card teams became a permanent fixture of the post-season. Regular-season interleague play was introduced in 1997 and the second-highest attendance mark for a full season was set.[80] In 2000, the National and American Leagues were dissolved as legal entities. While their identities were maintained for scheduling purposes (and the designated hitter distinction), the regulations and other functions—such as player discipline and umpire supervision—they had administered separately were consolidated under the rubric of MLB.[81]
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In 2001, Barry Bonds established the current record of 73 home runs in a single season. There had long been suspicions that the dramatic increase in power hitting was fueled in large part by the abuse of illegal steroids (as well as by the dilution of pitching talent due to expansion), but the issue only began attracting significant media attention in 2002 and there was no penalty for the use of performance-enhancing drugs before 2004.[82] In 2007, Bonds became MLB's all-time home run leader, surpassing Hank Aaron, as total major league and minor league attendance both reached all-time highs.[83][84]
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Widely known as America's pastime, baseball is well established in several other countries as well. As early as 1877, a professional league, the International Association, featured teams from both Canada and the US.[85] While baseball is widely played in Canada and many minor league teams have been based in the country,[86][87] the American major leagues did not include a Canadian club until 1969, when the Montreal Expos joined the National League as an expansion team. In 1977, the expansion Toronto Blue Jays joined the American League.[88]
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In 1847, American soldiers played what may have been the first baseball game in Mexico at Parque Los Berros in Xalapa, Veracruz.[89] The first formal baseball league outside of the United States and Canada was founded in 1878 in Cuba, which maintains a rich baseball tradition. The Dominican Republic held its first islandwide championship tournament in 1912.[90] Professional baseball tournaments and leagues began to form in other countries between the world wars, including the Netherlands (formed in 1922), Australia (1934), Japan (1936), Mexico (1937), and Puerto Rico (1938).[91] The Japanese major leagues have long been considered the highest quality professional circuits outside of the United States.[92]
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After World War II, professional leagues were founded in many Latin American countries, most prominently Venezuela (1946) and the Dominican Republic (1955).[93] Since the early 1970s, the annual Caribbean Series has matched the championship clubs from the four leading Latin American winter leagues: the Dominican Professional Baseball League, Mexican Pacific League, Puerto Rican Professional Baseball League, and Venezuelan Professional Baseball League. In Asia, South Korea (1982), Taiwan (1990) and China (2003) all have professional leagues.[94]
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Many European countries have professional leagues as well; the most successful, other than the Dutch league, is the Italian league, founded in 1948.[95] In 2004, Australia won a surprise silver medal at the Olympic Games.[96] The Confédération Européene de Baseball (European Baseball Confederation), founded in 1953, organizes a number of competitions between clubs from different countries. Other competitions between national teams, such as the Baseball World Cup and the Olympic baseball tournament, were administered by the International Baseball Federation (IBAF) from its formation in 1938 until its 2013 merger with the International Softball Federation to create the current joint governing body for both sports, the World Baseball Softball Confederation (WBSC).[97] Women's baseball is played on an organized amateur basis in numerous countries.[98]
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After being admitted to the Olympics as a medal sport beginning with the 1992 Games, baseball was dropped from the 2012 Summer Olympic Games at the 2005 International Olympic Committee meeting. It remained part of the 2008 Games.[99] While the sport's lack of a following in much of the world was a factor,[100] more important was MLB's reluctance to allow its players to participate during the major league season.[101] MLB initiated the World Baseball Classic, scheduled to precede its season, partly as a replacement, high-profile international tournament. The inaugural Classic, held in March 2006, was the first tournament involving national teams to feature a significant number of MLB participants.[102][103] The Baseball World Cup was discontinued after its 2011 edition in favor of an expanded World Baseball Classic.[104]
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Baseball has certain attributes that set it apart from the other popular team sports in the countries where it has a following. All of these sports use a clock,[105] play is less individual,[106] and the variation between playing fields is not as substantial or important.[107] The comparison between cricket and baseball demonstrates that many of baseball's distinctive elements are shared in various ways with its cousin sports.[108]
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In clock-limited sports, games often end with a team that holds the lead killing the clock rather than competing aggressively against the opposing team. In contrast, baseball has no clock, thus a team cannot win without getting the last batter out and rallies are not constrained by time. At almost any turn in any baseball game, the most advantageous strategy is some form of aggressive strategy.[109] Whereas, in the case of multi-day Test and first-class cricket, the possibility of a draw often encourages a team that is batting last and well behind, to bat defensively and run out the clock, giving up any faint chance at a win, to avoid an overall loss.[110]
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While nine innings has been the standard since the beginning of professional baseball, the duration of the average major league game has increased steadily through the years. At the turn of the 20th century, games typically took an hour and a half to play. In the 1920s, they averaged just less than two hours, which eventually ballooned to 2:38 in 1960.[111] By 1997, the average American League game lasted 2:57 (National League games were about 10 minutes shorter—pitchers at the plate making for quicker outs than designated hitters).[112] In 2004, Major League Baseball declared that its goal was an average game of 2:45.[111] By 2014, though, the average MLB game took over three hours to complete.[113] The lengthening of games is attributed to longer breaks between half-innings for television commercials, increased offense, more pitching changes, and a slower pace of play with pitchers taking more time between each delivery, and batters stepping out of the box more frequently.[111][112] Other leagues have experienced similar issues. In 2008, Nippon Professional Baseball took steps aimed at shortening games by 12 minutes from the preceding decade's average of 3:18.[114]
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In 2016, the average nine-inning playoff game in Major League baseball was 3 hours and 35 minutes. This was up 10 minutes from 2015 and 21 minutes from 2014.[115]
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Although baseball is a team sport, individual players are often placed under scrutiny and pressure. In 1915, a baseball instructional manual pointed out that every single pitch, of which there are often more than two hundred in a game, involves an individual, one-on-one contest: "the pitcher and the batter in a battle of wits".[116] Pitcher, batter, and fielder all act essentially independent of each other. While coaching staffs can signal pitcher or batter to pursue certain tactics, the execution of the play itself is a series of solitary acts. If the batter hits a line drive, the outfielder is solely responsible for deciding to try to catch it or play it on the bounce and for succeeding or failing. The statistical precision of baseball is both facilitated by this isolation and reinforces it.
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Cricket is more similar to baseball than many other team sports in this regard: while the individual focus in cricket is mitigated by the importance of the batting partnership and the practicalities of tandem running, it is enhanced by the fact that a batsman may occupy the wicket for an hour or much more. There is no statistical equivalent in cricket for the fielding error and thus less emphasis on personal responsibility in this area of play.[117]
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Unlike those of most sports, baseball playing fields can vary significantly in size and shape. While the dimensions of the infield are specifically regulated, the only constraint on outfield size and shape for professional teams, following the rules of MLB and Minor League Baseball, is that fields built or remodeled since June 1, 1958, must have a minimum distance of 325 feet (99 m) from home plate to the fences in left and right field and 400 feet (122 m) to center.[118] Major league teams often skirt even this rule. For example, at Minute Maid Park, which became the home of the Houston Astros in 2000, the Crawford Boxes in left field are only 315 feet (96 m) from home plate.[119] There are no rules at all that address the height of fences or other structures at the edge of the outfield. The most famously idiosyncratic outfield boundary is the left-field wall at Boston's Fenway Park, in use since 1912: the Green Monster is 310 feet (94 m) from home plate down the line and 37 feet (11 m) tall.[120]
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Similarly, there are no regulations at all concerning the dimensions of foul territory. Thus a foul fly ball may be entirely out of play in a park with little space between the foul lines and the stands, but a foulout in a park with more expansive foul ground.[121] A fence in foul territory that is close to the outfield line will tend to direct balls that strike it back toward the fielders, while one that is farther away may actually prompt more collisions, as outfielders run full speed to field balls deep in the corner. These variations can make the difference between a double and a triple or inside-the-park home run.[122] The surface of the field is also unregulated. While the adjacent image shows a traditional field surfacing arrangement (and the one used by virtually all MLB teams with naturally surfaced fields), teams are free to decide what areas will be grassed or bare.[123] Some fields—including several in MLB—use an artificial surface, such as AstroTurf. Surface variations can have a significant effect on how ground balls behave and are fielded as well as on baserunning. Similarly, the presence of a roof (seven major league teams play in stadiums with permanent or retractable roofs) can greatly affect how fly balls are played.[124] While football and soccer players deal with similar variations of field surface and stadium covering, the size and shape of their fields are much more standardized. The area out-of-bounds on a football or soccer field does not affect play the way foul territory in baseball does, so variations in that regard are largely insignificant.[125]
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These physical variations create a distinctive set of playing conditions at each ballpark. Other local factors, such as altitude and climate, can also significantly affect play. A given stadium may acquire a reputation as a pitcher's park or a hitter's park, if one or the other discipline notably benefits from its unique mix of elements. The most exceptional park in this regard is Coors Field, home of the Colorado Rockies. Its high altitude—5,282 feet (1,610 m) above sea level—is partly responsible for giving it the strongest hitter's park effect in the major leagues due to the low air pressure.[126] Wrigley Field, home of the Chicago Cubs, is known for its fickle disposition: a hitter's park when the strong winds off Lake Michigan are blowing out, it becomes more of a pitcher's park when they are blowing in.[127] The absence of a standardized field affects not only how particular games play out, but the nature of team rosters and players' statistical records. For example, hitting a fly ball 330 feet (100 m) into right field might result in an easy catch on the warning track at one park, and a home run at another. A team that plays in a park with a relatively short right field, such as the New York Yankees, will tend to stock its roster with left-handed pull hitters, who can best exploit it. On the individual level, a player who spends most of his career with a team that plays in a hitter's park will gain an advantage in batting statistics over time—even more so if his talents are especially suited to the park.[128]
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Organized baseball lends itself to statistics to a greater degree than many other sports. Each play is discrete and has a relatively small number of possible outcomes. In the late 19th century, a former cricket player, English-born Henry Chadwick of Brooklyn, was responsible for the "development of the box score, tabular standings, the annual baseball guide, the batting average, and most of the common statistics and tables used to describe baseball."[129] The statistical record is so central to the game's "historical essence" that Chadwick came to be known as Father Baseball.[129] In the 1920s, American newspapers began devoting more and more attention to baseball statistics, initiating what journalist and historian Alan Schwarz describes as a "tectonic shift in sports, as intrigue that once focused mostly on teams began to go to individual players and their statistics lines."[130]
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The Official Baseball Rules administered by MLB require the official scorer to categorize each baseball play unambiguously. The rules provide detailed criteria to promote consistency. The score report is the official basis for both the box score of the game and the relevant statistical records.[131] General managers, managers, and baseball scouts use statistics to evaluate players and make strategic decisions.
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Certain traditional statistics are familiar to most baseball fans. The basic batting statistics include:[132]
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The basic baserunning statistics include:[133]
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The basic pitching statistics include:[134]
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The basic fielding statistics include:[135]
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Among the many other statistics that are kept are those collectively known as situational statistics. For example, statistics can indicate which specific pitchers a certain batter performs best against. If a given situation statistically favors a certain batter, the manager of the fielding team may be more likely to change pitchers or have the pitcher intentionally walk the batter in order to face one who is less likely to succeed.[136]
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Sabermetrics refers to the field of baseball statistical study and the development of new statistics and analytical tools. The term is also used to refer directly to new statistics themselves. The term was coined around 1980 by one of the field's leading proponents, Bill James, and derives from the Society for American Baseball Research (SABR).[137]
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The growing popularity of sabermetrics since the early 1980s has brought more attention to two batting statistics that sabermetricians argue are much better gauges of a batter's skill than batting average:[138]
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Some of the new statistics devised by sabermetricians have gained wide use:
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Writing in 1919, philosopher Morris Raphael Cohen described baseball as America's national religion.[143] In the words of sports columnist Jayson Stark, baseball has long been "a unique paragon of American culture"—a status he sees as devastated by the steroid abuse scandal.[144] Baseball has an important place in other national cultures as well: Scholar Peter Bjarkman describes "how deeply the sport is ingrained in the history and culture of a nation such as Cuba, [and] how thoroughly it was radically reshaped and nativized in Japan."[145] Since the early 1980s, the Dominican Republic, in particular the city of San Pedro de Macorís, has been the major leagues' primary source of foreign talent.[146] In 2017, 83 of the 868 players on MLB Opening Day rosters (and disabled lists) were from the country. Among other Caribbean countries and territories, a combined 97 MLB players were born in Venezuela, Cuba, and Puerto Rico.[147] Hall-of-Famer Roberto Clemente remains one of the greatest national heroes in Puerto Rico's history.[148] While baseball has long been the island's primary athletic pastime, its once well-attended professional winter league has declined in popularity since 1990, when young Puerto Rican players began to be included in the major leagues' annual first-year player draft.[149] In Asia, baseball is among the most popular sports in Japan and South Korea.[150]
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The major league game in the United States was originally targeted toward a middle-class, white-collar audience: relative to other spectator pastimes, the National League's set ticket price of 50 cents in 1876 was high, while the location of playing fields outside the inner city and the workweek daytime scheduling of games were also obstacles to a blue-collar audience.[151] A century later, the situation was very different. With the rise in popularity of other team sports with much higher average ticket prices—football, basketball, and hockey—professional baseball had become among the most blue-collar-oriented of leading American spectator sports.[152]
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Overall, baseball has a large following in the United States; a 2006 poll found that nearly half of Americans are fans.[153] In the late 1900s and early 2000s, baseball's position compared to football in the United States moved in contradictory directions. In 2008, MLB set a revenue record of $6.5 billion, matching the NFL's revenue for the first time in decades.[154] A new MLB revenue record of more than $10 billion was set in 2017.[155] On the other hand, the percentage of American sports fans polled who named baseball as their favorite sport was 9%, compared to pro football at 37%.[156] In 1985, the respective figures were pro football 24%, baseball 23%.[157] Because there are so many more major league games played, there is no comparison in overall attendance.[158] In 2008, total attendance at major league games was the second-highest in history: 78.6 million, 0.7% off the record set the previous year.[83] The following year, amid the U.S. recession, attendance fell by 6.6% to 73.4 million.[159] Eight years later, it dropped under 73 million.[160] Attendance at games held under the Minor League Baseball umbrella set a record in 2008, with 43.3 million.[161] While MLB games have not drawn the same national TV viewership as football games, MLB games are dominant in teams' local markets and regularly lead all programs in primetime in their markets during the summer.[162]
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In Japan, where baseball is inarguably the leading spectator team sport, combined revenue for the twelve teams in Nippon Professional Baseball (NPB), the body that oversees both the Central and Pacific Leagues, was estimated at $1 billion in 2007. Total NPB attendance for the year was approximately 20 million. While in the preceding two decades, MLB attendance grew by 50 percent and revenue nearly tripled, the comparable NPB figures were stagnant. There are concerns that MLB's growing interest in acquiring star Japanese players will hurt the game in their home country.[163] In Cuba, where baseball is by every reckoning the national sport,[164] the national team overshadows the city and provincial teams that play in the top-level domestic leagues.[165] Revenue figures are not released for the country's amateur system. Similarly, according to one official pronouncement, the sport's governing authority "has never taken into account attendance ... because its greatest interest has always been the development of athletes".[166]
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As of 2018[update], Little League Baseball oversees leagues with close to 2.4 million participants in over 80 countries.[167] The number of players has fallen since the 1990s, when 3 million children took part in Little League Baseball annually.[168] Babe Ruth League teams have over 1 million participants.[169] According to the president of the International Baseball Federation, between 300,000 and 500,000 women and girls play baseball around the world, including Little League and the introductory game of Tee Ball.[170]
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A varsity baseball team is an established part of physical education departments at most high schools and colleges in the United States.[171] In 2015, nearly half a million high schoolers and over 34,000 collegians played on their schools' baseball teams.[172] By early in the 20th century, intercollegiate baseball was Japan's leading sport. Today, high school baseball in particular is immensely popular there.[173] The final rounds of the two annual tournaments—the National High School Baseball Invitational Tournament in the spring, and the even more important National High School Baseball Championship in the summer—are broadcast around the country. The tournaments are known, respectively, as Spring Koshien and Summer Koshien after the 55,000-capacity stadium where they are played.[174] In Cuba, baseball is a mandatory part of the state system of physical education, which begins at age six. Talented children as young as seven are sent to special district schools for more intensive training—the first step on a ladder whose acme is the national baseball team.[165]
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Baseball has had a broad impact on popular culture, both in the United States and elsewhere. Dozens of English-language idioms have been derived from baseball; in particular, the game is the source of a number of widely used sexual euphemisms.[176] The first networked radio broadcasts in North America were of the 1922 World Series: famed sportswriter Grantland Rice announced play-by-play from New York City's Polo Grounds on WJZ–Newark, New Jersey, which was connected by wire to WGY–Schenectady, New York, and WBZ–Springfield, Massachusetts.[177] The baseball cap has become a ubiquitous fashion item not only in the United States and Japan, but also in countries where the sport itself is not particularly popular, such as the United Kingdom.[178]
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Baseball has inspired many works of art and entertainment. One of the first major examples, Ernest Thayer's poem "Casey at the Bat", appeared in 1888. A wry description of the failure of a star player in what would now be called a "clutch situation", the poem became the source of vaudeville and other staged performances, audio recordings, film adaptations, and an opera, as well as a host of sequels and parodies in various media. There have been many baseball movies, including the Academy Award–winning The Pride of the Yankees (1942) and the Oscar nominees The Natural (1984) and Field of Dreams (1989). The American Film Institute's selection of the ten best sports movies includes The Pride of the Yankees at number 3 and Bull Durham (1988) at number 5.[179] Baseball has provided thematic material for hits on both stage—the Adler–Ross musical Damn Yankees—and record—George J. Gaskin's "Slide, Kelly, Slide", Simon and Garfunkel's "Mrs. Robinson", and John Fogerty's "Centerfield".[180] The baseball-inspired comedic sketch "Who's on First", popularized by Abbott and Costello in 1938, quickly became famous. Six decades later, Time named it the best comedy routine of the 20th century.[181]
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Literary works connected to the game include the short fiction of Ring Lardner and novels such as Bernard Malamud's The Natural (the source for the movie), Robert Coover's The Universal Baseball Association, Inc., J. Henry Waugh, Prop., and W. P. Kinsella's Shoeless Joe (the source for Field of Dreams). Baseball's literary canon also includes the beat reportage of Damon Runyon; the columns of Grantland Rice, Red Smith, Dick Young, and Peter Gammons; and the essays of Roger Angell. Among the celebrated nonfiction books in the field are Lawrence S. Ritter's The Glory of Their Times, Roger Kahn's The Boys of Summer, and Michael Lewis's Moneyball. The 1970 publication of major league pitcher Jim Bouton's tell-all chronicle Ball Four is considered a turning point in the reporting of professional sports.[182]
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Baseball has also inspired the creation of new cultural forms. Baseball cards were introduced in the late 19th century as trade cards. A typical example featured an image of a baseball player on one side and advertising for a business on the other. In the early 1900s they were produced widely as promotional items by tobacco and confectionery companies. The 1930s saw the popularization of the modern style of baseball card, with a player photograph accompanied on the rear by statistics and biographical data. Baseball cards—many of which are now prized collectibles—are the source of the much broader trading card industry, involving similar products for different sports and non-sports-related fields.[183]
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Modern fantasy sports began in 1980 with the invention of Rotisserie League Baseball by New York writer Daniel Okrent and several friends. Participants in a Rotisserie league draft notional teams from the list of active MLB players and play out an entire imaginary season with game outcomes based on the players' latest real-world statistics. Rotisserie-style play quickly became a phenomenon. Now known more generically as fantasy baseball, it has inspired similar games based on an array of different sports.[184] The field boomed with increasing Internet access and new fantasy sports-related websites. By 2008, 29.9 million people in the United States and Canada were playing fantasy sports, spending $800 million on the hobby.[185] The burgeoning popularity of fantasy baseball is also credited with the increasing attention paid to sabermetrics—first among fans, only later among baseball professionals.[186]
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Tennis is a racket sport that can be played individually against a single opponent (singles) or between two teams of two players each (doubles). Each player uses a tennis racket that is strung with cord to strike a hollow rubber ball covered with felt over or around a net and into the opponent's court. The object of the game is to maneuver the ball in such a way that the opponent is not able to play a valid return. The player who is unable to return the ball will not gain a point, while the opposite player will.
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Tennis is an Olympic sport and is played at all levels of society and at all ages. The sport can be played by anyone who can hold a racket, including wheelchair users. The modern game of tennis originated in Birmingham, England, in the late 19th century as lawn tennis.[1] It had close connections both to various field (lawn) games such as croquet and bowls as well as to the older racket sport today called real tennis. During most of the 19th century, in fact, the term tennis referred to real tennis, not lawn tennis.
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The rules of modern tennis have changed little since the 1890s. Two exceptions are that from 1908 to 1961 the server had to keep one foot on the ground at all times, and the adoption of the tiebreak in the 1970s. A recent addition to professional tennis has been the adoption of electronic review technology coupled with a point-challenge system, which allows a player to contest the line call of a point, a system known as Hawk-Eye.
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Tennis is played by millions of recreational players and is also a popular worldwide spectator sport. The four Grand Slam tournaments (also referred to as the Majors) are especially popular: the Australian Open played on hard courts, the French Open played on red clay courts, Wimbledon played on grass courts, and the US Open also played on hard courts.
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Historians believe that the game's ancient origin lay in 12th century northern France, where a ball was struck with the palm of the hand.[2] Louis X of France was a keen player of jeu de paume ("game of the palm"), which evolved into real tennis, and became notable as the first person to construct indoor tennis courts in the modern style. Louis was unhappy with playing tennis outdoors and accordingly had indoor, enclosed courts made in Paris "around the end of the 13th century".[3] In due course this design spread across royal palaces all over Europe.[3] In June 1316 at Vincennes, Val-de-Marne and following a particularly exhausting game, Louis drank a large quantity of cooled wine and subsequently died of either pneumonia or pleurisy, although there was also suspicion of poisoning.[4] Because of the contemporary accounts of his death, Louis X is history's first tennis player known by name.[4] Another of the early enthusiasts of the game was King Charles V of France, who had a court set up at the Louvre Palace.[5]
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It was not until the 16th century that rackets came into use and the game began to be called "tennis", from the French term tenez, which can be translated as "hold!", "receive!" or "take!", an interjection used as a call from the server to his opponent.[6] It was popular in England and France, although the game was only played indoors where the ball could be hit off the wall. Henry VIII of England was a big fan of this game, which is now known as real tennis.[7] During the 18th and early 19th centuries, as real tennis declined, new racket sports emerged in England.[8]
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The invention of the first lawn mower in 1830, in Britain, is believed to have been a catalyst, for the preparation of modern-style grass courts, sporting ovals, playing fields, pitches, greens, etc. This in turn led to the codification of modern rules for many sports, including lawn tennis, most football codes, lawn bowls and others.[9]
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Between 1859 and 1865 Harry Gem, a solicitor and his friend Augurio Perera developed a game that combined elements of racquets and the Basque ball game pelota, which they played on Perera's croquet lawn in Birmingham in England.[10][11] In 1872, along with two local doctors, they founded the world's first tennis club on Avenue Road, Leamington Spa.[12] This is where "lawn tennis" was used as a name of activity by a club for the first time. After Leamington, the second club to take up the game of lawn tennis appears to have been the Edgbaston Archery and Croquet Society, also in Birmingham.
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In Tennis: A Cultural History, Heiner Gillmeister reveals that on December 8, 1874, British army officer Walter Clopton Wingfield wrote to Harry Gem, commenting that he (Wingfield) had been experimenting with his version of lawn tennis “for a year and a half”.[13] In December 1873, Wingfield designed and patented a game which he called sphairistikè (Greek: σφαιριστική, meaning "ball-playing"), and was soon known simply as "sticky" – for the amusement of guests at a garden party on his friend's estate of Nantclwyd Hall, in Llanelidan, Wales.[14] According to R. D. C. Evans, turfgrass agronomist, "Sports historians all agree that [Wingfield] deserves much of the credit for the development of modern tennis."[8][15] According to Honor Godfrey, museum curator at Wimbledon, Wingfield "popularized this game enormously. He produced a boxed set which included a net, poles, rackets, balls for playing the game – and most importantly you had his rules. He was absolutely terrific at marketing and he sent his game all over the world. He had very good connections with the clergy, the law profession, and the aristocracy and he sent thousands of sets out in the first year or so, in 1874."[16] The world's oldest annual tennis tournament took place at Leamington Lawn Tennis Club in Birmingham in 1874.[17] This was three years before the All England Lawn Tennis and Croquet Club would hold its first championships at Wimbledon, in 1877. The first Championships culminated in a significant debate on how to standardise the rules.[16]
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In the U.S. in 1874 Mary Ewing Outerbridge, a young socialite, returned from Bermuda with a sphairistikè set. She became fascinated by the game of tennis after watching British army officers play.[18] She laid out a tennis court at the Staten Island Cricket Club at Camp Washington, Tompkinsville, Staten Island, New York. The first American National championship was played there in September 1880. An Englishman named O.E. Woodhouse won the singles title, and a silver cup worth $100, by defeating Canadian I. F. Hellmuth.[19] There was also a doubles match which was won by a local pair. There were different rules at each club. The ball in Boston was larger than the one normally used in New York.
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On 21 May 1881, the oldest nationwide tennis organization in the world[20] was formed, the United States National Lawn Tennis Association (now the United States Tennis Association) in order to standardize the rules and organize competitions.[21] The U.S. National Men's Singles Championship, now the US Open, was first held in 1881 at the Newport Casino, Newport, Rhode Island.[22] The U.S. National Women's Singles Championships were first held in 1887 in Philadelphia.[23]
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Tennis also became popular in France, where the French Championships dates to 1891 although until 1925 it was open only to tennis players who were members of French clubs.[24] Thus, Wimbledon, the US Open, the French Open, and the Australian Open (dating to 1905) became and have remained the most prestigious events in tennis.[25][26] Together these four events are called the Majors or Slams (a term borrowed from bridge rather than baseball).[27]
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In 1913, the International Lawn Tennis Federation (ILTF), now the International Tennis Federation (ITF), was founded and established three official tournaments as the major championships of the day. The World Grass Court Championships were awarded to Great Britain. The World Hard Court Championships were awarded to France; the term "hard court" was used for clay courts at the time. Some tournaments were held in Belgium instead. And the World Covered Court Championships for indoor courts was awarded annually; Sweden, France, Great Britain, Denmark, Switzerland and Spain each hosted the tournament.[28] At a meeting held on 16 March 1923 in Paris, the title 'World Championship' was dropped and a new category of Official Championship was created for events in Great Britain, France, the United States, and Australia – today's Grand Slam events.[28][29] The impact on the four recipient nations to replace the ‘world championships’ with ‘official
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championships’ was simple in a general sense: each became a major nation of the federation with enhanced voting power and each now operated a major event.[28]
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The comprehensive rules promulgated in 1924 by the ILTF, have remained largely stable in the ensuing eighty years, the one major change being the addition of the tiebreak system designed by Jimmy Van Alen.[30] That same year, tennis withdrew from the Olympics after the 1924 Games but returned 60 years later as a 21-and-under demonstration event in 1984. This reinstatement was credited by the efforts by the then ITF President Philippe Chatrier, ITF General Secretary David Gray and ITF Vice President Pablo Llorens, and support from IOC President Juan Antonio Samaranch. The success of the event was overwhelming and the IOC decided to reintroduce tennis as a full medal sport at Seoul in 1988.[31][32]
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The Davis Cup, an annual competition between men's national teams, dates to 1900.[33] The analogous competition for women's national teams, the Fed Cup, was founded as the Federation Cup in 1963 to celebrate the 50th anniversary of the founding of the ITF.[34]
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In 1926, promoter C. C. Pyle established the first professional tennis tour with a group of American and French tennis players playing exhibition matches to paying audiences.[26][35] The most notable of these early professionals were the American Vinnie Richards and the Frenchwoman Suzanne Lenglen.[26][36] Once a player turned pro he or she was no longer permitted to compete in the major (amateur) tournaments.[26]
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In 1968, commercial pressures and rumors of some amateurs taking money under the table led to the abandonment of this distinction, inaugurating the Open Era, in which all players could compete in all tournaments, and top players were able to make their living from tennis. With the beginning of the Open Era, the establishment of an international professional tennis circuit, and revenues from the sale of television rights, tennis's popularity has spread worldwide, and the sport has shed its middle-class English-speaking image[37] (although it is acknowledged that this stereotype still exists).[37][38]
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In 1954, Van Alen founded the International Tennis Hall of Fame, a non-profit museum in Newport, Rhode Island.[39] The building contains a large collection of tennis memorabilia as well as a hall of fame honouring prominent members and tennis players from all over the world. Each year, a grass court tournament and an induction ceremony honoring new Hall of Fame members are hosted on its grounds.
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Part of the appeal of tennis stems from the simplicity of equipment required for play. Beginners need only a racket and balls.
|
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The components of a tennis racket include a handle, known as the grip, connected to a neck which joins a roughly elliptical frame that holds a matrix of tightly pulled strings. For the first 100 years of the modern game, rackets were made of wood and of standard size, and strings were of animal gut. Laminated wood construction yielded more strength in rackets used through most of the 20th century until first metal and then composites of carbon graphite, ceramics, and lighter metals such as titanium were introduced. These stronger materials enabled the production of oversized rackets that yielded yet more power. Meanwhile, technology led to the use of synthetic strings that match the feel of gut yet with added durability.
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Under modern rules of tennis, the rackets must adhere to the following guidelines;[40]
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The rules regarding rackets have changed over time, as material and engineering advances have been made. For example, the maximum length of the frame had been 32 inches (81 cm) until 1997, when it was shortened to 29 inches (74 cm).[41]
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Many companies manufacture and distribute tennis rackets. Wilson, Head and Babolat are some of the more commonly used brands; however, many more companies exist.[example needed] The same companies sponsor players to use these rackets in the hopes that the company name will become more well known by the public.
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Tennis balls were originally made of cloth strips stitched together with thread and stuffed with feathers.[42] Modern tennis balls are made of hollow vulcanized rubber with a felt coating. Traditionally white, the predominant colour was gradually changed to optic yellow in the latter part of the 20th century to allow for improved visibility. Tennis balls must conform to certain criteria for size, weight, deformation, and bounce to be approved for regulation play. The International Tennis Federation (ITF) defines the official diameter as 65.41–68.58 mm (2.575–2.700 in). Balls must weigh between 56.0 and 59.4 g (1.98 and 2.10 oz).[43] Tennis balls were traditionally manufactured in the United States and Europe. Although the process of producing the balls has remained virtually unchanged for the past 100 years, the majority of manufacturing now takes place in the Far East. The relocation is due to cheaper labour costs and materials in the region.[44] Tournaments that are played under the ITF Rules of Tennis must use balls that are approved by the International Tennis Federation (ITF) and be named on the official ITF list of approved tennis balls.[45]
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Advanced players improve their performance through a number of accoutrements. Vibration dampeners may be interlaced in the proximal part of the string array for improved feel. Racket handles may be customized with absorbent or rubber-like materials to improve the players' grip. Players often use sweat bands on their wrists to keep their hands dry and head bands or bandanas to keep the sweat out of their eyes as well. Finally, although the game can be played in a variety of shoes, specialized tennis shoes have wide, flat soles for stability and a built-up front structure to avoid excess wear.
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Tennis is played on a rectangular, flat surface. The court is 78 feet (23.77 m) long, and 27 feet (8.2 m) wide for singles matches and 36 ft (11 m) for doubles matches.[46] Additional clear space around the court is required in order for players to reach overrun balls. A net is stretched across the full width of the court, parallel with the baselines, dividing it into two equal ends. It is held up by either a cord or metal cable of diameter no greater than 0.8 cm (1⁄3 in).[47] The net is 3 feet 6 inches (1.07 m) high at the posts and 3 feet (0.91 m) high in the center.[46] The net posts are 3 feet (0.91 m) outside the doubles court on each side or, for a singles net, 3 feet (0.91 m) outside the singles court on each side.
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The modern tennis court owes its design to Major Walter Clopton Wingfield. In 1873, Wingfield patented a court much the same as the current one for his stické tennis (sphairistike). This template was modified in 1875 to the court design that exists today, with markings similar to Wingfield's version, but with the hourglass shape of his court changed to a rectangle.[48]
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Tennis is unusual in that it is played on a variety of surfaces.[49] Grass, clay, and hardcourts of concrete or asphalt topped with acrylic are the most common. Occasionally carpet is used for indoor play, with hardwood flooring having been historically used. Artificial turf courts can also be found.
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The lines that delineate the width of the court are called the baseline (farthest back) and the service line (middle of the court). The short mark in the center of each baseline is referred to as either the hash mark or the center mark. The outermost lines that make up the length are called the doubles sidelines; they are the boundaries for doubles matches. The lines to the inside of the doubles sidelines are the singles sidelines, and are the boundaries in singles play. The area between a doubles sideline and the nearest singles sideline is called the doubles alley, playable in doubles play. The line that runs across the center of a player's side of the court is called the service line because the serve must be delivered into the area between the service line and the net on the receiving side. Despite its name, this is not where a player legally stands when making a serve.[50]
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The line dividing the service line in two is called the center line or center service line. The boxes this center line creates are called the service boxes; depending on a player's position, they have to hit the ball into one of these when serving.[51] A ball is out only if none of it has hit the area inside the lines, or the line, upon its first bounce. All lines are required to be between 1 and 2 inches (25 and 51 mm) in width, with the exception of the baseline which can be up to 4 inches (100 mm) wide, although in practice it is often the same width as the others.[50]
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The players or teams start on opposite sides of the net. One player is designated the server, and the opposing player is the receiver. The choice to be server or receiver in the first game and the choice of ends is decided by a coin toss before the warm-up starts. Service alternates game by game between the two players or teams. For each point, the server starts behind the baseline, between the center mark and the sideline. The receiver may start anywhere on their side of the net. When the receiver is ready, the server will serve, although the receiver must play to the pace of the server.
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For a service to be legal, the ball must travel over the net without touching it into the diagonally opposite service box. If the ball hits the net but lands in the service box, this is a let or net service, which is void, and the server retakes that serve. The player can serve any number of let services in a point and they are always treated as voids and not as faults. A fault is a serve that falls long or wide of the service box, or does not clear the net. There is also a "foot fault" when a player's foot touches the baseline or an extension of the center mark before the ball is hit. If the second service, after a fault, is also a fault, the server double faults, and the receiver wins the point. However, if the serve is in, it is considered a legal service.
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A legal service starts a rally, in which the players alternate hitting the ball across the net. A legal return consists of a player hitting the ball so that it falls in the server's court, before it has bounced twice or hit any fixtures except the net. A player or team cannot hit the ball twice in a row. The ball must travel over the net into the other players' court. A ball that hits the net during a rally is considered a legal return as long as it crosses into the opposite side of the court. The first player or team to fail to make a legal return loses the point. The server then moves to the other side of the service line at the start of a new point.[52]
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A game consists of a sequence of points played with the same player serving. A game is won by the first player to have won at least four points in total and at least two points more than the opponent. The running score of each game is described in a manner peculiar to tennis: scores from zero to three points are described as "love", "15", "30", and "40", respectively. If at least three points have been scored by each player, making the player's scores equal at 40 apiece, the score is not called out as "40–40", but rather as "deuce". If at least three points have been scored by each side and a player has one more point than his opponent, the score of the game is "advantage" for the player in the lead. During informal games, "advantage" can also be called "ad in" or "van in" when the serving player is ahead, and "ad out" or "van out" when the receiving player is ahead; alternatively, either player may simply call out "my ad" or "your ad" during informal play.
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The score of a tennis game during play is always read with the serving player's score first. In tournament play, the chair umpire calls the point count (e.g., "15-love") after each point. At the end of a game, the chair umpire also announces the winner of the game and the overall score.
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A set consists of a sequence of games played with service alternating between games, ending when the count of games won meets certain criteria. Typically, a player wins a set by winning at least six games and at least two games more than the opponent. If one player has won six games and the opponent five, an additional game is played. If the leading player wins that game, the player wins the set 7–5. If the trailing player wins the game (tying the set 6–6) a tie-break is played. A tie-break, played under a separate set of rules, allows one player to win one more game and thus the set, to give a final set score of 7–6. A "love" set means that the loser of the set won zero games, colloquially termed a 'jam donut' in the US.[53] In tournament play, the chair umpire announces the winner of the set and the overall score. The final score in sets is always read with the winning player's score first, e.g. "6–2, 4–6, 6–0, 7–5".
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A match consists of a sequence of sets. The outcome is determined through a best of three or five sets system. On the professional circuit, men play best-of-five-set matches at all four Grand Slam tournaments, Davis Cup, and the final of the Olympic Games and best-of-three-set matches at all other tournaments, while women play best-of-three-set matches at all tournaments. The first player to win two sets in a best-of-three, or three sets in a best-of-five, wins the match.[54] Only in the final sets of matches at the French Open, the Olympic Games, and Fed Cup are tie-breaks not played. In these cases, sets are played indefinitely until one player has a two-game lead, occasionally leading to some remarkably long matches.
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In tournament play, the chair umpire announces the end of the match with the well-known phrase "Game, set, match" followed by the winning person's or team's name.
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A game point occurs in tennis whenever the player who is in the lead in the game needs only one more point to win the game. The terminology is extended to sets (set point), matches (match point), and even championships (championship point). For example, if the player who is serving has a score of 40-love, the player has a triple game point (triple set point, etc.) as the player has three consecutive chances to win the game. Game points, set points, and match points are not part of official scoring and are not announced by the chair umpire in tournament play.
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A break point occurs if the receiver, not the server, has a chance to win the game with the next point. Break points are of particular importance because serving is generally considered advantageous, with servers being expected to win games in which they are serving. A receiver who has one (score of 30–40 or advantage), two (score of 15–40) or three (score of love-40) consecutive chances to win the game has break point, double break point or triple break point, respectively. If the receiver does, in fact, win their break point, the game is awarded to the receiver, and the receiver is said to have converted their break point. If the receiver fails to win their break point it is called a failure to convert. Winning break points, and thus the game, is also referred to as breaking serve, as the receiver has disrupted, or broken the natural advantage of the server. If in the following game the previous server also wins a break point it is referred to as breaking back. Except where tie-breaks apply, at least one break of serve is required to win a set (otherwise a two-game lead would never occur).
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Another, however informal, tennis format is called Canadian doubles. This involves three players, with one person playing against a doubles team. The single player gets to utilize the alleys normally reserved only for a doubles team. Conversely, the doubles team does not use the alleys when executing a shot. The scoring is the same as for a regular game. This format is not sanctioned by any official body.
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"Australian doubles", another informal and unsanctioned form of tennis, is played with similar rules to the Canadian doubles style, only in this version, players rotate court position after each game, each player taking a turn at playing alone against the other two. As such, each player plays doubles and singles over the course of a match, with the singles player always serving. Scoring styles vary, but one popular method is to assign a value of 2 points to each game, with the server taking both points if he or she holds serve and the doubles team each taking one if they break serve.
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Wheelchair tennis can be played by able-bodied players as well as people who require a wheelchair for mobility. An extra bounce is permitted. This rule makes it possible to have mixed wheelchair and able-bodied matches. It is possible for a doubles team to consist of a wheelchair player and an able-bodied player (referred to as "one-up, one-down"), or for a wheelchair player to play against an able-bodied player. In such cases, the extra bounce is permitted for the wheelchair users only.
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In most professional play and some amateur competition, there is an officiating head judge or chair umpire (usually referred to simply as the umpire), who sits in a raised chair to one side of the court. The umpire has absolute authority to make factual determinations. The umpire may be assisted by line judges, who determine whether the ball has landed within the required part of the court and who also call foot faults. There also may be a net judge who determines whether the ball has touched the net during service. The umpire has the right to overrule a line judge or a net judge if the umpire is sure that a clear mistake has been made.[57]
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In past tournaments, line judges tasked with calling the serve were sometimes assisted by electronic sensors that beeped to indicate an out-of-bounds serve; one such system was called "Cyclops".[58] Cyclops has since largely been replaced by the Hawk-Eye system.[59][60] In professional tournaments using this system, players are allowed three unsuccessful appeals per set, plus one additional appeal in the tie-break to challenge close line calls by means of an electronic review. The US Open, Miami Masters, US Open Series, and World Team Tennis started using this challenge system in 2006 and the Australian Open and Wimbledon introduced the system in 2007.[61] In clay-court matches, such as at the French Open, a call may be questioned by reference to the mark left by the ball's impact on the court surface.
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The referee, who is usually located off the court, is the final authority about tennis rules. When called to the court by a player or team captain, the referee may overrule the umpire's decision if the tennis rules were violated (question of law) but may not change the umpire's decision on a question of fact. If, however, the referee is on the court during play, the referee may overrule the umpire's decision. (This would only happen in Davis Cup or Fed Cup matches, not at the World Group level, when a chair umpire from a non-neutral country is in the chair).[57]
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Ball boys and girls may be employed to retrieve balls, pass them to the players, and hand players their towels. They have no adjudicative role. In rare events (e.g., if they are hurt or if they have caused a hindrance), the umpire may ask them for a statement of what actually happened. The umpire may consider their statements when making a decision. In some leagues, especially junior leagues, players make their own calls, trusting each other to be honest. This is the case for many school and university level matches. The referee or referee's assistant, however, can be called on court at a player's request, and the referee or assistant may change a player's call. In unofficiated matches, a ball is out only if the player entitled to make the call is sure that the ball is out.
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In tennis, a junior is a player under 18 who is still legally protected by a parent or guardian. Players on the main adult tour who are under 18 must have documents signed by a parent or guardian. These players, however, are still eligible to play in junior tournaments.
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The International Tennis Federation (ITF) conducts a junior tour that allows juniors to establish a world ranking and an Association of Tennis Professionals (ATP) or Women's Tennis Association (WTA) ranking. Most juniors who enter the international circuit do so by progressing through ITF, Satellite, Future, and Challenger tournaments before entering the main circuit. The latter three circuits also have adults competing in them. Some juniors, however, such as Australian Lleyton Hewitt and Frenchman Gaël Monfils, have catapulted directly from the junior tour to the ATP tour by dominating the junior scene or by taking advantage of opportunities given to them to participate in professional tournaments.
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In 2004, the ITF implemented a new rankings scheme to encourage greater participation in doubles, by combining two rankings (singles and doubles) into one combined tally.[62] Junior tournaments do not offer prize money except for the Grand Slam tournaments, which are the most prestigious junior events. Juniors may earn income from tennis by participating in the Future, Satellite, or Challenger tours. Tournaments are broken up into different tiers offering different amounts of ranking points, culminating with Grade A.
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Leading juniors are allowed to participate for their nation in the Junior Fed Cup and Davis Cup competitions. To succeed in tennis often means having to begin playing at a young age. To facilitate and nurture a junior's growth in tennis, almost all tennis playing nations have developed a junior development system. Juniors develop their play through a range of tournaments on all surfaces, accommodating all different standards of play. Talented juniors may also receive sponsorships from governing bodies or private institutions.
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A tennis match is intended to be continuous.[63] Because stamina is a relevant factor, arbitrary delays are not permitted. In most cases, service is required to occur no more than 20 seconds after the end of the previous point.[63] This is increased to 90 seconds when the players change ends (after every odd-numbered game), and a 2-minute break is permitted between sets.[63] Other than this, breaks are permitted only when forced by events beyond the players' control, such as rain, damaged footwear, damaged racket, or the need to retrieve an errant ball. Should a player be deemed to be stalling repeatedly, the chair umpire may initially give a warning followed by subsequent penalties of "point", "game", and default of the match for the player who is consistently taking longer than the allowed time limit.[64]
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In the event of a rain delay, darkness or other external conditions halting play, the match is resumed at a later time, with the same score as at the time of the delay, and each player at the same end of the court as when rain halted play, or as close to the same relative compass point if play is resumed on a different court.
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Balls wear out quickly in serious play and, therefore, in ATP and WTA tournaments, they are changed after every nine games with the first change occurring after only seven games, because the first set of balls is also used for the pre-match warm-up.[43] In ITF tournaments like Fed Cup, the balls are changed after every eleven games (rather than nine) with the first change occurring after only nine games (instead of seven). An exception is that a ball change may not take place at the beginning of a tiebreaker, in which case the ball change is delayed until the beginning of the second game of the next set.[47] As a courtesy to the receiver, the server will often signal to the receiver before the first serve of the game in which new balls are used as a reminder that they are using new balls. Continuity of the balls' condition is considered part of the game, so if a re-warm-up is required after an extended break in play (usually due to rain), then the re-warm-up is done using a separate set of balls, and use of the match balls is resumed only when play resumes.
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A recent rule change is to allow coaching on court on a limited basis during a match.[65][66][67][68] This has been introduced in women's tennis for WTA Tour events in 2009 and allows the player to request her coach once per set.[69]
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Stance refers to the way a player prepares themselves in order to best be able to return a shot. Essentially, it enables them to move quickly in order to achieve a particular stroke. There are four main stances in modern tennis: open, semi-open, closed, and neutral. All four stances involve the player crouching in some manner: as well as being a more efficient striking posture, it allows them to isometrically preload their muscles in order to play the stroke more dynamically. What stance is selected is strongly influenced by shot selection. A player may quickly alter their stance depending on the circumstances and the type of shot they intend to play. Any given stance also alters dramatically based upon the actual playing of the shot with dynamic movements and shifts of body weight occurring.[70][71]
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This is the most common stance in tennis. The player’s feet are placed parallel to the net. They may be pointing sideways, directly at the net or diagonally towards it. This stance allows for a high degree of torso rotation which can add significant power to the stroke. This process is sometimes likened to the coiling and uncoiling of a spring. i.e the torso is rotated as a means of preloading the muscular system in preparation for playing the stroke: this is the coiling phase. When the stroke is played the torso rotates to face forwards again, called uncoiling, and adds significant power to the stroke. A disadvantage of this stance is that it does not always allow ‘for proper weight transfer and maintenance of balance’[70] when making powerful strokes. It is commonly used for forehand strokes; double-handed backhands can also be made effectively from it.
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This stance is somewhere between open and closed and is a very flexible stance. The feet are aligned diagonally towards the net. It allows for a lot of shoulder rotation and the torso can be coiled, before being uncoiled into the shot in order to increase the power of the shot. It is commonly used in modern tennis especially by ‘top professional players on the forehand’.[72] Two-handed backhands can also be employed from this stance.
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The closed stance is the least commonly used of the three main stances. One foot is placed further towards the net with the other foot further from it; there is a diagonal alignment between the feet. It allows for effective torso rotation in order to increase the power of the shot. It is usually used to play backhand shots and it is rare to see forehand shots played from it. A stroke from this stance may entail the rear foot coming completely off the floor with bodyweight being transferred entirely to the front foot.[70]
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[71]
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This is sometimes also referred to as the square stance. One foot is positioned closer to the net and ahead of the other which is behind and in line with it. Both feet are aligned at a 90 degree angle to the net. The neutral stance is often taught early because ‘It allows beginners to learn about shifting weight and rotation of the body.’[71] Forehands and backhands may be made from it.[73]
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A competent tennis player has eight basic shots in his or her repertoire: the serve, forehand, backhand, volley, half-volley, overhead smash, drop shot, and lob.
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A grip is a way of holding the racket in order to hit shots during a match. The grip affects the angle of the racket face when it hits the ball and influences the pace, spin, and placement of the shot. Players use various grips during play, including the Continental (The "Handshake Grip"), Eastern (Can be either semi-eastern or full eastern. Usually used for backhands.), and Western (semi-western or full western, usually for forehand grips) grips. Most players change grips during a match depending on what shot they are hitting; for example, slice shots and serves call for a Continental grip.[74]
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A serve (or, more formally, a "service") in tennis is a shot to start a point. The serve is initiated by tossing the ball into the air and hitting it (usually near the apex of its trajectory) into the diagonally opposite service box without touching the net. The serve may be hit under- or overhand although underhand serving remains a rarity.[75] If the ball hits the net on the first serve and bounces over into the correct diagonal box then it is called a "let" and the server gets two more additional serves to get it in. There can also be a let if the server serves the ball and the receiver isn't prepared.[47] If the server misses his or her first serve and gets a let on the second serve, then they get one more try to get the serve in the box.
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Experienced players strive to master the conventional overhand serve to maximize its power and placement. The server may employ different types of serve including flat serve, topspin serve, slice serve, and kick (American twist) serve. A reverse type of spin serve is hit in a manner that spins the ball opposite the natural spin of the server, the spin direction depending upon right- or left-handedness. If the ball is spinning counterclockwise, it will curve right from the hitter's point of view and curve left if spinning clockwise.[76]
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Some servers are content to use the serve simply to initiate the point; however, advanced players often try to hit a winning shot with their serve. A winning serve that is not touched by the opponent is called an "ace".
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For a right-handed player, the forehand is a stroke that begins on the right side of the body, continues across the body as contact is made with the ball, and ends on the left side of the body. There are various grips for executing the forehand, and their popularity has fluctuated over the years. The most important ones are the continental, the eastern, the semi-western, and the western. For a number of years, the small, frail 1920s player Bill Johnston was considered by many to have had the best forehand of all time, a stroke that he hit shoulder-high using a western grip. Few top players used the western grip after the 1920s, but in the latter part of the 20th century, as shot-making techniques and equipment changed radically, the western forehand made a strong comeback and is now used by many modern players. No matter which grip is used, most forehands are generally executed with one hand holding the racket, but there have been fine players with two-handed forehands. In the 1940s and 50s, the Ecuadorian/American player Pancho Segura used a two-handed forehand to achieve a devastating effect against larger, more powerful players. Players such as Monica Seles or France's Fabrice Santoro and Marion Bartoli are also notable players known for their two-handed forehands.[77]
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For right-handed players, the backhand is a stroke that begins on the left side of their body, continues across their body as contact is made with the ball, and ends on the right side of their body. It can be executed with either one hand or with both and is generally considered more difficult to master than the forehand. For most of the 20th century, the backhand was performed with one hand, using either an eastern or a continental grip. The first notable players to use two hands were the 1930s Australians Vivian McGrath and John Bromwich, but they were lonely exceptions. The two-handed grip gained popularity in the 1970s as Björn Borg, Chris Evert, Jimmy Connors, and later Mats Wilander and Marat Safin used it to great effect, and it is now used by a large number of the world's best players, including Rafael Nadal and Serena Williams.[78]
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Two hands give the player more control, while one hand can generate a slice shot, applying backspin on the ball to produce a low trajectory bounce. Reach is also limited with the two-handed shot. The player long considered to have had the best backhand of all time, Don Budge, had a powerful one-handed stroke in the 1930s and 1940s that imparted topspin onto the ball. Ken Rosewall, another player noted for his one-handed backhand, used a very accurate slice backhand through the 1950s and 1960s. A small number of players, notably Monica Seles, use two hands on both the backhand and forehand sides.
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A volley is a shot returned to the opponent in mid-air before the ball bounces, generally performed near the net, and is usually made with a stiff-wristed punching motion to hit the ball into an open area of the opponent's court. The half volley is made by hitting the ball on the rise just after it has bounced, also generally in the vicinity of the net, and played with the racket close to the ground.[79] The swinging volley is hit out of the air as the player approaches the net. It is an offensive shot used to take preparation time away from the opponent, as it returns the ball into the opponent's court much faster than a standard volley.
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From a poor defensive position on the baseline, the lob can be used as either an offensive or defensive weapon, hitting the ball high and deep into the opponent's court to either enable the lobber to get into better defensive position or to win the point outright by hitting it over the opponent's head. If the lob is not hit deeply enough into the other court, however, an opponent near the net may then hit an overhead smash, a hard, serve-like shot, to try to end the point.
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A difficult shot in tennis is the return of an attempted lob over the backhand side of a player. When the contact point is higher than the reach of a two-handed backhand, most players will try to execute a high slice (under the ball or sideways). Fewer players attempt the backhand sky-hook or smash. Rarely, a player will go for a high topspin backhand, while themselves in the air. A successful execution of any of these alternatives requires balance and timing, with less margin of error than the lower contact point backhands, since this shot is a break in the regular pattern of play.
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If an opponent is deep in his court, a player may suddenly employ an unexpected drop shot, by softly tapping the ball just over the net so that the opponent is unable to run in fast enough to retrieve it. Advanced players will often apply back spin to a drop shot, causing the ball to "skid" upon landing and bounce sideways, with less forward momentum toward their opponent, or even backwards towards the net, thus making it even more difficult to return.
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Muscle strain is one of the most common injuries in tennis.[80] When an isolated large-energy appears during the muscle contraction and at the same time body weight apply huge amount of pressure to the lengthened muscle, muscle strain can occur.[81] Inflammation and bleeding are triggered when muscle strain occurs, which can result in redness, pain and swelling.[81] Overuse is also common in tennis players of all levels. Muscle, cartilage, nerves, bursae, ligaments and tendons may be damaged from overuse. The repetitive use of a particular muscle without time for repair and recovery is the most common cause of injury.[81]
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Tournaments are often organized by gender and number of players. Common tournament configurations include men's singles, women's singles, and doubles, where two players play on each side of the net. Tournaments may be organized for specific age groups, with upper age limits for youth and lower age limits for senior players. Example of this include the Orange Bowl and Les Petits As junior tournaments. There are also tournaments for players with disabilities, such as wheelchair tennis and deaf tennis.[82] In the four Grand Slam tournaments, the singles draws are limited to 128 players for each gender.
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Most large tournaments seed players, but players may also be matched by their skill level. According to how well a person does in sanctioned play, a player is given a rating that is adjusted periodically to maintain competitive matches. For example, the United States Tennis Association administers the National Tennis Rating Program (NTRP), which rates players between 1.0 and 7.0 in 1/2 point increments. Average club players under this system would rate 3.0–4.5 while world class players would be 7.0 on this scale.
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The four Grand Slam tournaments are considered to be the most prestigious tennis events in the world. They are held annually and comprise, in chronological order, the Australian Open, the French Open, Wimbledon, and the US Open. Apart from the Olympic Games, Davis Cup, Fed Cup, and Hopman Cup, they are the only tournaments regulated by the International Tennis Federation (ITF).[83] The ITF's national associations, Tennis Australia (Australian Open), the Fédération Française de Tennis (French Open), the Lawn Tennis Association (Wimbledon) and the United States Tennis Association (US Open) are delegated the responsibility to organize these events.[83]
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Aside from the historical significance of these events, they also carry larger prize funds than any other tour event and are worth double the number of ranking points to the champion than in the next echelon of tournaments, the Masters 1000 (men) and Premier events (women).[84][85] Another distinguishing feature is the number of players in the singles draw. There are 128, more than any other professional tennis tournament. This draw is composed of 32 seeded players, other players ranked in the world's top 100, qualifiers, and players who receive invitations through wild cards. Grand Slam men's tournaments have best-of-five set matches while the women play best-of-three. Grand Slam tournaments are among the small number of events that last two weeks, the others being the Indian Wells Masters and the Miami Masters.
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Currently, the Grand Slam tournaments are the only tour events that have mixed doubles contests. Grand Slam tournaments are held in conjunction with wheelchair tennis tournaments and junior tennis competitions. These tournaments also contain their own idiosyncrasies. For example, players at Wimbledon are required to wear predominantly white. Andre Agassi chose to skip Wimbledon from 1988 through 1990 citing the event's traditionalism, particularly its "predominantly white" dress code.[86] Wimbledon has its own particular methods for disseminating tickets, often leading tennis fans to follow complex procedures to obtain tickets.[87]
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* The international tournament began in 1925
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The ATP World Tour Masters 1000 is a group of nine tournaments that form the second-highest echelon in men's tennis. Each event is held annually, and a win at one of these events is worth 1000 ranking points. When the ATP, led by Hamilton Jordan, began running the men's tour in 1990, the directors designated the top nine tournaments, outside of the Grand Slam events, as "Super 9" events.[88] In 2000 this became the Tennis Masters Series and in 2004 the ATP Masters Series. In November at the end of the tennis year, the world's top eight players compete in the ATP World Tour Finals, a tournament with a rotating locale. It is currently held in London, England.[89]
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In August 2007 the ATP announced major changes to the tour that were introduced in 2009. The Masters Series was renamed to the "Masters 1000", the addition of the number 1000 referring to the number of ranking points earned by the winner of each tournament. Contrary to earlier plans, the number of tournaments was not reduced from nine to eight and the Monte Carlo Masters remains part of the series although, unlike the other events, it does not have a mandatory player commitment. The Hamburg Masters has been downgraded to a 500-point event. The Madrid Masters moved to May and onto clay courts, and a new tournament in Shanghai took over Madrid's former indoor October slot. As of 2011 six of the nine "1000" level tournaments are combined ATP and WTA events.[90]
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The third and fourth tier of men's tennis tournaments are formed by the ATP World Tour 500 series, consisting of 11 tournaments, and the ATP World Tour 250 series with 40 tournaments.[91] Like the ATP World Tour Masters 1000, these events offer various amounts of prize money and the numbers refer to the amount of ranking points earned by the winner of a tournament.[84] The Dubai Tennis Championships offer the largest financial incentive to players, with total prize money of US$2,313,975 (2012).[92] These series have various draws of 28, 32, 48 and 56 for singles and 16 and 24 for doubles. It is mandatory for leading players to enter at least four 500 events, including at least one after the US Open.
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The Challenger Tour for men is the lowest level of tournament administered by the ATP. It is composed of about 150 events and, as a result, features a more diverse range of countries hosting events.[93] The majority of players use the Challenger Series at the beginning of their career to work their way up the rankings. Andre Agassi, between winning Grand Slam tournaments, plummeted to World No. 141 and used Challenger Series events for match experience and to progress back up the rankings.[94] The Challenger Series offers prize funds of between US$25,000 and US$150,000.
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Below the Challenger Tour are the Futures tournaments, events on the ITF Men's Circuit. These tournaments also contribute towards a player's ATP rankings points. Futures Tournaments offer prize funds of between US$10,000 and US$15,000.[95] Approximately 530 Futures Tournaments are played each year.
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Premier events for women form the most prestigious level of events on the Women's Tennis Association Tour after the Grand Slam tournaments. These events offer the largest rewards in terms of points and prize money. Within the Premier category are Premier Mandatory, Premier 5, and Premier tournaments. The Premier events were introduced in 2009 replacing the previous Tier I and II tournament categories. Currently four tournaments are Premier Mandatory, five tournaments are Premier 5, and twelve tournaments are Premier. The first tiering system in women's tennis was introduced in 1988. At the time of its creation, only two tournaments, the Lipton International Players Championships in Florida and the German Open in Berlin, comprised the Tier I category.
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International tournaments are the second main tier of the WTA tour and consist of 31 tournaments, with a prize money for every event at U.S.$220,000, except for the year-ending Commonwealth Bank Tournament of Champions in Bali, which has prize money of U.S.$600,000.
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Professional tennis players enjoy the same relative perks as most top sports personalities: clothing, equipment and endorsements. Like players of other individual sports such as golf, they are not salaried, but must play and finish highly in tournaments to obtain prize money.
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In recent years, some controversy has surrounded the involuntary or deliberate noise caused by players' grunting.[citation needed]
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While players are gradually less competitive in singles by their late 20s and early 30s, they can still continue competitively in doubles (as instanced by Martina Navratilova and John McEnroe, who won doubles titles in their 40s).
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In the Open Era, several female players such as Martina Navratilova, Margaret Court, Martina Hingis, Serena Williams, and Venus Williams (the latter two sisters playing together) have been prolific at both singles and doubles events throughout their careers. John McEnroe is one of the very few professional male players to be top ranked in both singles and doubles at the same time,[96][97][98] and Yevgeny Kafelnikov is the most recent male player to win multiple Grand Slams in both singles and doubles during the same period of his career.
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In terms of public attention and earnings (see below), singles champions have far surpassed their doubles counterparts. The Open Era, particularly the men's side, has seen many top-ranked singles players that only sparingly compete in doubles, while having "doubles specialists" who are typically being eliminated early in the singles draw but do well in the doubles portion of a tournament. Notable doubles pairings include The Woodies (Todd Woodbridge and Mark Woodforde) and the Bryan Brothers (identical twin brothers Robert Charles "Bob" Bryan and Michael Carl "Mike" Bryan). Woodbridge has disliked the term "doubles ‘specialists’", saying that he and Woodforde "set a singles schedule and doubles fitted in around that", although later in Woodbridge's career he focused exclusively on doubles as his singles ranking fell too low that it was no longer financially viable to recover at that age. Woodbridge noted that while top singles players earn enough that they don't need to nor want to play doubles, he suggested that lower-ranked singles players outside the Top Ten should play doubles to earn more playing time and money.[99][100]
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The Olympics doubles tennis tournament necessitates that both members of a doubles pairing be from the same country, hence several top professional pairs such as Jamie Murray and Bruno Soares cannot compete in the Olympics. Top-ranked singles players that are usually rivals on the professional circuit, such as Boris Becker and Michael Stich, and Roger Federer and Stan Wawrinka have formed a rare doubles partnership for the Olympics. Unlike professional tennis tournaments (see below) where singles players receive much more prize money than doubles players, an Olympic medal for both singles and doubles has similar prestige. The Olympics is more of a priority for doubles champions while singles champions often skip the tournament.[99][100] While the ATP has voted for Olympic results to count towards player ranking points, WTA players voted against it.[101]
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For the 2000 Olympics, Lisa Raymond was passed over for Team USA in favor of Serena Williams by captain Billie Jean King, even though Raymond was the top-ranked doubles player in the world at the time, and Raymond unsuccessfully challenged the selection.[101]
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In professional tennis tournaments such as Wimbledon, the singles competition receives the most prize money and coverage, followed by doubles, and then mixed doubles usually receive the lowest monetary awards.[102] For instance in the US Open as of 2018, the men's and women's singles prize money (US$40,912,000) accounts for 80.9 percent of total player base compensation, while men's and women's doubles (US$6,140,840), men's and women's singles qualifying (US$3,008,000), and mixed doubles (US$505,000) account for 12.1 percent, 5.9 percent, and 1.0 percent, respectively. The singles winner receives US$3,800,000, while the doubles winning pair receives $700,000 and the mixed doubles winning pair receives US$155,000.[103]
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The following players have won at least five singles titles at Grand Slam tournaments:
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A frequent topic of discussion among tennis fans and commentators is who was the greatest male singles player of all time. By a large margin, an Associated Press poll in 1950 named Bill Tilden as the greatest player of the first half of the 20th century.[104] From 1920 to 1930, Tilden won singles titles at Wimbledon three times and the U.S. Championships seven times. In 1938, however, Donald Budge became the first person to win all four major singles titles during the same calendar year, the Grand Slam, and won six consecutive major titles in 1937 and 1938. Tilden called Budge "the finest player 365 days a year that ever lived."[105] In his 1979 autobiography, Jack Kramer said that, based on consistent play, Budge was the greatest player ever.[106] Some observers, however, also felt that Kramer deserved consideration for the title. Kramer was among the few who dominated amateur and professional tennis during the late 1940s and early 1950s. Tony Trabert has said that of the players he saw before the start of the Open Era, Kramer was the best male champion.[107]
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By the 1960s, Budge and others had added Pancho Gonzales and Lew Hoad to the list of contenders. Budge reportedly believed that Gonzales was the greatest player ever.[108] Gonzales said about Hoad, "When Lew's game was at its peak nobody could touch him. ... I think his game was the best game ever. Better than mine. He was capable of making more shots than anybody. His two volleys were great. His overhead was enormous. He had the most natural tennis mind with the most natural tennis physique."[109]
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Before and during the Open Era, Rod Laver remains the only male player in history to have won the calendar year Grand Slam twice in 1962 and 1969 [110] and also the calendar year Professional Grand Slam in 1967.[111] More recently Björn Borg and Pete Sampras were regarded by many of their contemporaries as among the greatest ever. Andre Agassi, the first of two male players in history to have achieved a Career Golden Slam in singles tennis (followed by Rafael Nadal), has been called the best service returner in the history of the game.[112][113][114][115] He is the first man to win grand slams on all modern surfaces (previous holders of all grand slam tournaments played in an era of grass and clay only), and is regarded by a number of critics and fellow players to be among the greatest players of all time.[112][116][117] Both Rod Laver and Ken Rosewall also won major Pro Slam tournaments on all three surfaces (grass, clay, hard court) Rosewall in 1963 and Laver in 1967.[118]
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By the early twenty-first century, Roger Federer is considered by many observers to have the most "complete" game in modern tennis. He has won 20 grand slam titles and 6 World Tour Finals, the most for any male player. Many experts of tennis, former tennis players and his own tennis peers believe Federer is the greatest player in the history of the game.[119][120][121][122][123][124] Federer's biggest rival Rafael Nadal is regarded as the greatest competitor in tennis history by some former players and is regarded to have the potential to be the greatest of all time.[125][126] Nadal is regarded as the greatest clay court player of all time.[127] Novak Djokovic, a rival of both Nadal and Federer, is also considered to be one of the greatest tennis players of all time.[128]
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As with the men there are frequent discussions about who is the greatest female singles player of all time with Steffi Graf, Martina Navratilova and Serena Williams being the three players most often nominated.
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In March 2012 the TennisChannel published a combined list of the 100 greatest men and women tennis players of all time.[129] It ranked Steffi Graf as the greatest female player (in 3rd place overall), followed by Martina Navratilova (4th place) and Margaret Court (8th place). The rankings were determined by an international panel.
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Sportswriter John Wertheim of Sports Illustrated stated in an article in July 2010 that Serena Williams is the greatest female tennis player ever with the argument that "Head-to-head, on a neutral surface (i.e. hard courts), everyone at their best, I can't help feeling that she crushes the other legends.".[130] In a reaction to this article Yahoo sports blog Busted Racket published a list of the top-10 women's tennis players of all time placing Martina Navratilova in first spot.[131] This top-10 list was similar to the one published in June 2008 by the Bleacher Report who also ranked Martina Navratilova as the top female player of all time.[132]
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Steffi Graf is considered by some to be the greatest female player. Billie Jean King said in 1999, "Steffi is definitely the greatest women's tennis player of all time."[133] Martina Navratilova has included Graf on her list of great players.[133] In December 1999, Graf was named the greatest female tennis player of the 20th century by a panel of experts assembled by the Associated Press.[134] Tennis writer Steve Flink, in his book The Greatest Tennis Matches of the Twentieth Century, named her as the best female player of the 20th century, directly followed by Martina Navratilova.[135]
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Tennis magazine selected Martina Navratilova as the greatest female tennis player for the years 1965 through 2005.[136][137] Tennis historian and journalist Bud Collins has called Navratilova "arguably, the greatest player of all time."[138] Billie Jean King said about Navratilova in 2006, "She's the greatest singles, doubles and mixed doubles player who's ever lived."[139]
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A tenor is a type of classical male singing voice whose vocal range lies between the countertenor and baritone voice types. It is one of the highest of the male voice types.[1] The tenor's vocal range extends up to C5. The low extreme for tenors is roughly A♭2 (two A♭s below middle C). At the highest extreme, some tenors can sing up to the second F above middle C (F5).[2][page needed] The tenor voice type is generally divided into the leggero tenor, lyric tenor, spinto tenor, dramatic tenor, heldentenor, and tenor buffo or spieltenor.
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The name "tenor" derives from the Latin word tenere, which means "to hold". As Fallows, Jander, Forbes, Steane, Harris and Waldman note in the "Tenor" article at Grove Music Online:
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In polyphony between about 1250 and 1500, the [tenor was the] structurally fundamental (or 'holding') voice, vocal or instrumental; by the 15th century it came to signify the male voice that sang such parts.[3]
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All other voices were normally calculated in relation to the tenor, which often proceeded in longer note values and carried a borrowed Cantus firmus melody.[citation needed] Until the late 16th century introduction of the contratenor singers, the tenor was usually the highest voice, assuming the role of providing a foundation.[citation needed] It was also in the 18th century that "tenor" came to signify the male voice that sang such parts. Thus, for earlier repertoire, a line marked 'tenor' indicated the part's role, and not the required voice type; indeed, even as late as the eighteenth century, partbooks labelled 'tenor' might contain parts for a range of voice types.[4][page needed]
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The vocal range of the tenor is one of the highest of the male voice types. Within opera, the lowest note in the standard tenor repertoire is probably[weasel words] A♭2 in Rossini's rarely performed La donna del lago in the role of Rodrigo di Dhu, written for Andrea Nozzari. Within more frequently performed repertoire, Mime and Herod both call for an A2. A few tenor roles in the standard repertoire call for a "tenor C" (C5, one octave above middle C). Some, if not all, of the few top Cs in the standard operatic repertoire are either optional—such as in "Che gelida manina" in Puccini's La bohème—or interpolated (added) by tradition, such as in "Di quella pira" from Verdi's Il trovatore);[citation needed] however, the highest demanded note in the standard tenor operatic repertoire is D5, found in "Mes amis, écoutez l'histoire"[5] from Adolphe Adam's Le postillon de Lonjumeau and "Loin de son amie"[6] from Fromental Halévy's La Juive). In the leggero repertoire, the highest note is F5 (Arturo in "Credeasi, misera" from Bellini's I puritani),[7][original research?] therefore, very few tenors have this role in their repertoire without transposition (given the raising of concert pitch since its composition),[8][page needed] or resorting to falsetto.
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In SATB four-part mixed chorus, the tenor is the second lowest vocal range, above the bass and below the alto and soprano. Men's chorus usually denotes an ensemble of TTBB in which the first tenor is the highest voice. Whilst certain choral music does require the first tenors to ascend the full tenor range, the majority of choral music places the tenors in the range from approximately B2 up to A4. The requirements of the tenor voice in choral music are also tied to the style of music most often performed by a given choir. Orchestra choruses require tenors with fully resonant voices, but chamber or a cappella choral music (sung with no instrumental accompaniment) can sometimes rely on light baritones singing in falsetto.[9][page needed]
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Even so, one nearly ubiquitous facet of choral singing is the shortage of tenor voices.[10][better source needed][11][page needed] Most men tend to have baritone voices and for this reason the majority of men tend to prefer singing in the bass section of a choir. (However true basses are even rarer than tenors.) Some men sing tenor even if they lack the full range, and sometimes low altos sing the tenor part.[9][page needed] In men's choruses that consist of four male vocal parts TTBB (tenor 1, tenor 2, bass 1, bass 2), tenors will often sing both in chest tone and falsetto, extending the vocal range of the choir.
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Within the tenor voice type category are seven generally recognized subcategories: leggero tenor, lyric tenor, spinto tenor, dramatic tenor, heldentenor, Mozart tenor, and tenor buffo or spieltenor. There is considerable overlap between the various categories of role and of voice-type; some tenor singers have begun with lyric voices but have transformed with time into spinto or even dramatic tenors.
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Also known as the tenore di grazia, the leggero tenor is essentially the male equivalent of a lyric coloratura. This voice is light, agile, and capable of executing difficult passages of fioritura. The typical leggero tenor possesses a range spanning from approximately C3 to E5, with a few being able to sing up to F5 or higher in full voice. In some cases, the chest register of the leggero tenor may extend below C3. Voices of this type are utilized frequently in the operas of Rossini, Donizetti, Bellini and in music dating from the Baroque period.[citation needed]
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Leggero tenor roles in operas:[8][page needed]
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The lyric tenor is a warm graceful voice with a bright, full timbre that is strong but not heavy and can be heard over an orchestra. Lyric tenors have a range from approximately the C one octave below middle C (C3) to the D one octave above middle C (D5). Similarly, their lower range may extend a few notes below the C3. There are many vocal shades to the lyric tenor group, repertoire should be selected according to the weight, colors, and abilities of the voice.
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Lyric tenor roles in operas:[8][page needed]
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The spinto tenor has the brightness and height of a lyric tenor, but with a heavier vocal weight enabling the voice to be "pushed" to dramatic climaxes with less strain than the lighter-voice counterparts. Spinto tenors have a darker timbre than a lyric tenor, without having a vocal color as dark as many (not all) dramatic tenors. The German equivalent of the Spinto fach is the Jugendlicher Heldentenor and encompasses many of the Dramatic tenor roles as well as some Wagner roles such as Lohengrin and Stolzing. The difference is often the depth and metal in the voice where some lyric tenors age or push their way into singing as a Spinto giving them a lighter tone and a Jugendlicher Heldentenor tends to be either a young heldentenor or true lyric spinto. Spinto tenors have a range from approximately the C one octave below middle C (C3) to the C one octave above middle C (C5).
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Spinto tenor roles in operas:[8][page needed]
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Also "tenore di forza" or "robusto", the dramatic tenor has an emotive, ringing and very powerful, clarion, heroic tenor sound. The dramatic tenor's approximate range is from the B one octave below middle C (B2) to the B one octave above middle C (B4) with some able to sing up to the C one octave above middle C (C5).[8][page needed] Many successful dramatic tenors though have historically avoided the coveted high C in performance. Their lower range tends to extend into the baritone tessitura or, a few notes below the C3, even down to A♭2. Some dramatic tenors have a rich and dark tonal colour to their voice (such as the mature Enrico Caruso) while others (like Francesco Tamagno) possess a bright, steely timbre.
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Dramatic tenor roles in operas:[8][page needed]
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The heldentenor (English: heroic tenor) has a rich, dark, powerful and dramatic voice. As its name implies, the heldentenor vocal fach features in the German romantic operatic repertoire. The heldentenor is the German equivalent of the tenore drammatico, however with a more baritonal quality: the typical Wagnerian protagonist. The keystone of the heldentenor's repertoire is arguably Wagner's Siegfried, an extremely demanding role requiring a wide vocal range and great power, plus tremendous stamina and acting ability. Often the heldentenor is a baritone who has transitioned to this fach or tenors who have been misidentified as baritones. Therefore, the heldentenor voice might or might not have facility up to high B or C. The repertoire, however, rarely calls for such high notes.
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Heldentenor roles in operas:[8][page needed]
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A Mozart tenor is yet another distinct tenor type. In Mozart singing, the most important element is the instrumental approach of the vocal sound which implies: flawless and slender emission of sound, perfect intonation, legato, diction and phrasing, capability to cope with the dynamic requirements of the score, beauty of timbre, secure line of singing through perfect support and absolute breath control, musical intelligence, body discipline, elegance, nobility, agility and, most importantly, ability for dramatic expressiveness within the narrow borders imposed by the strict Mozartian style.
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The German Mozart tenor tradition goes back to the end of the 1920s, when Mozart tenors started making use of Caruso's technique (a tenor who rarely sang Mozart) to achieve and improve the required dynamics and dramatic expressiveness.
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Mozart tenor roles in operas:[8][page needed]
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A Tenor buffo or spieltenor is a tenor with good acting ability, and the ability to create distinct voices for his characters. This voice specializes in smaller comic roles. The range of the tenor buffo is from the C one octave below middle C (C3) to the C one octave above middle C (C5).[13] The tessitura of these parts ranges from lower than other tenor roles to very high and broad. These parts are often played by younger tenors who have not yet reached their full vocal potential or older tenors who are beyond their prime singing years. Only rarely will a singer specialize in these roles for an entire career.[8][page needed] In French opéra comique, supporting roles requiring a thin voice but good acting are sometimes described as 'trial', after the singer Antoine Trial (1737–1795), examples being in the operas of Ravel and in The Tales of Hoffmann.[14][page needed]
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Tenor buffo or spieltenor roles in operas:[8][page needed]
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All of Gilbert and Sullivan's Savoy operas have at least one lead lyric tenor character. Notable operetta roles are:
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There are four parts in Barbershop harmony: bass, baritone, lead, and tenor (lowest to highest), with "tenor" referring to the highest part. The tenor generally sings in falsetto voice, corresponding roughly to the countertenor in classical music, and harmonizes above the lead, who sings the melody. The barbershop tenor range is Middle C to A one octave above Middle C, though it is written an octave lower. The "lead" in barbershop music is equivalent to the normal tenor range.[15][page needed]
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In bluegrass music, the melody line is called the lead. Tenor is sung an interval of a third above the lead. Baritone is the fifth of the scale that has the lead as a tonic, and may be sung below the lead, or even above the lead (and the tenor), in which case it is called "high baritone."[16][page needed]
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1 |
+
Voltage, electric potential difference, electric pressure or electric tension is the difference in electric potential between two points, which (in a static electric field) is defined as the work needed per unit of charge to move a test charge between the two points. In the International System of Units, the derived unit for voltage (potential difference) is named volt.[1]:166 In SI units, work per unit charge is expressed as joules per coulomb, where 1 volt = 1 joule (of work) per 1 coulomb (of charge). The old SI definition for volt used power and current, since 1990 quantum Hall and Josephson effect were used and recently (2019) fundamental physical constants were introduced for the definition of all SI units and derived units.[1]:177f, 197f Voltage or electric potential difference is denoted symbolically by ∆V, simplified V,[2] or U,[3] for instance in the context of Ohm's or Kirchhoff's circuit laws.
|
2 |
+
|
3 |
+
Electric potential differences between points can be caused by electric charge, by electric current through a magnetic field, by time-varying magnetic fields, or some combination of these three.[4][5] A voltmeter can be used to measure the voltage (or potential difference) between two points in a system; often a common reference potential such as the ground of the system is used as one of the points. A voltage may represent either a source of energy (electromotive force) or lost, used, or stored energy (potential drop).
|
4 |
+
|
5 |
+
There are multiple useful ways to define voltage, including the standard definition mentioned at the start of this page. There are also other useful definitions of work per charge (see this section).
|
6 |
+
|
7 |
+
Voltage is defined so that negatively charged objects are pulled towards higher voltages, while positively charged objects are pulled towards lower voltages. Therefore, the conventional current in a wire or resistor always flows from higher voltage to lower voltage.
|
8 |
+
|
9 |
+
Historically, voltage has been referred to using terms like "tension" and "pressure". Even today, the term "tension" is still used, for example within the phrase "high tension" (HT) which is commonly used in thermionic valve (vacuum tube) based electronics.
|
10 |
+
|
11 |
+
The voltage increase from some point
|
12 |
+
|
13 |
+
|
14 |
+
|
15 |
+
|
16 |
+
x
|
17 |
+
|
18 |
+
A
|
19 |
+
|
20 |
+
|
21 |
+
|
22 |
+
|
23 |
+
{\textstyle x_{A}}
|
24 |
+
|
25 |
+
to some point
|
26 |
+
|
27 |
+
|
28 |
+
|
29 |
+
|
30 |
+
x
|
31 |
+
|
32 |
+
B
|
33 |
+
|
34 |
+
|
35 |
+
|
36 |
+
|
37 |
+
{\textstyle x_{B}}
|
38 |
+
|
39 |
+
is given by
|
40 |
+
|
41 |
+
In this case, the voltage increase from point A to point B is equal to the work done per unit charge, against the electric field, to move the charge from A to B without causing any acceleration. Mathematically, this is expressed as the line integral of the electric field along that path.
|
42 |
+
Under this definition, the voltage difference between two points is not uniquely defined when there are time-varying magnetic fields since the electric force is not a conservative force in such cases.
|
43 |
+
|
44 |
+
If this definition of voltage is used, any circuit where there are time-varying magnetic fields,[note 1] such as circuits containing inductors, will not have a well-defined voltage between nodes in the circuit. However, if magnetic fields are suitably contained to each component, then the electric field is conservative in the region exterior[note 2] to the components, and voltages are well-defined in that region.[6] In this case, the voltage across an inductor, viewed externally, turns out to be
|
45 |
+
|
46 |
+
U
|
47 |
+
=
|
48 |
+
Δ
|
49 |
+
V
|
50 |
+
=
|
51 |
+
−
|
52 |
+
L
|
53 |
+
|
54 |
+
|
55 |
+
|
56 |
+
d
|
57 |
+
I
|
58 |
+
|
59 |
+
|
60 |
+
d
|
61 |
+
t
|
62 |
+
|
63 |
+
|
64 |
+
|
65 |
+
|
66 |
+
|
67 |
+
{\displaystyle U=\Delta V=-L{\frac {dI}{dt}}}
|
68 |
+
|
69 |
+
despite the fact that, internally, the electric field in the coil is zero[6] (assuming it is a perfect conductor).
|
70 |
+
|
71 |
+
Using the above definition, the electric potential is not defined whenever magnetic fields change with time. In physics, it's sometimes useful to generalize the electric potential by only considering the conservative part of the electric field. This is done by the following decomposition used in electrodynamics:
|
72 |
+
|
73 |
+
E
|
74 |
+
→
|
75 |
+
|
76 |
+
|
77 |
+
|
78 |
+
=
|
79 |
+
−
|
80 |
+
∇
|
81 |
+
V
|
82 |
+
−
|
83 |
+
|
84 |
+
|
85 |
+
|
86 |
+
∂
|
87 |
+
|
88 |
+
|
89 |
+
|
90 |
+
A
|
91 |
+
→
|
92 |
+
|
93 |
+
|
94 |
+
|
95 |
+
|
96 |
+
|
97 |
+
∂
|
98 |
+
t
|
99 |
+
|
100 |
+
|
101 |
+
|
102 |
+
|
103 |
+
|
104 |
+
{\displaystyle {\vec {E}}=-\nabla V-{\frac {\partial {\vec {A}}}{\partial t}}}
|
105 |
+
|
106 |
+
where
|
107 |
+
|
108 |
+
|
109 |
+
|
110 |
+
|
111 |
+
|
112 |
+
|
113 |
+
A
|
114 |
+
→
|
115 |
+
|
116 |
+
|
117 |
+
|
118 |
+
|
119 |
+
|
120 |
+
{\textstyle {\vec {A}}}
|
121 |
+
|
122 |
+
is the magnetic vector potential. The above decomposition is justified by Helmholtz's theorem.
|
123 |
+
|
124 |
+
In this case, the voltage increase from
|
125 |
+
|
126 |
+
|
127 |
+
|
128 |
+
|
129 |
+
x
|
130 |
+
|
131 |
+
A
|
132 |
+
|
133 |
+
|
134 |
+
|
135 |
+
|
136 |
+
{\textstyle x_{A}}
|
137 |
+
|
138 |
+
to
|
139 |
+
|
140 |
+
|
141 |
+
|
142 |
+
|
143 |
+
x
|
144 |
+
|
145 |
+
B
|
146 |
+
|
147 |
+
|
148 |
+
|
149 |
+
|
150 |
+
{\textstyle x_{B}}
|
151 |
+
|
152 |
+
is given by
|
153 |
+
|
154 |
+
Δ
|
155 |
+
|
156 |
+
V
|
157 |
+
|
158 |
+
A
|
159 |
+
B
|
160 |
+
|
161 |
+
|
162 |
+
|
163 |
+
|
164 |
+
|
165 |
+
=
|
166 |
+
−
|
167 |
+
|
168 |
+
∫
|
169 |
+
|
170 |
+
|
171 |
+
x
|
172 |
+
|
173 |
+
A
|
174 |
+
|
175 |
+
|
176 |
+
|
177 |
+
|
178 |
+
|
179 |
+
x
|
180 |
+
|
181 |
+
B
|
182 |
+
|
183 |
+
|
184 |
+
|
185 |
+
|
186 |
+
|
187 |
+
|
188 |
+
|
189 |
+
|
190 |
+
E
|
191 |
+
→
|
192 |
+
|
193 |
+
|
194 |
+
|
195 |
+
|
196 |
+
|
197 |
+
c
|
198 |
+
o
|
199 |
+
n
|
200 |
+
s
|
201 |
+
e
|
202 |
+
r
|
203 |
+
v
|
204 |
+
a
|
205 |
+
t
|
206 |
+
i
|
207 |
+
v
|
208 |
+
e
|
209 |
+
|
210 |
+
|
211 |
+
|
212 |
+
⋅
|
213 |
+
d
|
214 |
+
|
215 |
+
|
216 |
+
|
217 |
+
l
|
218 |
+
→
|
219 |
+
|
220 |
+
|
221 |
+
|
222 |
+
|
223 |
+
|
224 |
+
|
225 |
+
|
226 |
+
|
227 |
+
|
228 |
+
=
|
229 |
+
−
|
230 |
+
|
231 |
+
∫
|
232 |
+
|
233 |
+
|
234 |
+
x
|
235 |
+
|
236 |
+
A
|
237 |
+
|
238 |
+
|
239 |
+
|
240 |
+
|
241 |
+
|
242 |
+
x
|
243 |
+
|
244 |
+
B
|
245 |
+
|
246 |
+
|
247 |
+
|
248 |
+
|
249 |
+
|
250 |
+
(
|
251 |
+
|
252 |
+
|
253 |
+
|
254 |
+
|
255 |
+
E
|
256 |
+
→
|
257 |
+
|
258 |
+
|
259 |
+
|
260 |
+
+
|
261 |
+
|
262 |
+
|
263 |
+
|
264 |
+
∂
|
265 |
+
|
266 |
+
|
267 |
+
|
268 |
+
A
|
269 |
+
→
|
270 |
+
|
271 |
+
|
272 |
+
|
273 |
+
|
274 |
+
|
275 |
+
∂
|
276 |
+
t
|
277 |
+
|
278 |
+
|
279 |
+
|
280 |
+
|
281 |
+
)
|
282 |
+
|
283 |
+
⋅
|
284 |
+
d
|
285 |
+
|
286 |
+
|
287 |
+
|
288 |
+
l
|
289 |
+
→
|
290 |
+
|
291 |
+
|
292 |
+
|
293 |
+
|
294 |
+
|
295 |
+
|
296 |
+
|
297 |
+
|
298 |
+
|
299 |
+
=
|
300 |
+
−
|
301 |
+
|
302 |
+
∫
|
303 |
+
|
304 |
+
|
305 |
+
x
|
306 |
+
|
307 |
+
A
|
308 |
+
|
309 |
+
|
310 |
+
|
311 |
+
|
312 |
+
|
313 |
+
x
|
314 |
+
|
315 |
+
B
|
316 |
+
|
317 |
+
|
318 |
+
|
319 |
+
|
320 |
+
(
|
321 |
+
|
322 |
+
|
323 |
+
|
324 |
+
E
|
325 |
+
→
|
326 |
+
|
327 |
+
|
328 |
+
|
329 |
+
−
|
330 |
+
|
331 |
+
|
332 |
+
|
333 |
+
|
334 |
+
E
|
335 |
+
→
|
336 |
+
|
337 |
+
|
338 |
+
|
339 |
+
|
340 |
+
|
341 |
+
i
|
342 |
+
n
|
343 |
+
d
|
344 |
+
u
|
345 |
+
c
|
346 |
+
e
|
347 |
+
d
|
348 |
+
|
349 |
+
|
350 |
+
|
351 |
+
)
|
352 |
+
⋅
|
353 |
+
d
|
354 |
+
|
355 |
+
|
356 |
+
|
357 |
+
l
|
358 |
+
→
|
359 |
+
|
360 |
+
|
361 |
+
|
362 |
+
|
363 |
+
|
364 |
+
|
365 |
+
|
366 |
+
|
367 |
+
|
368 |
+
{\displaystyle {\begin{aligned}\Delta V_{AB}&=-\int _{x_{A}}^{x_{B}}{\vec {E}}_{\mathrm {conservative} }\cdot d{\vec {l}}\\&=-\int _{x_{A}}^{x_{B}}\left({\vec {E}}+{\frac {\partial {\vec {A}}}{\partial t}}\right)\cdot d{\vec {l}}\\&=-\int _{x_{A}}^{x_{B}}({\vec {E}}-{\vec {E}}_{\mathrm {induced} })\cdot d{\vec {l}}\end{aligned}}}
|
369 |
+
|
370 |
+
where
|
371 |
+
|
372 |
+
|
373 |
+
|
374 |
+
|
375 |
+
|
376 |
+
|
377 |
+
|
378 |
+
E
|
379 |
+
→
|
380 |
+
|
381 |
+
|
382 |
+
|
383 |
+
|
384 |
+
|
385 |
+
i
|
386 |
+
n
|
387 |
+
d
|
388 |
+
u
|
389 |
+
c
|
390 |
+
e
|
391 |
+
d
|
392 |
+
|
393 |
+
|
394 |
+
|
395 |
+
|
396 |
+
|
397 |
+
{\textstyle {\vec {E}}_{\mathrm {induced} }}
|
398 |
+
|
399 |
+
is the rotational electric field due to time-varying magnetic fields. In this case, the voltage between points is always uniquely defined.
|
400 |
+
|
401 |
+
In circuit analysis and electrical engineering, the voltage across an inductor is not considered to be zero or undefined, as the standard definition would suggest. This is because electrical engineers use a lumped element model to represent and analyze circuits.
|
402 |
+
|
403 |
+
When using a lumped element model, it is assumed that there are no magnetic fields in the region surrounding the circuit and that the effects of these are contained in 'lumped elements', which are idealized and self-contained circuit elements used to model physical components.[7] If the assumption of negligible leaked fields is too inaccurate, their effects can be modelled by parasitic components.
|
404 |
+
|
405 |
+
In the case of a physical inductor though, the ideal lumped representation is often accurate. This is because the leaked fields of the inductor are generally negligible, especially if the inductor is a toroid. If leaked fields are negligible, we find that
|
406 |
+
|
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+
∫
|
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+
|
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+
|
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|
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|
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|
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|
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|
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|
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|
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|
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l
|
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+
→
|
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=
|
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|
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|
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|
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|
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+
|
456 |
+
{\displaystyle \int _{\mathrm {exterior} }{\vec {E}}\cdot d{\vec {l}}=-L{\frac {dI}{dt}}}
|
457 |
+
|
458 |
+
is path-independent, and there is a well-defined voltage across the inductor's terminals.[6] This is the reason that measurements with a voltmeter across an inductor are often reasonably independent of the placement of the test leads.
|
459 |
+
|
460 |
+
The volt (symbol: V) is the derived unit for electric potential, electric potential difference, and electromotive force. The volt is named in honour of the Italian physicist Alessandro Volta (1745–1827), who invented the voltaic pile, possibly the first chemical battery.
|
461 |
+
|
462 |
+
A simple analogy for an electric circuit is water flowing in a closed circuit of pipework, driven by a mechanical pump. This can be called a "water circuit". Potential difference between two points corresponds to the pressure difference between two points. If the pump creates a pressure difference between two points, then water flowing from one point to the other will be able to do work, such as driving a turbine. Similarly, work can be done by an electric current driven by the potential difference provided by a battery. For example, the voltage provided by a sufficiently-charged automobile battery can "push" a large current through the windings of an automobile's starter motor. If the pump isn't working, it produces no pressure difference, and the turbine will not rotate. Likewise, if the automobile's battery is very weak or "dead" (or "flat"), then it will not turn the starter motor.
|
463 |
+
|
464 |
+
The hydraulic analogy is a useful way of understanding many electrical concepts. In such a system, the work done to move water is equal to the pressure multiplied by the volume of water moved. Similarly, in an electrical circuit, the work done to move electrons or other charge-carriers is equal to "electrical pressure" multiplied by the quantity of electrical charges moved. In relation to "flow", the larger the "pressure difference" between two points (potential difference or water pressure difference), the greater the flow between them (electric current or water flow). (See "electric power".)
|
465 |
+
|
466 |
+
Specifying a voltage measurement requires explicit or implicit specification of the points across which the voltage is measured. When using a voltmeter to measure potential difference, one electrical lead of the voltmeter must be connected to the first point, one to the second point.
|
467 |
+
|
468 |
+
A common use of the term "voltage" is in describing the voltage dropped across an electrical device (such as a resistor). The voltage drop across the device can be understood as the difference between measurements at each terminal of the device with respect to a common reference point (or ground). The voltage drop is the difference between the two readings. Two points in an electric circuit that are connected by an ideal conductor without resistance and not within a changing magnetic field have a voltage of zero. Any two points with the same potential may be connected by a conductor and no current will flow between them.
|
469 |
+
|
470 |
+
The voltage between A and C is the sum of the voltage between A and B and the voltage between B and C. The various voltages in a circuit can be computed using Kirchhoff's circuit laws.
|
471 |
+
|
472 |
+
When talking about alternating current (AC) there is a difference between instantaneous voltage and average voltage. Instantaneous voltages can be added for direct current (DC) and AC, but average voltages can be meaningfully added only when they apply to signals that all have the same frequency and phase.
|
473 |
+
|
474 |
+
Instruments for measuring voltages include the voltmeter, the potentiometer, and the oscilloscope. Analog voltmeters, such as moving-coil instruments, work by measuring the current through a fixed resistor, which, according to Ohm's Law, is proportional to the voltage across the resistor. The potentiometer works by balancing the unknown voltage against a known voltage in a bridge circuit. The cathode-ray oscilloscope works by amplifying the voltage and using it to deflect an electron beam from a straight path, so that the deflection of the beam is proportional to the voltage.
|
475 |
+
|
476 |
+
A common voltage for flashlight batteries is 1.5 volts (DC).
|
477 |
+
A common voltage for automobile batteries is 12 volts (DC).
|
478 |
+
|
479 |
+
Common voltages supplied by power companies to consumers are 110 to 120 volts (AC) and 220 to 240 volts (AC). The voltage in electric power transmission lines used to distribute electricity from power stations can be several hundred times greater than consumer voltages, typically 110 to 1200 kV (AC).
|
480 |
+
|
481 |
+
The voltage used in overhead lines to power railway locomotives is between 12 kV and 50 kV (AC) or between 0.75 kV and 3 kV (DC).
|
482 |
+
|
483 |
+
Inside a conductive material, the energy of an electron is affected not only by the average electric potential, but also by the specific thermal and atomic environment that it is in.
|
484 |
+
When a voltmeter is connected between two different types of metal, it measures not the electrostatic potential difference, but instead something else that is affected by thermodynamics.[8]
|
485 |
+
The quantity measured by a voltmeter is the negative of the difference of the electrochemical potential of electrons (Fermi level) divided by the electron charge and commonly referred to as the voltage difference, while the pure unadjusted electrostatic potential (not measurable with a voltmeter) is sometimes called Galvani potential.
|
486 |
+
The terms "voltage" and "electric potential" are ambiguous in that, in practice, they can refer to either of these in different contexts.
|
487 |
+
|
488 |
+
The term electromotive force was first used by Volta in a letter to Giovanni Aldini in 1798, and first appeared in a published paper in 1801 in Annales de chimie et de physique.[9]:408 Volta meant by this a force that was not an electrostatic force, specifically, an electrochemical force.[9]:405 The term was taken up by Michael Faraday in connection with electromagnetic induction in the 1820s. However, a clear definition of voltage and method of measuring it had not been developed at this time.[10]:554 Volta distinguished electromotive force (emf) from tension (potential difference): the observed potential difference at the terminals of an electrochemical cell when it was open circuit must exactly balance the emf of the cell so that no current flowed.[9]:405
|
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|
1 |
+
Voltage, electric potential difference, electric pressure or electric tension is the difference in electric potential between two points, which (in a static electric field) is defined as the work needed per unit of charge to move a test charge between the two points. In the International System of Units, the derived unit for voltage (potential difference) is named volt.[1]:166 In SI units, work per unit charge is expressed as joules per coulomb, where 1 volt = 1 joule (of work) per 1 coulomb (of charge). The old SI definition for volt used power and current, since 1990 quantum Hall and Josephson effect were used and recently (2019) fundamental physical constants were introduced for the definition of all SI units and derived units.[1]:177f, 197f Voltage or electric potential difference is denoted symbolically by ∆V, simplified V,[2] or U,[3] for instance in the context of Ohm's or Kirchhoff's circuit laws.
|
2 |
+
|
3 |
+
Electric potential differences between points can be caused by electric charge, by electric current through a magnetic field, by time-varying magnetic fields, or some combination of these three.[4][5] A voltmeter can be used to measure the voltage (or potential difference) between two points in a system; often a common reference potential such as the ground of the system is used as one of the points. A voltage may represent either a source of energy (electromotive force) or lost, used, or stored energy (potential drop).
|
4 |
+
|
5 |
+
There are multiple useful ways to define voltage, including the standard definition mentioned at the start of this page. There are also other useful definitions of work per charge (see this section).
|
6 |
+
|
7 |
+
Voltage is defined so that negatively charged objects are pulled towards higher voltages, while positively charged objects are pulled towards lower voltages. Therefore, the conventional current in a wire or resistor always flows from higher voltage to lower voltage.
|
8 |
+
|
9 |
+
Historically, voltage has been referred to using terms like "tension" and "pressure". Even today, the term "tension" is still used, for example within the phrase "high tension" (HT) which is commonly used in thermionic valve (vacuum tube) based electronics.
|
10 |
+
|
11 |
+
The voltage increase from some point
|
12 |
+
|
13 |
+
|
14 |
+
|
15 |
+
|
16 |
+
x
|
17 |
+
|
18 |
+
A
|
19 |
+
|
20 |
+
|
21 |
+
|
22 |
+
|
23 |
+
{\textstyle x_{A}}
|
24 |
+
|
25 |
+
to some point
|
26 |
+
|
27 |
+
|
28 |
+
|
29 |
+
|
30 |
+
x
|
31 |
+
|
32 |
+
B
|
33 |
+
|
34 |
+
|
35 |
+
|
36 |
+
|
37 |
+
{\textstyle x_{B}}
|
38 |
+
|
39 |
+
is given by
|
40 |
+
|
41 |
+
In this case, the voltage increase from point A to point B is equal to the work done per unit charge, against the electric field, to move the charge from A to B without causing any acceleration. Mathematically, this is expressed as the line integral of the electric field along that path.
|
42 |
+
Under this definition, the voltage difference between two points is not uniquely defined when there are time-varying magnetic fields since the electric force is not a conservative force in such cases.
|
43 |
+
|
44 |
+
If this definition of voltage is used, any circuit where there are time-varying magnetic fields,[note 1] such as circuits containing inductors, will not have a well-defined voltage between nodes in the circuit. However, if magnetic fields are suitably contained to each component, then the electric field is conservative in the region exterior[note 2] to the components, and voltages are well-defined in that region.[6] In this case, the voltage across an inductor, viewed externally, turns out to be
|
45 |
+
|
46 |
+
U
|
47 |
+
=
|
48 |
+
Δ
|
49 |
+
V
|
50 |
+
=
|
51 |
+
−
|
52 |
+
L
|
53 |
+
|
54 |
+
|
55 |
+
|
56 |
+
d
|
57 |
+
I
|
58 |
+
|
59 |
+
|
60 |
+
d
|
61 |
+
t
|
62 |
+
|
63 |
+
|
64 |
+
|
65 |
+
|
66 |
+
|
67 |
+
{\displaystyle U=\Delta V=-L{\frac {dI}{dt}}}
|
68 |
+
|
69 |
+
despite the fact that, internally, the electric field in the coil is zero[6] (assuming it is a perfect conductor).
|
70 |
+
|
71 |
+
Using the above definition, the electric potential is not defined whenever magnetic fields change with time. In physics, it's sometimes useful to generalize the electric potential by only considering the conservative part of the electric field. This is done by the following decomposition used in electrodynamics:
|
72 |
+
|
73 |
+
E
|
74 |
+
→
|
75 |
+
|
76 |
+
|
77 |
+
|
78 |
+
=
|
79 |
+
−
|
80 |
+
∇
|
81 |
+
V
|
82 |
+
−
|
83 |
+
|
84 |
+
|
85 |
+
|
86 |
+
∂
|
87 |
+
|
88 |
+
|
89 |
+
|
90 |
+
A
|
91 |
+
→
|
92 |
+
|
93 |
+
|
94 |
+
|
95 |
+
|
96 |
+
|
97 |
+
∂
|
98 |
+
t
|
99 |
+
|
100 |
+
|
101 |
+
|
102 |
+
|
103 |
+
|
104 |
+
{\displaystyle {\vec {E}}=-\nabla V-{\frac {\partial {\vec {A}}}{\partial t}}}
|
105 |
+
|
106 |
+
where
|
107 |
+
|
108 |
+
|
109 |
+
|
110 |
+
|
111 |
+
|
112 |
+
|
113 |
+
A
|
114 |
+
→
|
115 |
+
|
116 |
+
|
117 |
+
|
118 |
+
|
119 |
+
|
120 |
+
{\textstyle {\vec {A}}}
|
121 |
+
|
122 |
+
is the magnetic vector potential. The above decomposition is justified by Helmholtz's theorem.
|
123 |
+
|
124 |
+
In this case, the voltage increase from
|
125 |
+
|
126 |
+
|
127 |
+
|
128 |
+
|
129 |
+
x
|
130 |
+
|
131 |
+
A
|
132 |
+
|
133 |
+
|
134 |
+
|
135 |
+
|
136 |
+
{\textstyle x_{A}}
|
137 |
+
|
138 |
+
to
|
139 |
+
|
140 |
+
|
141 |
+
|
142 |
+
|
143 |
+
x
|
144 |
+
|
145 |
+
B
|
146 |
+
|
147 |
+
|
148 |
+
|
149 |
+
|
150 |
+
{\textstyle x_{B}}
|
151 |
+
|
152 |
+
is given by
|
153 |
+
|
154 |
+
Δ
|
155 |
+
|
156 |
+
V
|
157 |
+
|
158 |
+
A
|
159 |
+
B
|
160 |
+
|
161 |
+
|
162 |
+
|
163 |
+
|
164 |
+
|
165 |
+
=
|
166 |
+
−
|
167 |
+
|
168 |
+
∫
|
169 |
+
|
170 |
+
|
171 |
+
x
|
172 |
+
|
173 |
+
A
|
174 |
+
|
175 |
+
|
176 |
+
|
177 |
+
|
178 |
+
|
179 |
+
x
|
180 |
+
|
181 |
+
B
|
182 |
+
|
183 |
+
|
184 |
+
|
185 |
+
|
186 |
+
|
187 |
+
|
188 |
+
|
189 |
+
|
190 |
+
E
|
191 |
+
→
|
192 |
+
|
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|
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|
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|
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|
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|
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|
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|
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⋅
|
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|
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l
|
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→
|
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|
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|
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=
|
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−
|
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|
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∫
|
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|
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|
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x
|
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|
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A
|
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|
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|
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|
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|
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|
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x
|
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B
|
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|
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(
|
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|
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|
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|
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E
|
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→
|
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|
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∂
|
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A
|
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→
|
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∂
|
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t
|
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)
|
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⋅
|
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|
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l
|
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→
|
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|
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=
|
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−
|
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|
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∫
|
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|
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|
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x
|
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|
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A
|
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|
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|
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x
|
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B
|
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(
|
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E
|
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→
|
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|
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|
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−
|
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E
|
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→
|
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|
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⋅
|
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|
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|
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|
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{\displaystyle {\begin{aligned}\Delta V_{AB}&=-\int _{x_{A}}^{x_{B}}{\vec {E}}_{\mathrm {conservative} }\cdot d{\vec {l}}\\&=-\int _{x_{A}}^{x_{B}}\left({\vec {E}}+{\frac {\partial {\vec {A}}}{\partial t}}\right)\cdot d{\vec {l}}\\&=-\int _{x_{A}}^{x_{B}}({\vec {E}}-{\vec {E}}_{\mathrm {induced} })\cdot d{\vec {l}}\end{aligned}}}
|
369 |
+
|
370 |
+
where
|
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|
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→
|
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|
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|
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|
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|
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{\textstyle {\vec {E}}_{\mathrm {induced} }}
|
398 |
+
|
399 |
+
is the rotational electric field due to time-varying magnetic fields. In this case, the voltage between points is always uniquely defined.
|
400 |
+
|
401 |
+
In circuit analysis and electrical engineering, the voltage across an inductor is not considered to be zero or undefined, as the standard definition would suggest. This is because electrical engineers use a lumped element model to represent and analyze circuits.
|
402 |
+
|
403 |
+
When using a lumped element model, it is assumed that there are no magnetic fields in the region surrounding the circuit and that the effects of these are contained in 'lumped elements', which are idealized and self-contained circuit elements used to model physical components.[7] If the assumption of negligible leaked fields is too inaccurate, their effects can be modelled by parasitic components.
|
404 |
+
|
405 |
+
In the case of a physical inductor though, the ideal lumped representation is often accurate. This is because the leaked fields of the inductor are generally negligible, especially if the inductor is a toroid. If leaked fields are negligible, we find that
|
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+
|
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∫
|
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|
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|
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|
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|
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|
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|
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|
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|
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E
|
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→
|
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|
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|
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⋅
|
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|
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|
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→
|
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|
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|
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|
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=
|
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|
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L
|
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|
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|
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|
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I
|
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|
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|
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|
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|
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|
454 |
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|
455 |
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|
456 |
+
{\displaystyle \int _{\mathrm {exterior} }{\vec {E}}\cdot d{\vec {l}}=-L{\frac {dI}{dt}}}
|
457 |
+
|
458 |
+
is path-independent, and there is a well-defined voltage across the inductor's terminals.[6] This is the reason that measurements with a voltmeter across an inductor are often reasonably independent of the placement of the test leads.
|
459 |
+
|
460 |
+
The volt (symbol: V) is the derived unit for electric potential, electric potential difference, and electromotive force. The volt is named in honour of the Italian physicist Alessandro Volta (1745–1827), who invented the voltaic pile, possibly the first chemical battery.
|
461 |
+
|
462 |
+
A simple analogy for an electric circuit is water flowing in a closed circuit of pipework, driven by a mechanical pump. This can be called a "water circuit". Potential difference between two points corresponds to the pressure difference between two points. If the pump creates a pressure difference between two points, then water flowing from one point to the other will be able to do work, such as driving a turbine. Similarly, work can be done by an electric current driven by the potential difference provided by a battery. For example, the voltage provided by a sufficiently-charged automobile battery can "push" a large current through the windings of an automobile's starter motor. If the pump isn't working, it produces no pressure difference, and the turbine will not rotate. Likewise, if the automobile's battery is very weak or "dead" (or "flat"), then it will not turn the starter motor.
|
463 |
+
|
464 |
+
The hydraulic analogy is a useful way of understanding many electrical concepts. In such a system, the work done to move water is equal to the pressure multiplied by the volume of water moved. Similarly, in an electrical circuit, the work done to move electrons or other charge-carriers is equal to "electrical pressure" multiplied by the quantity of electrical charges moved. In relation to "flow", the larger the "pressure difference" between two points (potential difference or water pressure difference), the greater the flow between them (electric current or water flow). (See "electric power".)
|
465 |
+
|
466 |
+
Specifying a voltage measurement requires explicit or implicit specification of the points across which the voltage is measured. When using a voltmeter to measure potential difference, one electrical lead of the voltmeter must be connected to the first point, one to the second point.
|
467 |
+
|
468 |
+
A common use of the term "voltage" is in describing the voltage dropped across an electrical device (such as a resistor). The voltage drop across the device can be understood as the difference between measurements at each terminal of the device with respect to a common reference point (or ground). The voltage drop is the difference between the two readings. Two points in an electric circuit that are connected by an ideal conductor without resistance and not within a changing magnetic field have a voltage of zero. Any two points with the same potential may be connected by a conductor and no current will flow between them.
|
469 |
+
|
470 |
+
The voltage between A and C is the sum of the voltage between A and B and the voltage between B and C. The various voltages in a circuit can be computed using Kirchhoff's circuit laws.
|
471 |
+
|
472 |
+
When talking about alternating current (AC) there is a difference between instantaneous voltage and average voltage. Instantaneous voltages can be added for direct current (DC) and AC, but average voltages can be meaningfully added only when they apply to signals that all have the same frequency and phase.
|
473 |
+
|
474 |
+
Instruments for measuring voltages include the voltmeter, the potentiometer, and the oscilloscope. Analog voltmeters, such as moving-coil instruments, work by measuring the current through a fixed resistor, which, according to Ohm's Law, is proportional to the voltage across the resistor. The potentiometer works by balancing the unknown voltage against a known voltage in a bridge circuit. The cathode-ray oscilloscope works by amplifying the voltage and using it to deflect an electron beam from a straight path, so that the deflection of the beam is proportional to the voltage.
|
475 |
+
|
476 |
+
A common voltage for flashlight batteries is 1.5 volts (DC).
|
477 |
+
A common voltage for automobile batteries is 12 volts (DC).
|
478 |
+
|
479 |
+
Common voltages supplied by power companies to consumers are 110 to 120 volts (AC) and 220 to 240 volts (AC). The voltage in electric power transmission lines used to distribute electricity from power stations can be several hundred times greater than consumer voltages, typically 110 to 1200 kV (AC).
|
480 |
+
|
481 |
+
The voltage used in overhead lines to power railway locomotives is between 12 kV and 50 kV (AC) or between 0.75 kV and 3 kV (DC).
|
482 |
+
|
483 |
+
Inside a conductive material, the energy of an electron is affected not only by the average electric potential, but also by the specific thermal and atomic environment that it is in.
|
484 |
+
When a voltmeter is connected between two different types of metal, it measures not the electrostatic potential difference, but instead something else that is affected by thermodynamics.[8]
|
485 |
+
The quantity measured by a voltmeter is the negative of the difference of the electrochemical potential of electrons (Fermi level) divided by the electron charge and commonly referred to as the voltage difference, while the pure unadjusted electrostatic potential (not measurable with a voltmeter) is sometimes called Galvani potential.
|
486 |
+
The terms "voltage" and "electric potential" are ambiguous in that, in practice, they can refer to either of these in different contexts.
|
487 |
+
|
488 |
+
The term electromotive force was first used by Volta in a letter to Giovanni Aldini in 1798, and first appeared in a published paper in 1801 in Annales de chimie et de physique.[9]:408 Volta meant by this a force that was not an electrostatic force, specifically, an electrochemical force.[9]:405 The term was taken up by Michael Faraday in connection with electromagnetic induction in the 1820s. However, a clear definition of voltage and method of measuring it had not been developed at this time.[10]:554 Volta distinguished electromotive force (emf) from tension (potential difference): the observed potential difference at the terminals of an electrochemical cell when it was open circuit must exactly balance the emf of the cell so that no current flowed.[9]:405
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|
1 |
+
|
2 |
+
|
3 |
+
|
4 |
+
|
5 |
+
Martialinae
|
6 |
+
|
7 |
+
Leptanillinae
|
8 |
+
|
9 |
+
Amblyoponinae
|
10 |
+
|
11 |
+
Paraponerinae
|
12 |
+
|
13 |
+
Agroecomyrmecinae
|
14 |
+
|
15 |
+
Ponerinae
|
16 |
+
|
17 |
+
Proceratiinae
|
18 |
+
|
19 |
+
Ecitoninae‡
|
20 |
+
|
21 |
+
Aenictinae‡
|
22 |
+
|
23 |
+
Dorylini‡
|
24 |
+
|
25 |
+
Aenictogitoninae‡
|
26 |
+
|
27 |
+
Cerapachyinae‡*
|
28 |
+
|
29 |
+
Leptanilloidinae‡
|
30 |
+
|
31 |
+
Dolichoderinae
|
32 |
+
|
33 |
+
Aneuretinae
|
34 |
+
|
35 |
+
Pseudomyrmecinae
|
36 |
+
|
37 |
+
Myrmeciinae
|
38 |
+
|
39 |
+
Ectatomminae
|
40 |
+
|
41 |
+
Heteroponerinae
|
42 |
+
|
43 |
+
Myrmicinae
|
44 |
+
|
45 |
+
Formicinae
|
46 |
+
|
47 |
+
A phylogeny of the extant ant subfamilies.[2][3]
|
48 |
+
*Cerapachyinae is paraphyletic
|
49 |
+
‡ The previous dorylomorph subfamilies were synonymized under Dorylinae by Brady et al. in 2014[4]
|
50 |
+
|
51 |
+
Ants are eusocial insects of the family Formicidae and, along with the related wasps and bees, belong to the order Hymenoptera. Ants appear in the fossil record across the globe in considerable diversity during the latest Early Cretaceous and early Late Cretaceous, suggesting an earlier origin. Ants evolved from vespoid wasp ancestors in the Cretaceous period, and diversified after the rise of flowering plants. More than 12,500 of an estimated total of 22,000 species have been classified. They are easily identified by their elbowed antennae and the distinctive node-like structure that forms their slender waists.
|
52 |
+
|
53 |
+
Ants form colonies that range in size from a few dozen predatory individuals living in small natural cavities to highly organised colonies that may occupy large territories and consist of millions of individuals. Larger colonies consist of various castes of sterile, wingless females, most of which are workers (ergates), as well as soldiers (dinergates) and other specialised groups. Nearly all ant colonies also have some fertile males called "drones" (aner) and one or more fertile females called "queens" (gynes). The colonies are described as superorganisms because the ants appear to operate as a unified entity, collectively working together to support the colony.
|
54 |
+
|
55 |
+
Ants have colonised almost every landmass on Earth. The only places lacking indigenous ants are Antarctica and a few remote or inhospitable islands. Ants thrive in most ecosystems and may form 15–25% of the terrestrial animal biomass. Their success in so many environments has been attributed to their social organisation and their ability to modify habitats, tap resources, and defend themselves. Their long co-evolution with other species has led to mimetic, commensal, parasitic, and mutualistic relationships.
|
56 |
+
|
57 |
+
Ant societies have division of labour, communication between individuals, and an ability to solve complex problems. These parallels with human societies have long been an inspiration and subject of study. Many human cultures make use of ants in cuisine, medication, and rituals. Some species are valued in their role as biological pest control agents. Their ability to exploit resources may bring ants into conflict with humans, however, as they can damage crops and invade buildings. Some species, such as the red imported fire ant (Solenopsis invicta), are regarded as invasive species, establishing themselves in areas where they have been introduced accidentally.
|
58 |
+
|
59 |
+
The word ant and its chiefly dialectal form emmet[5] come from ante, emete of Middle English, which come from ǣmette of Old English, and these are all related to the dialectal Dutch emt and the Old High German āmeiza, from which comes the modern German Ameise. All of these words come from West Germanic *ǣmaitjōn, and the original meaning of the word was "the biter" (from Proto-Germanic *ai-, "off, away" + *mait- "cut").[6][7] The family name Formicidae is derived from the Latin formīca ("ant")[8] from which the words in other Romance languages, such as the Portuguese formiga, Italian formica, Spanish hormiga, Romanian furnică, and French fourmi are derived. It has been hypothesised that a Proto-Indo-European word *morwi- was used, cf. Sanskrit vamrah, Latin formīca, Greek μύρμηξ mýrmēx, Old Church Slavonic mraviji, Old Irish moirb, Old Norse maurr, Dutch mier.[9]
|
60 |
+
|
61 |
+
Chrysidoidea
|
62 |
+
|
63 |
+
Vespidae
|
64 |
+
|
65 |
+
Rhopalosomatidae
|
66 |
+
|
67 |
+
Pompilidae
|
68 |
+
|
69 |
+
Tiphiidae
|
70 |
+
|
71 |
+
Scolioidea
|
72 |
+
|
73 |
+
Apoidea
|
74 |
+
|
75 |
+
Formicidae
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The family Formicidae belongs to the order Hymenoptera, which also includes sawflies, bees, and wasps. Ants evolved from a lineage within the stinging wasps, and a 2013 study suggests that they are a sister group of the Apoidea.[10] In 1966, E. O. Wilson and his colleagues identified the fossil remains of an ant (Sphecomyrma) that lived in the Cretaceous period. The specimen, trapped in amber dating back to around 92 million years ago, has features found in some wasps, but not found in modern ants.[11] Sphecomyrma was possibly a ground forager, while Haidomyrmex and Haidomyrmodes, related genera in subfamily Sphecomyrminae, are reconstructed as active arboreal predators.[12] Older ants in the genus Sphecomyrmodes have been found in 99 million year-old amber from Myanmar.[13][14] A 2006 study suggested that ants arose tens of millions of years earlier than previously thought, up to 168 million years ago.[1] After the rise of flowering plants about 100 million years ago they diversified and assumed ecological dominance around 60 million years ago.[15][1][16][17] Some groups, such as the Leptanillinae and Martialinae, are suggested to have diversified from early primitive ants that were likely to have been predators underneath the surface of the soil.[3][18]
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During the Cretaceous period, a few species of primitive ants ranged widely on the Laurasian supercontinent (the Northern Hemisphere). They were scarce in comparison to the populations of other insects, representing only about 1% of the entire insect population. Ants became dominant after adaptive radiation at the beginning of the Paleogene period. By the Oligocene and Miocene, ants had come to represent 20–40% of all insects found in major fossil deposits. Of the species that lived in the Eocene epoch, around one in 10 genera survive to the present. Genera surviving today comprise 56% of the genera in Baltic amber fossils (early Oligocene), and 92% of the genera in Dominican amber fossils (apparently early Miocene).[15][19]
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Termites live in colonies and are sometimes called ‘white ants’, but termites are not ants. They are the sub-order Isoptera, and together with cockroaches they form the order Blattodea. Blattodeans are related to mantids, crickets, and other winged insects that do not undergo full metamorphosis. Like ants, termites are eusocial, with sterile workers, but they differ greatly in the genetics of reproduction. The similarity of their social structure to that of ants is attributed to convergent evolution.[20] Velvet ants look like large ants, but are wingless female wasps.[21][22]
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Ants are found on all continents except Antarctica, and only a few large islands, such as Greenland, Iceland, parts of Polynesia and the Hawaiian Islands lack native ant species.[24][25] Ants occupy a wide range of ecological niches and exploit many different food resources as direct or indirect herbivores, predators and scavengers. Most ant species are omnivorous generalists, but a few are specialist feeders. Their ecological dominance is demonstrated by their biomass: ants are estimated to contribute 15–20 % (on average and nearly 25% in the tropics) of terrestrial animal biomass, exceeding that of the vertebrates.[26]
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Ants range in size from 0.75 to 52 millimetres (0.030–2.0 in),[27][28] the largest species being the fossil Titanomyrma giganteum, the queen of which was 6 centimetres (2.4 in) long with a wingspan of 15 centimetres (5.9 in).[29] Ants vary in colour; most ants are red or black, but a few species are green and some tropical species have a metallic lustre. More than 12,000 species are currently known (with upper estimates of the potential existence of about 22,000) (see the article List of ant genera), with the greatest diversity in the tropics. Taxonomic studies continue to resolve the classification and systematics of ants. Online databases of ant species, including AntBase and the Hymenoptera Name Server, help to keep track of the known and newly described species.[30] The relative ease with which ants may be sampled and studied in ecosystems has made them useful as indicator species in biodiversity studies.[31][32]
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Ants are distinct in their morphology from other insects in having elbowed antennae, metapleural glands, and a strong constriction of their second abdominal segment into a node-like petiole. The head, mesosoma, and metasoma are the three distinct body segments (formally tagmata). The petiole forms a narrow waist between their mesosoma (thorax plus the first abdominal segment, which is fused to it) and gaster (abdomen less the abdominal segments in the petiole). The petiole may be formed by one or two nodes (the second alone, or the second and third abdominal segments).[33]
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Like other insects, ants have an exoskeleton, an external covering that provides a protective casing around the body and a point of attachment for muscles, in contrast to the internal skeletons of humans and other vertebrates. Insects do not have lungs; oxygen and other gases, such as carbon dioxide, pass through their exoskeleton via tiny valves called spiracles. Insects also lack closed blood vessels; instead, they have a long, thin, perforated tube along the top of the body (called the "dorsal aorta") that functions like a heart, and pumps haemolymph toward the head, thus driving the circulation of the internal fluids. The nervous system consists of a ventral nerve cord that runs the length of the body, with several ganglia and branches along the way reaching into the extremities of the appendages.[34]
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An ant's head contains many sensory organs. Like most insects, ants have compound eyes made from numerous tiny lenses attached together. Ant eyes are good for acute movement detection, but do not offer a high resolution image. They also have three small ocelli (simple eyes) on the top of the head that detect light levels and polarization.[35] Compared to vertebrates, ants tend to have blurrier eyesight, particularly in smaller species,[36] and a few subterranean taxa are completely blind.[2] However, some ants, such as Australia's bulldog ant, have excellent vision and are capable of discriminating the distance and size of objects moving nearly a metre away.[37]
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Two antennae ("feelers") are attached to the head; these organs detect chemicals, air currents, and vibrations; they also are used to transmit and receive signals through touch. The head has two strong jaws, the mandibles, used to carry food, manipulate objects, construct nests, and for defence.[34] In some species, a small pocket (infrabuccal chamber) inside the mouth stores food, so it may be passed to other ants or their larvae.[38]
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Both the legs and wings of the ant are attached to the mesosoma ("thorax"). The legs terminate in a hooked claw which allows them to hook on and climb surfaces.[39] Only reproductive ants, queens, and males, have wings. Queens shed their wings after the nuptial flight, leaving visible stubs, a distinguishing feature of queens. In a few species, wingless queens (ergatoids) and males occur.[34]
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The metasoma (the "abdomen") of the ant houses important internal organs, including those of the reproductive, respiratory (tracheae), and excretory systems. Workers of many species have their egg-laying structures modified into stings that are used for subduing prey and defending their nests.[34]
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In the colonies of a few ant species, there are physical castes—workers in distinct size-classes, called minor, median, and major ergates. Often, the larger ants have disproportionately larger heads, and correspondingly stronger mandibles. These are known as macrergates while smaller workers are known as micrergates.[40] Although formally known as dinergates, such individuals are sometimes called "soldier" ants because their stronger mandibles make them more effective in fighting, although they still are workers and their "duties" typically do not vary greatly from the minor or median workers. In a few species, the median workers are absent, creating a sharp divide between the minors and majors.[41] Weaver ants, for example, have a distinct bimodal size distribution.[42][43] Some other species show continuous variation in the size of workers. The smallest and largest workers in Pheidologeton diversus show nearly a 500-fold difference in their dry-weights.[44]
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Workers cannot mate; however, because of the haplodiploid sex-determination system in ants, workers of a number of species can lay unfertilised eggs that become fully fertile, haploid males. The role of workers may change with their age and in some species, such as honeypot ants, young workers are fed until their gasters are distended, and act as living food storage vessels. These food storage workers are called repletes.[45] For instance, these replete workers develop in the North American honeypot ant Myrmecocystus mexicanus. Usually the largest workers in the colony develop into repletes; and, if repletes are removed from the colony, other workers become repletes, demonstrating the flexibility of this particular polymorphism.[46] This polymorphism in morphology and behaviour of workers initially was thought to be determined by environmental factors such as nutrition and hormones that led to different developmental paths; however, genetic differences between worker castes have been noted in Acromyrmex sp.[47] These polymorphisms are caused by relatively small genetic changes; differences in a single gene of Solenopsis invicta can decide whether the colony will have single or multiple queens.[48] The Australian jack jumper ant (Myrmecia pilosula) has only a single pair of chromosomes (with the males having just one chromosome as they are haploid), the lowest number known for any animal, making it an interesting subject for studies in the genetics and developmental biology of social insects.[49][50]
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The life of an ant starts from an egg; if the egg is fertilised, the progeny will be female diploid, if not, it will be male haploid. Ants develop by complete metamorphosis with the larva stages passing through a pupal stage before emerging as an adult. The larva is largely immobile and is fed and cared for by workers. Food is given to the larvae by trophallaxis, a process in which an ant regurgitates liquid food held in its crop. This is also how adults share food, stored in the "social stomach". Larvae, especially in the later stages, may also be provided solid food, such as trophic eggs, pieces of prey, and seeds brought by workers.[51]
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The larvae grow through a series of four or five moults and enter the pupal stage. The pupa has the appendages free and not fused to the body as in a butterfly pupa.[52] The differentiation into queens and workers (which are both female), and different castes of workers, is influenced in some species by the nutrition the larvae obtain. Genetic influences and the control of gene expression by the developmental environment are complex and the determination of caste continues to be a subject of research.[53] Winged male ants, called drones, emerge from pupae along with the usually winged breeding females. Some species, such as army ants, have wingless queens. Larvae and pupae need to be kept at fairly constant temperatures to ensure proper development, and so often are moved around among the various brood chambers within the colony.[54]
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A new ergate spends the first few days of its adult life caring for the queen and young. She then graduates to digging and other nest work, and later to defending the nest and foraging. These changes are sometimes fairly sudden, and define what are called temporal castes. An explanation for the sequence is suggested by the high casualties involved in foraging, making it an acceptable risk only for ants who are older and are likely to die soon of natural causes.[55][56]
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Ant colonies can be long-lived. The queens can live for up to 30 years, and workers live from 1 to 3 years. Males, however, are more transitory, being quite short-lived and surviving for only a few weeks.[57] Ant queens are estimated to live 100 times as long as solitary insects of a similar size.[58]
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Ants are active all year long in the tropics, but, in cooler regions, they survive the winter in a state of dormancy known as hibernation. The forms of inactivity are varied and some temperate species have larvae going into the inactive state (diapause), while in others, the adults alone pass the winter in a state of reduced activity.[59]
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A wide range of reproductive strategies have been noted in ant species. Females of many species are known to be capable of reproducing asexually through thelytokous parthenogenesis.[60] Secretions from the male accessory glands in some species can plug the female genital opening and prevent females from re-mating.[61] Most ant species have a system in which only the queen and breeding females have the ability to mate. Contrary to popular belief, some ant nests have multiple queens, while others may exist without queens. Workers with the ability to reproduce are called "gamergates" and colonies that lack queens are then called gamergate colonies; colonies with queens are said to be queen-right.[62]
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Drones can also mate with existing queens by entering a foreign colony. When the drone is initially attacked by the workers, it releases a mating pheromone. If recognized as a mate, it will be carried to the queen to mate.[63] Males may also patrol the nest and fight others by grabbing them with their mandibles, piercing their exoskeleton and then marking them with a pheromone. The marked male is interpreted as an invader by worker ants and is killed.[64]
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Most ants are univoltine, producing a new generation each year.[65] During the species-specific breeding period, winged females and winged males, known to entomologists as alates, leave the colony in what is called a nuptial flight. The nuptial flight usually takes place in the late spring or early summer when the weather is hot and humid. Heat makes flying easier and freshly fallen rain makes the ground softer for mated queens to dig nests.[66] Males typically take flight before the females. Males then use visual cues to find a common mating ground, for example, a landmark such as a pine tree to which other males in the area converge. Males secrete a mating pheromone that females follow. Males will mount females in the air, but the actual mating process usually takes place on the ground. Females of some species mate with just one male but in others they may mate with as many as ten or more different males, storing the sperm in their spermathecae.[67]
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Mated females then seek a suitable place to begin a colony. There, they break off their wings and begin to lay and care for eggs. The females can selectively fertilise future eggs with the sperm stored to produce diploid workers or lay unfertilized haploid eggs to produce drones. The first workers to hatch are known as nanitics,[68] and are weaker and smaller than later workers, but they begin to serve the colony immediately. They enlarge the nest, forage for food, and care for the other eggs. Species that have multiple queens may have a queen leaving the nest along with some workers to found a colony at a new site,[67] a process akin to swarming in honeybees.
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Ants communicate with each other using pheromones, sounds, and touch.[69] The use of pheromones as chemical signals is more developed in ants, such as the red harvester ant, than in other hymenopteran groups. Like other insects, ants perceive smells with their long, thin, and mobile antennae. The paired antennae provide information about the direction and intensity of scents. Since most ants live on the ground, they use the soil surface to leave pheromone trails that may be followed by other ants. In species that forage in groups, a forager that finds food marks a trail on the way back to the colony; this trail is followed by other ants, these ants then reinforce the trail when they head back with food to the colony. When the food source is exhausted, no new trails are marked by returning ants and the scent slowly dissipates. This behaviour helps ants deal with changes in their environment. For instance, when an established path to a food source is blocked by an obstacle, the foragers leave the path to explore new routes. If an ant is successful, it leaves a new trail marking the shortest route on its return. Successful trails are followed by more ants, reinforcing better routes and gradually identifying the best path.[70]
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Ants use pheromones for more than just making trails. A crushed ant emits an alarm pheromone that sends nearby ants into an attack frenzy and attracts more ants from farther away. Several ant species even use "propaganda pheromones" to confuse enemy ants and make them fight among themselves.[71] Pheromones are produced by a wide range of structures including Dufour's glands, poison glands and glands on the hindgut, pygidium, rectum, sternum, and hind tibia.[58] Pheromones also are exchanged, mixed with food, and passed by trophallaxis, transferring information within the colony.[72] This allows other ants to detect what task group (e.g., foraging or nest maintenance) other colony members belong to.[73] In ant species with queen castes, when the dominant queen stops producing a specific pheromone, workers begin to raise new queens in the colony.[74]
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Some ants produce sounds by stridulation, using the gaster segments and their mandibles. Sounds may be used to communicate with colony members or with other species.[75][76]
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Ants attack and defend themselves by biting and, in many species, by stinging, often injecting or spraying chemicals, such as formic acid in the case of formicine ants, alkaloids and piperidines in fire ants, and a variety of protein components in other ants. Bullet ants (Paraponera), located in Central and South America, are considered to have the most painful sting of any insect, although it is usually not fatal to humans. This sting is given the highest rating on the Schmidt sting pain index.[77]
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The sting of jack jumper ants can be fatal,[78] and an antivenom has been developed for it.[79]
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Fire ants, Solenopsis spp., are unique in having a venom sac containing piperidine alkaloids.[80] Their stings are painful and can be dangerous to hypersensitive people.[81]
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Trap-jaw ants of the genus Odontomachus are equipped with mandibles called trap-jaws, which snap shut faster than any other predatory appendages within the animal kingdom.[82] One study of Odontomachus bauri recorded peak speeds of between 126 and 230 km/h (78 and 143 mph), with the jaws closing within 130 microseconds on average.
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The ants were also observed to use their jaws as a catapult to eject intruders or fling themselves backward to escape a threat.[82] Before striking, the ant opens its mandibles extremely widely and locks them in this position by an internal mechanism. Energy is stored in a thick band of muscle and explosively released when triggered by the stimulation of sensory organs resembling hairs on the inside of the mandibles. The mandibles also permit slow and fine movements for other tasks. Trap-jaws also are seen in the following genera: Anochetus, Orectognathus, and Strumigenys,[82] plus some members of the Dacetini tribe,[83] which are viewed as examples of convergent evolution.
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A Malaysian species of ant in the Camponotus cylindricus group has enlarged mandibular glands that extend into their gaster. If combat takes a turn for the worse, a worker may perform a final act of suicidal altruism by rupturing the membrane of its gaster, causing the content of its mandibular glands to burst from the anterior region of its head, spraying a poisonous, corrosive secretion containing acetophenones and other chemicals that immobilise small insect attackers. The worker subsequently dies.[84]
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Suicidal defences by workers are also noted in a Brazilian ant, Forelius pusillus, where a small group of ants leaves the security of the nest after sealing the entrance from the outside each evening.[85]
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In addition to defence against predators, ants need to protect their colonies from pathogens. Some worker ants maintain the hygiene of the colony and their activities include undertaking or necrophory, the disposal of dead nest-mates.[86] Oleic acid has been identified as the compound released from dead ants that triggers necrophoric behaviour in Atta mexicana[87] while workers of Linepithema humile react to the absence of characteristic chemicals (dolichodial and iridomyrmecin) present on the cuticle of their living nestmates to trigger similar behaviour.[88]
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Nests may be protected from physical threats such as flooding and overheating by elaborate nest architecture.[89][90] Workers of Cataulacus muticus, an arboreal species that lives in plant hollows, respond to flooding by drinking water inside the nest, and excreting it outside.[91] Camponotus anderseni, which nests in the cavities of wood in mangrove habitats, deals with submergence under water by switching to anaerobic respiration.[92]
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Many animals can learn behaviours by imitation, but ants may be the only group apart from mammals where interactive teaching has been observed. A knowledgeable forager of Temnothorax albipennis can lead a naïve nest-mate to newly discovered food by the process of tandem running. The follower obtains knowledge through its leading tutor. The leader is acutely sensitive to the progress of the follower and slows down when the follower lags and speeds up when the follower gets too close.[93]
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Controlled experiments with colonies of Cerapachys biroi suggest that an individual may choose nest roles based on her previous experience. An entire generation of identical workers was divided into two groups whose outcome in food foraging was controlled. One group was continually rewarded with prey, while it was made certain that the other failed. As a result, members of the successful group intensified their foraging attempts while the unsuccessful group ventured out fewer and fewer times. A month later, the successful foragers continued in their role while the others had moved to specialise in brood care.[94]
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Complex nests are built by many ant species, but other species are nomadic and do not build permanent structures. Ants may form subterranean nests or build them on trees. These nests may be found in the ground, under stones or logs, inside logs, hollow stems, or even acorns. The materials used for construction include soil and plant matter,[67] and ants carefully select their nest sites; Temnothorax albipennis will avoid sites with dead ants, as these may indicate the presence of pests or disease. They are quick to abandon established nests at the first sign of threats.[95]
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The army ants of South America, such as the Eciton burchellii species, and the driver ants of Africa do not build permanent nests, but instead, alternate between nomadism and stages where the workers form a temporary nest (bivouac) from their own bodies, by holding each other together.[96]
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Weaver ant (Oecophylla spp.) workers build nests in trees by attaching leaves together, first pulling them together with bridges of workers and then inducing their larvae to produce silk as they are moved along the leaf edges. Similar forms of nest construction are seen in some species of Polyrhachis.[97]
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Formica polyctena, among other ant species, constructs nests that maintain a relatively constant interior temperature that aids in the development of larvae. The ants maintain the nest temperature by choosing the location, nest materials, controlling ventilation and maintaining the heat from solar radiation, worker activity and metabolism, and in some moist nests, microbial activity in the nest materials.[98]
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Some ant species, such as those that use natural cavities, can be opportunistic and make use of the controlled micro-climate provided inside human dwellings and other artificial structures to house their colonies and nest structures.[99][100]
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Most ants are generalist predators, scavengers, and indirect herbivores,[16] but a few have evolved specialised ways of obtaining nutrition. It is believed that many ant species that engage in indirect herbivory rely on specialized symbiosis with their gut microbes [101] to upgrade the nutritional value of the food they collect [102] and allow them to survive in nitrogen poor regions, such as rainforest canopies.[103] Leafcutter ants (Atta and Acromyrmex) feed exclusively on a fungus that grows only within their colonies. They continually collect leaves which are taken to the colony, cut into tiny pieces and placed in fungal gardens. Ergates specialise in related tasks according to their sizes. The largest ants cut stalks, smaller workers chew the leaves and the smallest tend the fungus. Leafcutter ants are sensitive enough to recognise the reaction of the fungus to different plant material, apparently detecting chemical signals from the fungus. If a particular type of leaf is found to be toxic to the fungus, the colony will no longer collect it. The ants feed on structures produced by the fungi called gongylidia. Symbiotic bacteria on the exterior surface of the ants produce antibiotics that kill bacteria introduced into the nest that may harm the fungi.[104]
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Foraging ants travel distances of up to 200 metres (700 ft) from their nest [105] and scent trails allow them to find their way back even in the dark. In hot and arid regions, day-foraging ants face death by desiccation, so the ability to find the shortest route back to the nest reduces that risk. Diurnal desert ants of the genus Cataglyphis such as the Sahara desert ant navigate by keeping track of direction as well as distance travelled. Distances travelled are measured using an internal pedometer that keeps count of the steps taken[106] and also by evaluating the movement of objects in their visual field (optical flow).[107] Directions are measured using the position of the sun.[108]
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They integrate this information to find the shortest route back to their nest.[109]
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Like all ants, they can also make use of visual landmarks when available[110] as well as olfactory and tactile cues to navigate.[111][112] Some species of ant are able to use the Earth's magnetic field for navigation.[113] The compound eyes of ants have specialised cells that detect polarised light from the Sun, which is used to determine direction.[114][115]
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These polarization detectors are sensitive in the ultraviolet region of the light spectrum.[116] In some army ant species, a group of foragers who become separated from the main column may sometimes turn back on themselves and form a circular ant mill. The workers may then run around continuously until they die of exhaustion.[117]
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The female worker ants do not have wings and reproductive females lose their wings after their mating flights in order to begin their colonies. Therefore, unlike their wasp ancestors, most ants travel by walking. Some species are capable of leaping. For example, Jerdon's jumping ant (Harpegnathos saltator) is able to jump by synchronising the action of its mid and hind pairs of legs.[118] There are several species of gliding ant including Cephalotes atratus; this may be a common trait among arboreal ants with small colonies. Ants with this ability are able to control their horizontal movement so as to catch tree trunks when they fall from atop the forest canopy.[119]
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Other species of ants can form chains to bridge gaps over water, underground, or through spaces in vegetation. Some species also form floating rafts that help them survive floods.[120] These rafts may also have a role in allowing ants to colonise islands.[121] Polyrhachis sokolova, a species of ant found in Australian mangrove swamps, can swim and live in underwater nests. Since they lack gills, they go to trapped pockets of air in the submerged nests to breathe.[122]
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Not all ants have the same kind of societies. The Australian bulldog ants are among the biggest and most basal of ants. Like virtually all ants, they are eusocial, but their social behaviour is poorly developed compared to other species. Each individual hunts alone, using her large eyes instead of chemical senses to find prey.[123]
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Some species (such as Tetramorium caespitum) attack and take over neighbouring ant colonies. Others are less expansionist, but just as aggressive; they invade colonies to steal eggs or larvae, which they either eat or raise as workers or slaves. Extreme specialists among these slave-raiding ants, such as the Amazon ants, are incapable of feeding themselves and need captured workers to survive.[124] Captured workers of enslaved Temnothorax species have evolved a counter strategy, destroying just the female pupae of the slave-making Temnothorax americanus, but sparing the males (who don't take part in slave-raiding as adults).[125]
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Ants identify kin and nestmates through their scent, which comes from hydrocarbon-laced secretions that coat their exoskeletons. If an ant is separated from its original colony, it will eventually lose the colony scent. Any ant that enters a colony without a matching scent will be attacked.[126] Also, the reason why two separate colonies of ants will attack each other even if they are of the same species is because the genes responsible for pheromone production are different between them. The Argentine ant, however, does not have this characteristic, due to lack of genetic diversity, and has become a global pest because of it.
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Parasitic ant species enter the colonies of host ants and establish themselves as social parasites; species such as Strumigenys xenos are entirely parasitic and do not have workers, but instead, rely on the food gathered by their Strumigenys perplexa hosts.[127][128] This form of parasitism is seen across many ant genera, but the parasitic ant is usually a species that is closely related to its host. A variety of methods are employed to enter the nest of the host ant. A parasitic queen may enter the host nest before the first brood has hatched, establishing herself prior to development of a colony scent. Other species use pheromones to confuse the host ants or to trick them into carrying the parasitic queen into the nest. Some simply fight their way into the nest.[129]
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A conflict between the sexes of a species is seen in some species of ants with these reproducers apparently competing to produce offspring that are as closely related to them as possible. The most extreme form involves the production of clonal offspring. An extreme of sexual conflict is seen in Wasmannia auropunctata, where the queens produce diploid daughters by thelytokous parthenogenesis and males produce clones by a process whereby a diploid egg loses its maternal contribution to produce haploid males who are clones of the father.[130]
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Ants form symbiotic associations with a range of species, including other ant species, other insects, plants, and fungi. They also are preyed on by many animals and even certain fungi. Some arthropod species spend part of their lives within ant nests, either preying on ants, their larvae, and eggs, consuming the food stores of the ants, or avoiding predators. These inquilines may bear a close resemblance to ants. The nature of this ant mimicry (myrmecomorphy) varies, with some cases involving Batesian mimicry, where the mimic reduces the risk of predation. Others show Wasmannian mimicry, a form of mimicry seen only in inquilines.[131][132]
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Aphids and other hemipteran insects secrete a sweet liquid called honeydew, when they feed on plant sap. The sugars in honeydew are a high-energy food source, which many ant species collect.[133] In some cases, the aphids secrete the honeydew in response to ants tapping them with their antennae. The ants in turn keep predators away from the aphids and will move them from one feeding location to another. When migrating to a new area, many colonies will take the aphids with them, to ensure a continued supply of honeydew. Ants also tend mealybugs to harvest their honeydew. Mealybugs may become a serious pest of pineapples if ants are present to protect mealybugs from their natural enemies.[134]
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Myrmecophilous (ant-loving) caterpillars of the butterfly family Lycaenidae (e.g., blues, coppers, or hairstreaks) are herded by the ants, led to feeding areas in the daytime, and brought inside the ants' nest at night. The caterpillars have a gland which secretes honeydew when the ants massage them. Some caterpillars produce vibrations and sounds that are perceived by the ants.[135] A similar adaptation can be seen in Grizzled skipper butterflies that emit vibrations by expanding their wings in order to communicate with ants, which are natural predators of these butterflies.[136] Other caterpillars have evolved from ant-loving to ant-eating: these myrmecophagous caterpillars secrete a pheromone that makes the ants act as if the caterpillar is one of their own larvae. The caterpillar is then taken into the ant nest where it feeds on the ant larvae.[137] A number of specialized bacterial have been found as endosymbionts in ant guts. Some of the dominant bacteria belong to the order Rhizobiales whose members are known for being nitrogen-fixing symbionts in legumes but the species found in ant lack the ability to fix nitrogen.[138][139] Fungus-growing ants that make up the tribe Attini, including leafcutter ants, cultivate certain species of fungus in the genera Leucoagaricus or Leucocoprinus of the family Agaricaceae. In this ant-fungus mutualism, both species depend on each other for survival. The ant Allomerus decemarticulatus has evolved a three-way association with the host plant, Hirtella physophora (Chrysobalanaceae), and a sticky fungus which is used to trap their insect prey.[140]
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+
Lemon ants make devil's gardens by killing surrounding plants with their stings and leaving a pure patch of lemon ant trees, (Duroia hirsuta). This modification of the forest provides the ants with more nesting sites inside the stems of the Duroia trees.[141] Although some ants obtain nectar from flowers, pollination by ants is somewhat rare, one example being of the pollination of the orchid Leporella fimbriata which induces male Myrmecia urens to pseudocopulate with the flowers, transferring pollen in the process.[142] One theory that has been proposed for the rarity of pollination is that the secretions of the metapleural gland inactivate and reduce the viability of pollen.[143][144] Some plants have special nectar exuding structures, extrafloral nectaries, that provide food for ants, which in turn protect the plant from more damaging herbivorous insects.[145] Species such as the bullhorn acacia (Acacia cornigera) in Central America have hollow thorns that house colonies of stinging ants (Pseudomyrmex ferruginea) who defend the tree against insects, browsing mammals, and epiphytic vines. Isotopic labelling studies suggest that plants also obtain nitrogen from the ants.[146] In return, the ants obtain food from protein- and lipid-rich Beltian bodies. In Fiji Philidris nagasau (Dolichoderinae) are known to selectively grow species of epiphytic Squamellaria (Rubiaceae) which produce large domatia inside which the ant colonies nest. The ants plant the seeds and the domatia of young seedling are immediately occupied and the ant faeces in them contribute to rapid growth.[147] Similar dispersal associations are found with other dolichoderines in the region as well.[148] Another example of this type of ectosymbiosis comes from the Macaranga tree, which has stems adapted to house colonies of Crematogaster ants.[149]
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+
Many plant species have seeds that are adapted for dispersal by ants.[150] Seed dispersal by ants or myrmecochory is widespread, and new estimates suggest that nearly 9% of all plant species may have such ant associations.[151][150] Often, seed-dispersing ants perform directed dispersal, depositing the seeds in locations that increase the likelihood of seed survival to reproduction.[152] Some plants in arid, fire-prone systems are particularly dependent on ants for their survival and dispersal as the seeds are transported to safety below the ground.[153] Many ant-dispersed seeds have special external structures, elaiosomes, that are sought after by ants as food.[154]
|
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A convergence, possibly a form of mimicry, is seen in the eggs of stick insects. They have an edible elaiosome-like structure and are taken into the ant nest where the young hatch.[155]
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Most ants are predatory and some prey on and obtain food from other social insects including other ants. Some species specialise in preying on termites (Megaponera and Termitopone) while a few Cerapachyinae prey on other ants.[105] Some termites, including Nasutitermes corniger, form associations with certain ant species to keep away predatory ant species.[156] The tropical wasp Mischocyttarus drewseni coats the pedicel of its nest with an ant-repellent chemical.[157] It is suggested that many tropical wasps may build their nests in trees and cover them to protect themselves from ants. Other wasps, such as A. multipicta, defend against ants by blasting them off the nest with bursts of wing buzzing.[158] Stingless bees (Trigona and Melipona) use chemical defences against ants.[105]
|
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Flies in the Old World genus Bengalia (Calliphoridae) prey on ants and are kleptoparasites, snatching prey or brood from the mandibles of adult ants.[159] Wingless and legless females of the Malaysian phorid fly (Vestigipoda myrmolarvoidea) live in the nests of ants of the genus Aenictus and are cared for by the ants.[159]
|
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Fungi in the genera Cordyceps and Ophiocordyceps infect ants. Ants react to their infection by climbing up plants and sinking their mandibles into plant tissue. The fungus kills the ants, grows on their remains, and produces a fruiting body. It appears that the fungus alters the behaviour of the ant to help disperse its spores [160] in a microhabitat that best suits the fungus.[161] Strepsipteran parasites also manipulate their ant host to climb grass stems, to help the parasite find mates.[162]
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A nematode (Myrmeconema neotropicum) that infects canopy ants (Cephalotes atratus) causes the black-coloured gasters of workers to turn red. The parasite also alters the behaviour of the ant, causing them to carry their gasters high. The conspicuous red gasters are mistaken by birds for ripe fruits, such as Hyeronima alchorneoides, and eaten. The droppings of the bird are collected by other ants and fed to their young, leading to further spread of the nematode.[163]
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South American poison dart frogs in the genus Dendrobates feed mainly on ants, and the toxins in their skin may come from the ants.[164]
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Army ants forage in a wide roving column, attacking any animals in that path that are unable to escape. In Central and South America, Eciton burchellii is the swarming ant most commonly attended by "ant-following" birds such as antbirds and woodcreepers.[165][166] This behaviour was once considered mutualistic, but later studies found the birds to be parasitic. Direct kleptoparasitism (birds stealing food from the ants' grasp) is rare and has been noted in Inca doves which pick seeds at nest entrances as they are being transported by species of Pogonomyrmex.[167] Birds that follow ants eat many prey insects and thus decrease the foraging success of ants.[168] Birds indulge in a peculiar behaviour called anting that, as yet, is not fully understood. Here birds rest on ant nests, or pick and drop ants onto their wings and feathers; this may be a means to remove ectoparasites from the birds.
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Anteaters, aardvarks, pangolins, echidnas and numbats have special adaptations for living on a diet of ants. These adaptations include long, sticky tongues to capture ants and strong claws to break into ant nests. Brown bears (Ursus arctos) have been found to feed on ants. About 12%, 16%, and 4% of their faecal volume in spring, summer, and autumn, respectively, is composed of ants.[169]
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Ants perform many ecological roles that are beneficial to humans, including the suppression of pest populations and aeration of the soil. The use of weaver ants in citrus cultivation in southern China is considered one of the oldest known applications of biological control.[170] On the other hand, ants may become nuisances when they invade buildings, or cause economic losses.
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In some parts of the world (mainly Africa and South America), large ants, especially army ants, are used as surgical sutures. The wound is pressed together and ants are applied along it. The ant seizes the edges of the wound in its mandibles and locks in place. The body is then cut off and the head and mandibles remain in place to close the wound.[171][172][173] The large heads of the dinergates (soldiers) of the leafcutting ant Atta cephalotes are also used by native surgeons in closing wounds.[174]
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Some ants have toxic venom and are of medical importance. The species include Paraponera clavata (tocandira) and Dinoponera spp. (false tocandiras) of South America [175] and the Myrmecia ants of Australia.[176]
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In South Africa, ants are used to help harvest the seeds of rooibos (Aspalathus linearis), a plant used to make a herbal tea. The plant disperses its seeds widely, making manual collection difficult. Black ants collect and store these and other seeds in their nest, where humans can gather them en masse. Up to half a pound (200 g) of seeds may be collected from one ant-heap.[177][178]
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Although most ants survive attempts by humans to eradicate them, a few are highly endangered. These tend to be island species that have evolved specialized traits and risk being displaced by introduced ant species. Examples include the critically endangered Sri Lankan relict ant (Aneuretus simoni) and Adetomyrma venatrix of Madagascar.[179]
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It has been estimated by E.O. Wilson that the total number of individual ants alive in the world at any one time is between one and ten quadrillion (short scale) (i.e., between 1015 and 1016). According to this estimate, the total biomass of all the ants in the world is approximately equal to the total biomass of the entire human race.[180] Also, according to this estimate, there are approximately 1 million ants for every human on Earth.[181]
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Ants and their larvae are eaten in different parts of the world. The eggs of two species of ants are used in Mexican escamoles. They are considered a form of insect caviar and can sell for as much as US$40 per pound ($90/kg) because they are seasonal and hard to find. In the Colombian department of Santander, hormigas culonas (roughly interpreted as "large-bottomed ants") Atta laevigata are toasted alive and eaten.[182]
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In areas of India, and throughout Burma and Thailand, a paste of the green weaver ant (Oecophylla smaragdina) is served as a condiment with curry.[183] Weaver ant eggs and larvae, as well as the ants, may be used in a Thai salad, yam (Thai: ยำ), in a dish called yam khai mot daeng (Thai: ยำไข่มดแดง) or red ant egg salad, a dish that comes from the Issan or north-eastern region of Thailand. Saville-Kent, in the Naturalist in Australia wrote "Beauty, in the case of the green ant, is more than skin-deep. Their attractive, almost sweetmeat-like translucency possibly invited the first essays at their consumption by the human species". Mashed up in water, after the manner of lemon squash, "these ants form a pleasant acid drink which is held in high favor by the natives of North Queensland, and is even appreciated by many European palates".[184]
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In his First Summer in the Sierra, John Muir notes that the Digger Indians of California ate the tickling, acid gasters of the large jet-black carpenter ants. The Mexican Indians eat the replete workers, or living honey-pots, of the honey ant (Myrmecocystus).[184]
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Some ant species are considered as pests, primarily those that occur in human habitations, where their presence is often problematic. For example, the presence of ants would be undesirable in sterile places such as hospitals or kitchens. Some species or genera commonly categorized as pests include the Argentine ant, pavement ant, yellow crazy ant, banded sugar ant, pharaoh ant, red ant, carpenter ant, odorous house ant, red imported fire ant, and European fire ant. Some ants will raid stored food, some will seek water sources, others may damage indoor structures, some may damage agricultural crops directly (or by aiding sucking pests). Some will sting or bite.[185] The adaptive nature of ant colonies make it nearly impossible to eliminate entire colonies and most pest management practices aim to control local populations and tend to be temporary solutions. Ant populations are managed by a combination of approaches that make use of chemical, biological, and physical methods. Chemical methods include the use of insecticidal bait which is gathered by ants as food and brought back to the nest where the poison is inadvertently spread to other colony members through trophallaxis. Management is based on the species and techniques may vary according to the location and circumstance.[185]
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Observed by humans since the dawn of history, the behaviour of ants has been documented and the subject of early writings and fables passed from one century to another. Those using scientific methods, myrmecologists, study ants in the laboratory and in their natural conditions. Their complex and variable social structures have made ants ideal model organisms. Ultraviolet vision was first discovered in ants by Sir John Lubbock in 1881.[186] Studies on ants have tested hypotheses in ecology and sociobiology, and have been particularly important in examining the predictions of theories of kin selection and evolutionarily stable strategies.[187] Ant colonies may be studied by rearing or temporarily maintaining them in formicaria, specially constructed glass framed enclosures.[188] Individuals may be tracked for study by marking them with dots of colours.[189]
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The successful techniques used by ant colonies have been studied in computer science and robotics to produce distributed and fault-tolerant systems for solving problems, for example Ant colony optimization and Ant robotics. This area of biomimetics has led to studies of ant locomotion, search engines that make use of "foraging trails", fault-tolerant storage, and networking algorithms.[190]
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From the late 1950s through the late 1970s, ant farms were popular educational children's toys in the United States. Some later commercial versions use transparent gel instead of soil, allowing greater visibility at the cost of stressing the ants with unnatural light.[191]
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Anthropomorphised ants have often been used in fables and children's stories to represent industriousness and cooperative effort. They also are mentioned in religious texts.[192][193] In the Book of Proverbs in the Bible, ants are held up as a good example for humans for their hard work and cooperation.[194] Aesop did the same in his fable The Ant and the Grasshopper. In the Quran, Sulayman is said to have heard and understood an ant warning other ants to return home to avoid being accidentally crushed by Sulayman and his marching army.[Quran 27:18][195] In parts of Africa, ants are considered to be the messengers of the deities. Some Native American mythology, such as the Hopi mythology, considers ants as the very first animals. Ant bites are often said to have curative properties. The sting of some species of Pseudomyrmex is claimed to give fever relief.[196] Ant bites are used in the initiation ceremonies of some Amazon Indian cultures as a test of endurance.[197][198]
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Ant society has always fascinated humans and has been written about both humorously and seriously. Mark Twain wrote about ants in his 1880 book A Tramp Abroad.[199] Some modern authors have used the example of the ants to comment on the relationship between society and the individual. Examples are Robert Frost in his poem "Departmental" and T. H. White in his fantasy novel The Once and Future King. The plot in French entomologist and writer Bernard Werber's Les Fourmis science-fiction trilogy is divided between the worlds of ants and humans; ants and their behaviour is described using contemporary scientific knowledge. H.G. Wells wrote about intelligent ants destroying human settlements in Brazil and threatening human civilization in his 1905 science-fiction short story, The Empire of the Ants. In more recent times, animated cartoons and 3-D animated films featuring ants have been produced including Antz, A Bug's Life, The Ant Bully, The Ant and the Aardvark, Ferdy the Ant and Atom Ant. Renowned myrmecologist E. O. Wilson wrote a short story, "Trailhead" in 2010 for The New Yorker magazine, which describes the life and death of an ant-queen and the rise and fall of her colony, from an ants' point of view.[200] The French neuroanatomist, psychiatrist and eugenicist Auguste Forel believed that ant societies were models for human society. He published a five volume work from 1921 to 1923 that examined ant biology and society.[201]
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In the early 1990s, the video game SimAnt, which simulated an ant colony, won the 1992 Codie award for "Best Simulation Program".[202]
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Ants also are quite popular inspiration for many science-fiction insectoids, such as the Formics of Ender's Game, the Bugs of Starship Troopers, the giant ants in the films Them! and Empire of the Ants, Marvel Comics' super hero Ant-Man, and ants mutated into super-intelligence in Phase IV. In computer strategy games, ant-based species often benefit from increased production rates due to their single-minded focus, such as the Klackons in the Master of Orion series of games or the ChCht in Deadlock II. These characters are often credited with a hive mind, a common misconception about ant colonies.[203]
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1 |
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2 |
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3 |
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4 |
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5 |
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Martialinae
|
6 |
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7 |
+
Leptanillinae
|
8 |
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|
9 |
+
Amblyoponinae
|
10 |
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|
11 |
+
Paraponerinae
|
12 |
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|
13 |
+
Agroecomyrmecinae
|
14 |
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15 |
+
Ponerinae
|
16 |
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17 |
+
Proceratiinae
|
18 |
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|
19 |
+
Ecitoninae‡
|
20 |
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|
21 |
+
Aenictinae‡
|
22 |
+
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23 |
+
Dorylini‡
|
24 |
+
|
25 |
+
Aenictogitoninae‡
|
26 |
+
|
27 |
+
Cerapachyinae‡*
|
28 |
+
|
29 |
+
Leptanilloidinae‡
|
30 |
+
|
31 |
+
Dolichoderinae
|
32 |
+
|
33 |
+
Aneuretinae
|
34 |
+
|
35 |
+
Pseudomyrmecinae
|
36 |
+
|
37 |
+
Myrmeciinae
|
38 |
+
|
39 |
+
Ectatomminae
|
40 |
+
|
41 |
+
Heteroponerinae
|
42 |
+
|
43 |
+
Myrmicinae
|
44 |
+
|
45 |
+
Formicinae
|
46 |
+
|
47 |
+
A phylogeny of the extant ant subfamilies.[2][3]
|
48 |
+
*Cerapachyinae is paraphyletic
|
49 |
+
‡ The previous dorylomorph subfamilies were synonymized under Dorylinae by Brady et al. in 2014[4]
|
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+
|
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+
Ants are eusocial insects of the family Formicidae and, along with the related wasps and bees, belong to the order Hymenoptera. Ants appear in the fossil record across the globe in considerable diversity during the latest Early Cretaceous and early Late Cretaceous, suggesting an earlier origin. Ants evolved from vespoid wasp ancestors in the Cretaceous period, and diversified after the rise of flowering plants. More than 12,500 of an estimated total of 22,000 species have been classified. They are easily identified by their elbowed antennae and the distinctive node-like structure that forms their slender waists.
|
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+
|
53 |
+
Ants form colonies that range in size from a few dozen predatory individuals living in small natural cavities to highly organised colonies that may occupy large territories and consist of millions of individuals. Larger colonies consist of various castes of sterile, wingless females, most of which are workers (ergates), as well as soldiers (dinergates) and other specialised groups. Nearly all ant colonies also have some fertile males called "drones" (aner) and one or more fertile females called "queens" (gynes). The colonies are described as superorganisms because the ants appear to operate as a unified entity, collectively working together to support the colony.
|
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+
|
55 |
+
Ants have colonised almost every landmass on Earth. The only places lacking indigenous ants are Antarctica and a few remote or inhospitable islands. Ants thrive in most ecosystems and may form 15–25% of the terrestrial animal biomass. Their success in so many environments has been attributed to their social organisation and their ability to modify habitats, tap resources, and defend themselves. Their long co-evolution with other species has led to mimetic, commensal, parasitic, and mutualistic relationships.
|
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+
|
57 |
+
Ant societies have division of labour, communication between individuals, and an ability to solve complex problems. These parallels with human societies have long been an inspiration and subject of study. Many human cultures make use of ants in cuisine, medication, and rituals. Some species are valued in their role as biological pest control agents. Their ability to exploit resources may bring ants into conflict with humans, however, as they can damage crops and invade buildings. Some species, such as the red imported fire ant (Solenopsis invicta), are regarded as invasive species, establishing themselves in areas where they have been introduced accidentally.
|
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|
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+
The word ant and its chiefly dialectal form emmet[5] come from ante, emete of Middle English, which come from ǣmette of Old English, and these are all related to the dialectal Dutch emt and the Old High German āmeiza, from which comes the modern German Ameise. All of these words come from West Germanic *ǣmaitjōn, and the original meaning of the word was "the biter" (from Proto-Germanic *ai-, "off, away" + *mait- "cut").[6][7] The family name Formicidae is derived from the Latin formīca ("ant")[8] from which the words in other Romance languages, such as the Portuguese formiga, Italian formica, Spanish hormiga, Romanian furnică, and French fourmi are derived. It has been hypothesised that a Proto-Indo-European word *morwi- was used, cf. Sanskrit vamrah, Latin formīca, Greek μύρμηξ mýrmēx, Old Church Slavonic mraviji, Old Irish moirb, Old Norse maurr, Dutch mier.[9]
|
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+
|
61 |
+
Chrysidoidea
|
62 |
+
|
63 |
+
Vespidae
|
64 |
+
|
65 |
+
Rhopalosomatidae
|
66 |
+
|
67 |
+
Pompilidae
|
68 |
+
|
69 |
+
Tiphiidae
|
70 |
+
|
71 |
+
Scolioidea
|
72 |
+
|
73 |
+
Apoidea
|
74 |
+
|
75 |
+
Formicidae
|
76 |
+
|
77 |
+
The family Formicidae belongs to the order Hymenoptera, which also includes sawflies, bees, and wasps. Ants evolved from a lineage within the stinging wasps, and a 2013 study suggests that they are a sister group of the Apoidea.[10] In 1966, E. O. Wilson and his colleagues identified the fossil remains of an ant (Sphecomyrma) that lived in the Cretaceous period. The specimen, trapped in amber dating back to around 92 million years ago, has features found in some wasps, but not found in modern ants.[11] Sphecomyrma was possibly a ground forager, while Haidomyrmex and Haidomyrmodes, related genera in subfamily Sphecomyrminae, are reconstructed as active arboreal predators.[12] Older ants in the genus Sphecomyrmodes have been found in 99 million year-old amber from Myanmar.[13][14] A 2006 study suggested that ants arose tens of millions of years earlier than previously thought, up to 168 million years ago.[1] After the rise of flowering plants about 100 million years ago they diversified and assumed ecological dominance around 60 million years ago.[15][1][16][17] Some groups, such as the Leptanillinae and Martialinae, are suggested to have diversified from early primitive ants that were likely to have been predators underneath the surface of the soil.[3][18]
|
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During the Cretaceous period, a few species of primitive ants ranged widely on the Laurasian supercontinent (the Northern Hemisphere). They were scarce in comparison to the populations of other insects, representing only about 1% of the entire insect population. Ants became dominant after adaptive radiation at the beginning of the Paleogene period. By the Oligocene and Miocene, ants had come to represent 20–40% of all insects found in major fossil deposits. Of the species that lived in the Eocene epoch, around one in 10 genera survive to the present. Genera surviving today comprise 56% of the genera in Baltic amber fossils (early Oligocene), and 92% of the genera in Dominican amber fossils (apparently early Miocene).[15][19]
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Termites live in colonies and are sometimes called ‘white ants’, but termites are not ants. They are the sub-order Isoptera, and together with cockroaches they form the order Blattodea. Blattodeans are related to mantids, crickets, and other winged insects that do not undergo full metamorphosis. Like ants, termites are eusocial, with sterile workers, but they differ greatly in the genetics of reproduction. The similarity of their social structure to that of ants is attributed to convergent evolution.[20] Velvet ants look like large ants, but are wingless female wasps.[21][22]
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Ants are found on all continents except Antarctica, and only a few large islands, such as Greenland, Iceland, parts of Polynesia and the Hawaiian Islands lack native ant species.[24][25] Ants occupy a wide range of ecological niches and exploit many different food resources as direct or indirect herbivores, predators and scavengers. Most ant species are omnivorous generalists, but a few are specialist feeders. Their ecological dominance is demonstrated by their biomass: ants are estimated to contribute 15–20 % (on average and nearly 25% in the tropics) of terrestrial animal biomass, exceeding that of the vertebrates.[26]
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Ants range in size from 0.75 to 52 millimetres (0.030–2.0 in),[27][28] the largest species being the fossil Titanomyrma giganteum, the queen of which was 6 centimetres (2.4 in) long with a wingspan of 15 centimetres (5.9 in).[29] Ants vary in colour; most ants are red or black, but a few species are green and some tropical species have a metallic lustre. More than 12,000 species are currently known (with upper estimates of the potential existence of about 22,000) (see the article List of ant genera), with the greatest diversity in the tropics. Taxonomic studies continue to resolve the classification and systematics of ants. Online databases of ant species, including AntBase and the Hymenoptera Name Server, help to keep track of the known and newly described species.[30] The relative ease with which ants may be sampled and studied in ecosystems has made them useful as indicator species in biodiversity studies.[31][32]
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Ants are distinct in their morphology from other insects in having elbowed antennae, metapleural glands, and a strong constriction of their second abdominal segment into a node-like petiole. The head, mesosoma, and metasoma are the three distinct body segments (formally tagmata). The petiole forms a narrow waist between their mesosoma (thorax plus the first abdominal segment, which is fused to it) and gaster (abdomen less the abdominal segments in the petiole). The petiole may be formed by one or two nodes (the second alone, or the second and third abdominal segments).[33]
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Like other insects, ants have an exoskeleton, an external covering that provides a protective casing around the body and a point of attachment for muscles, in contrast to the internal skeletons of humans and other vertebrates. Insects do not have lungs; oxygen and other gases, such as carbon dioxide, pass through their exoskeleton via tiny valves called spiracles. Insects also lack closed blood vessels; instead, they have a long, thin, perforated tube along the top of the body (called the "dorsal aorta") that functions like a heart, and pumps haemolymph toward the head, thus driving the circulation of the internal fluids. The nervous system consists of a ventral nerve cord that runs the length of the body, with several ganglia and branches along the way reaching into the extremities of the appendages.[34]
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An ant's head contains many sensory organs. Like most insects, ants have compound eyes made from numerous tiny lenses attached together. Ant eyes are good for acute movement detection, but do not offer a high resolution image. They also have three small ocelli (simple eyes) on the top of the head that detect light levels and polarization.[35] Compared to vertebrates, ants tend to have blurrier eyesight, particularly in smaller species,[36] and a few subterranean taxa are completely blind.[2] However, some ants, such as Australia's bulldog ant, have excellent vision and are capable of discriminating the distance and size of objects moving nearly a metre away.[37]
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Two antennae ("feelers") are attached to the head; these organs detect chemicals, air currents, and vibrations; they also are used to transmit and receive signals through touch. The head has two strong jaws, the mandibles, used to carry food, manipulate objects, construct nests, and for defence.[34] In some species, a small pocket (infrabuccal chamber) inside the mouth stores food, so it may be passed to other ants or their larvae.[38]
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Both the legs and wings of the ant are attached to the mesosoma ("thorax"). The legs terminate in a hooked claw which allows them to hook on and climb surfaces.[39] Only reproductive ants, queens, and males, have wings. Queens shed their wings after the nuptial flight, leaving visible stubs, a distinguishing feature of queens. In a few species, wingless queens (ergatoids) and males occur.[34]
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The metasoma (the "abdomen") of the ant houses important internal organs, including those of the reproductive, respiratory (tracheae), and excretory systems. Workers of many species have their egg-laying structures modified into stings that are used for subduing prey and defending their nests.[34]
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In the colonies of a few ant species, there are physical castes—workers in distinct size-classes, called minor, median, and major ergates. Often, the larger ants have disproportionately larger heads, and correspondingly stronger mandibles. These are known as macrergates while smaller workers are known as micrergates.[40] Although formally known as dinergates, such individuals are sometimes called "soldier" ants because their stronger mandibles make them more effective in fighting, although they still are workers and their "duties" typically do not vary greatly from the minor or median workers. In a few species, the median workers are absent, creating a sharp divide between the minors and majors.[41] Weaver ants, for example, have a distinct bimodal size distribution.[42][43] Some other species show continuous variation in the size of workers. The smallest and largest workers in Pheidologeton diversus show nearly a 500-fold difference in their dry-weights.[44]
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Workers cannot mate; however, because of the haplodiploid sex-determination system in ants, workers of a number of species can lay unfertilised eggs that become fully fertile, haploid males. The role of workers may change with their age and in some species, such as honeypot ants, young workers are fed until their gasters are distended, and act as living food storage vessels. These food storage workers are called repletes.[45] For instance, these replete workers develop in the North American honeypot ant Myrmecocystus mexicanus. Usually the largest workers in the colony develop into repletes; and, if repletes are removed from the colony, other workers become repletes, demonstrating the flexibility of this particular polymorphism.[46] This polymorphism in morphology and behaviour of workers initially was thought to be determined by environmental factors such as nutrition and hormones that led to different developmental paths; however, genetic differences between worker castes have been noted in Acromyrmex sp.[47] These polymorphisms are caused by relatively small genetic changes; differences in a single gene of Solenopsis invicta can decide whether the colony will have single or multiple queens.[48] The Australian jack jumper ant (Myrmecia pilosula) has only a single pair of chromosomes (with the males having just one chromosome as they are haploid), the lowest number known for any animal, making it an interesting subject for studies in the genetics and developmental biology of social insects.[49][50]
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The life of an ant starts from an egg; if the egg is fertilised, the progeny will be female diploid, if not, it will be male haploid. Ants develop by complete metamorphosis with the larva stages passing through a pupal stage before emerging as an adult. The larva is largely immobile and is fed and cared for by workers. Food is given to the larvae by trophallaxis, a process in which an ant regurgitates liquid food held in its crop. This is also how adults share food, stored in the "social stomach". Larvae, especially in the later stages, may also be provided solid food, such as trophic eggs, pieces of prey, and seeds brought by workers.[51]
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The larvae grow through a series of four or five moults and enter the pupal stage. The pupa has the appendages free and not fused to the body as in a butterfly pupa.[52] The differentiation into queens and workers (which are both female), and different castes of workers, is influenced in some species by the nutrition the larvae obtain. Genetic influences and the control of gene expression by the developmental environment are complex and the determination of caste continues to be a subject of research.[53] Winged male ants, called drones, emerge from pupae along with the usually winged breeding females. Some species, such as army ants, have wingless queens. Larvae and pupae need to be kept at fairly constant temperatures to ensure proper development, and so often are moved around among the various brood chambers within the colony.[54]
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A new ergate spends the first few days of its adult life caring for the queen and young. She then graduates to digging and other nest work, and later to defending the nest and foraging. These changes are sometimes fairly sudden, and define what are called temporal castes. An explanation for the sequence is suggested by the high casualties involved in foraging, making it an acceptable risk only for ants who are older and are likely to die soon of natural causes.[55][56]
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Ant colonies can be long-lived. The queens can live for up to 30 years, and workers live from 1 to 3 years. Males, however, are more transitory, being quite short-lived and surviving for only a few weeks.[57] Ant queens are estimated to live 100 times as long as solitary insects of a similar size.[58]
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Ants are active all year long in the tropics, but, in cooler regions, they survive the winter in a state of dormancy known as hibernation. The forms of inactivity are varied and some temperate species have larvae going into the inactive state (diapause), while in others, the adults alone pass the winter in a state of reduced activity.[59]
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A wide range of reproductive strategies have been noted in ant species. Females of many species are known to be capable of reproducing asexually through thelytokous parthenogenesis.[60] Secretions from the male accessory glands in some species can plug the female genital opening and prevent females from re-mating.[61] Most ant species have a system in which only the queen and breeding females have the ability to mate. Contrary to popular belief, some ant nests have multiple queens, while others may exist without queens. Workers with the ability to reproduce are called "gamergates" and colonies that lack queens are then called gamergate colonies; colonies with queens are said to be queen-right.[62]
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Drones can also mate with existing queens by entering a foreign colony. When the drone is initially attacked by the workers, it releases a mating pheromone. If recognized as a mate, it will be carried to the queen to mate.[63] Males may also patrol the nest and fight others by grabbing them with their mandibles, piercing their exoskeleton and then marking them with a pheromone. The marked male is interpreted as an invader by worker ants and is killed.[64]
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Most ants are univoltine, producing a new generation each year.[65] During the species-specific breeding period, winged females and winged males, known to entomologists as alates, leave the colony in what is called a nuptial flight. The nuptial flight usually takes place in the late spring or early summer when the weather is hot and humid. Heat makes flying easier and freshly fallen rain makes the ground softer for mated queens to dig nests.[66] Males typically take flight before the females. Males then use visual cues to find a common mating ground, for example, a landmark such as a pine tree to which other males in the area converge. Males secrete a mating pheromone that females follow. Males will mount females in the air, but the actual mating process usually takes place on the ground. Females of some species mate with just one male but in others they may mate with as many as ten or more different males, storing the sperm in their spermathecae.[67]
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Mated females then seek a suitable place to begin a colony. There, they break off their wings and begin to lay and care for eggs. The females can selectively fertilise future eggs with the sperm stored to produce diploid workers or lay unfertilized haploid eggs to produce drones. The first workers to hatch are known as nanitics,[68] and are weaker and smaller than later workers, but they begin to serve the colony immediately. They enlarge the nest, forage for food, and care for the other eggs. Species that have multiple queens may have a queen leaving the nest along with some workers to found a colony at a new site,[67] a process akin to swarming in honeybees.
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Ants communicate with each other using pheromones, sounds, and touch.[69] The use of pheromones as chemical signals is more developed in ants, such as the red harvester ant, than in other hymenopteran groups. Like other insects, ants perceive smells with their long, thin, and mobile antennae. The paired antennae provide information about the direction and intensity of scents. Since most ants live on the ground, they use the soil surface to leave pheromone trails that may be followed by other ants. In species that forage in groups, a forager that finds food marks a trail on the way back to the colony; this trail is followed by other ants, these ants then reinforce the trail when they head back with food to the colony. When the food source is exhausted, no new trails are marked by returning ants and the scent slowly dissipates. This behaviour helps ants deal with changes in their environment. For instance, when an established path to a food source is blocked by an obstacle, the foragers leave the path to explore new routes. If an ant is successful, it leaves a new trail marking the shortest route on its return. Successful trails are followed by more ants, reinforcing better routes and gradually identifying the best path.[70]
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Ants use pheromones for more than just making trails. A crushed ant emits an alarm pheromone that sends nearby ants into an attack frenzy and attracts more ants from farther away. Several ant species even use "propaganda pheromones" to confuse enemy ants and make them fight among themselves.[71] Pheromones are produced by a wide range of structures including Dufour's glands, poison glands and glands on the hindgut, pygidium, rectum, sternum, and hind tibia.[58] Pheromones also are exchanged, mixed with food, and passed by trophallaxis, transferring information within the colony.[72] This allows other ants to detect what task group (e.g., foraging or nest maintenance) other colony members belong to.[73] In ant species with queen castes, when the dominant queen stops producing a specific pheromone, workers begin to raise new queens in the colony.[74]
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Some ants produce sounds by stridulation, using the gaster segments and their mandibles. Sounds may be used to communicate with colony members or with other species.[75][76]
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Ants attack and defend themselves by biting and, in many species, by stinging, often injecting or spraying chemicals, such as formic acid in the case of formicine ants, alkaloids and piperidines in fire ants, and a variety of protein components in other ants. Bullet ants (Paraponera), located in Central and South America, are considered to have the most painful sting of any insect, although it is usually not fatal to humans. This sting is given the highest rating on the Schmidt sting pain index.[77]
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The sting of jack jumper ants can be fatal,[78] and an antivenom has been developed for it.[79]
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Fire ants, Solenopsis spp., are unique in having a venom sac containing piperidine alkaloids.[80] Their stings are painful and can be dangerous to hypersensitive people.[81]
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Trap-jaw ants of the genus Odontomachus are equipped with mandibles called trap-jaws, which snap shut faster than any other predatory appendages within the animal kingdom.[82] One study of Odontomachus bauri recorded peak speeds of between 126 and 230 km/h (78 and 143 mph), with the jaws closing within 130 microseconds on average.
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The ants were also observed to use their jaws as a catapult to eject intruders or fling themselves backward to escape a threat.[82] Before striking, the ant opens its mandibles extremely widely and locks them in this position by an internal mechanism. Energy is stored in a thick band of muscle and explosively released when triggered by the stimulation of sensory organs resembling hairs on the inside of the mandibles. The mandibles also permit slow and fine movements for other tasks. Trap-jaws also are seen in the following genera: Anochetus, Orectognathus, and Strumigenys,[82] plus some members of the Dacetini tribe,[83] which are viewed as examples of convergent evolution.
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A Malaysian species of ant in the Camponotus cylindricus group has enlarged mandibular glands that extend into their gaster. If combat takes a turn for the worse, a worker may perform a final act of suicidal altruism by rupturing the membrane of its gaster, causing the content of its mandibular glands to burst from the anterior region of its head, spraying a poisonous, corrosive secretion containing acetophenones and other chemicals that immobilise small insect attackers. The worker subsequently dies.[84]
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Suicidal defences by workers are also noted in a Brazilian ant, Forelius pusillus, where a small group of ants leaves the security of the nest after sealing the entrance from the outside each evening.[85]
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In addition to defence against predators, ants need to protect their colonies from pathogens. Some worker ants maintain the hygiene of the colony and their activities include undertaking or necrophory, the disposal of dead nest-mates.[86] Oleic acid has been identified as the compound released from dead ants that triggers necrophoric behaviour in Atta mexicana[87] while workers of Linepithema humile react to the absence of characteristic chemicals (dolichodial and iridomyrmecin) present on the cuticle of their living nestmates to trigger similar behaviour.[88]
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Nests may be protected from physical threats such as flooding and overheating by elaborate nest architecture.[89][90] Workers of Cataulacus muticus, an arboreal species that lives in plant hollows, respond to flooding by drinking water inside the nest, and excreting it outside.[91] Camponotus anderseni, which nests in the cavities of wood in mangrove habitats, deals with submergence under water by switching to anaerobic respiration.[92]
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Many animals can learn behaviours by imitation, but ants may be the only group apart from mammals where interactive teaching has been observed. A knowledgeable forager of Temnothorax albipennis can lead a naïve nest-mate to newly discovered food by the process of tandem running. The follower obtains knowledge through its leading tutor. The leader is acutely sensitive to the progress of the follower and slows down when the follower lags and speeds up when the follower gets too close.[93]
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Controlled experiments with colonies of Cerapachys biroi suggest that an individual may choose nest roles based on her previous experience. An entire generation of identical workers was divided into two groups whose outcome in food foraging was controlled. One group was continually rewarded with prey, while it was made certain that the other failed. As a result, members of the successful group intensified their foraging attempts while the unsuccessful group ventured out fewer and fewer times. A month later, the successful foragers continued in their role while the others had moved to specialise in brood care.[94]
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Complex nests are built by many ant species, but other species are nomadic and do not build permanent structures. Ants may form subterranean nests or build them on trees. These nests may be found in the ground, under stones or logs, inside logs, hollow stems, or even acorns. The materials used for construction include soil and plant matter,[67] and ants carefully select their nest sites; Temnothorax albipennis will avoid sites with dead ants, as these may indicate the presence of pests or disease. They are quick to abandon established nests at the first sign of threats.[95]
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The army ants of South America, such as the Eciton burchellii species, and the driver ants of Africa do not build permanent nests, but instead, alternate between nomadism and stages where the workers form a temporary nest (bivouac) from their own bodies, by holding each other together.[96]
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Weaver ant (Oecophylla spp.) workers build nests in trees by attaching leaves together, first pulling them together with bridges of workers and then inducing their larvae to produce silk as they are moved along the leaf edges. Similar forms of nest construction are seen in some species of Polyrhachis.[97]
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Formica polyctena, among other ant species, constructs nests that maintain a relatively constant interior temperature that aids in the development of larvae. The ants maintain the nest temperature by choosing the location, nest materials, controlling ventilation and maintaining the heat from solar radiation, worker activity and metabolism, and in some moist nests, microbial activity in the nest materials.[98]
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Some ant species, such as those that use natural cavities, can be opportunistic and make use of the controlled micro-climate provided inside human dwellings and other artificial structures to house their colonies and nest structures.[99][100]
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Most ants are generalist predators, scavengers, and indirect herbivores,[16] but a few have evolved specialised ways of obtaining nutrition. It is believed that many ant species that engage in indirect herbivory rely on specialized symbiosis with their gut microbes [101] to upgrade the nutritional value of the food they collect [102] and allow them to survive in nitrogen poor regions, such as rainforest canopies.[103] Leafcutter ants (Atta and Acromyrmex) feed exclusively on a fungus that grows only within their colonies. They continually collect leaves which are taken to the colony, cut into tiny pieces and placed in fungal gardens. Ergates specialise in related tasks according to their sizes. The largest ants cut stalks, smaller workers chew the leaves and the smallest tend the fungus. Leafcutter ants are sensitive enough to recognise the reaction of the fungus to different plant material, apparently detecting chemical signals from the fungus. If a particular type of leaf is found to be toxic to the fungus, the colony will no longer collect it. The ants feed on structures produced by the fungi called gongylidia. Symbiotic bacteria on the exterior surface of the ants produce antibiotics that kill bacteria introduced into the nest that may harm the fungi.[104]
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Foraging ants travel distances of up to 200 metres (700 ft) from their nest [105] and scent trails allow them to find their way back even in the dark. In hot and arid regions, day-foraging ants face death by desiccation, so the ability to find the shortest route back to the nest reduces that risk. Diurnal desert ants of the genus Cataglyphis such as the Sahara desert ant navigate by keeping track of direction as well as distance travelled. Distances travelled are measured using an internal pedometer that keeps count of the steps taken[106] and also by evaluating the movement of objects in their visual field (optical flow).[107] Directions are measured using the position of the sun.[108]
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They integrate this information to find the shortest route back to their nest.[109]
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Like all ants, they can also make use of visual landmarks when available[110] as well as olfactory and tactile cues to navigate.[111][112] Some species of ant are able to use the Earth's magnetic field for navigation.[113] The compound eyes of ants have specialised cells that detect polarised light from the Sun, which is used to determine direction.[114][115]
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These polarization detectors are sensitive in the ultraviolet region of the light spectrum.[116] In some army ant species, a group of foragers who become separated from the main column may sometimes turn back on themselves and form a circular ant mill. The workers may then run around continuously until they die of exhaustion.[117]
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The female worker ants do not have wings and reproductive females lose their wings after their mating flights in order to begin their colonies. Therefore, unlike their wasp ancestors, most ants travel by walking. Some species are capable of leaping. For example, Jerdon's jumping ant (Harpegnathos saltator) is able to jump by synchronising the action of its mid and hind pairs of legs.[118] There are several species of gliding ant including Cephalotes atratus; this may be a common trait among arboreal ants with small colonies. Ants with this ability are able to control their horizontal movement so as to catch tree trunks when they fall from atop the forest canopy.[119]
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Other species of ants can form chains to bridge gaps over water, underground, or through spaces in vegetation. Some species also form floating rafts that help them survive floods.[120] These rafts may also have a role in allowing ants to colonise islands.[121] Polyrhachis sokolova, a species of ant found in Australian mangrove swamps, can swim and live in underwater nests. Since they lack gills, they go to trapped pockets of air in the submerged nests to breathe.[122]
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Not all ants have the same kind of societies. The Australian bulldog ants are among the biggest and most basal of ants. Like virtually all ants, they are eusocial, but their social behaviour is poorly developed compared to other species. Each individual hunts alone, using her large eyes instead of chemical senses to find prey.[123]
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Some species (such as Tetramorium caespitum) attack and take over neighbouring ant colonies. Others are less expansionist, but just as aggressive; they invade colonies to steal eggs or larvae, which they either eat or raise as workers or slaves. Extreme specialists among these slave-raiding ants, such as the Amazon ants, are incapable of feeding themselves and need captured workers to survive.[124] Captured workers of enslaved Temnothorax species have evolved a counter strategy, destroying just the female pupae of the slave-making Temnothorax americanus, but sparing the males (who don't take part in slave-raiding as adults).[125]
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Ants identify kin and nestmates through their scent, which comes from hydrocarbon-laced secretions that coat their exoskeletons. If an ant is separated from its original colony, it will eventually lose the colony scent. Any ant that enters a colony without a matching scent will be attacked.[126] Also, the reason why two separate colonies of ants will attack each other even if they are of the same species is because the genes responsible for pheromone production are different between them. The Argentine ant, however, does not have this characteristic, due to lack of genetic diversity, and has become a global pest because of it.
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Parasitic ant species enter the colonies of host ants and establish themselves as social parasites; species such as Strumigenys xenos are entirely parasitic and do not have workers, but instead, rely on the food gathered by their Strumigenys perplexa hosts.[127][128] This form of parasitism is seen across many ant genera, but the parasitic ant is usually a species that is closely related to its host. A variety of methods are employed to enter the nest of the host ant. A parasitic queen may enter the host nest before the first brood has hatched, establishing herself prior to development of a colony scent. Other species use pheromones to confuse the host ants or to trick them into carrying the parasitic queen into the nest. Some simply fight their way into the nest.[129]
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A conflict between the sexes of a species is seen in some species of ants with these reproducers apparently competing to produce offspring that are as closely related to them as possible. The most extreme form involves the production of clonal offspring. An extreme of sexual conflict is seen in Wasmannia auropunctata, where the queens produce diploid daughters by thelytokous parthenogenesis and males produce clones by a process whereby a diploid egg loses its maternal contribution to produce haploid males who are clones of the father.[130]
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Ants form symbiotic associations with a range of species, including other ant species, other insects, plants, and fungi. They also are preyed on by many animals and even certain fungi. Some arthropod species spend part of their lives within ant nests, either preying on ants, their larvae, and eggs, consuming the food stores of the ants, or avoiding predators. These inquilines may bear a close resemblance to ants. The nature of this ant mimicry (myrmecomorphy) varies, with some cases involving Batesian mimicry, where the mimic reduces the risk of predation. Others show Wasmannian mimicry, a form of mimicry seen only in inquilines.[131][132]
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Aphids and other hemipteran insects secrete a sweet liquid called honeydew, when they feed on plant sap. The sugars in honeydew are a high-energy food source, which many ant species collect.[133] In some cases, the aphids secrete the honeydew in response to ants tapping them with their antennae. The ants in turn keep predators away from the aphids and will move them from one feeding location to another. When migrating to a new area, many colonies will take the aphids with them, to ensure a continued supply of honeydew. Ants also tend mealybugs to harvest their honeydew. Mealybugs may become a serious pest of pineapples if ants are present to protect mealybugs from their natural enemies.[134]
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Myrmecophilous (ant-loving) caterpillars of the butterfly family Lycaenidae (e.g., blues, coppers, or hairstreaks) are herded by the ants, led to feeding areas in the daytime, and brought inside the ants' nest at night. The caterpillars have a gland which secretes honeydew when the ants massage them. Some caterpillars produce vibrations and sounds that are perceived by the ants.[135] A similar adaptation can be seen in Grizzled skipper butterflies that emit vibrations by expanding their wings in order to communicate with ants, which are natural predators of these butterflies.[136] Other caterpillars have evolved from ant-loving to ant-eating: these myrmecophagous caterpillars secrete a pheromone that makes the ants act as if the caterpillar is one of their own larvae. The caterpillar is then taken into the ant nest where it feeds on the ant larvae.[137] A number of specialized bacterial have been found as endosymbionts in ant guts. Some of the dominant bacteria belong to the order Rhizobiales whose members are known for being nitrogen-fixing symbionts in legumes but the species found in ant lack the ability to fix nitrogen.[138][139] Fungus-growing ants that make up the tribe Attini, including leafcutter ants, cultivate certain species of fungus in the genera Leucoagaricus or Leucocoprinus of the family Agaricaceae. In this ant-fungus mutualism, both species depend on each other for survival. The ant Allomerus decemarticulatus has evolved a three-way association with the host plant, Hirtella physophora (Chrysobalanaceae), and a sticky fungus which is used to trap their insect prey.[140]
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Lemon ants make devil's gardens by killing surrounding plants with their stings and leaving a pure patch of lemon ant trees, (Duroia hirsuta). This modification of the forest provides the ants with more nesting sites inside the stems of the Duroia trees.[141] Although some ants obtain nectar from flowers, pollination by ants is somewhat rare, one example being of the pollination of the orchid Leporella fimbriata which induces male Myrmecia urens to pseudocopulate with the flowers, transferring pollen in the process.[142] One theory that has been proposed for the rarity of pollination is that the secretions of the metapleural gland inactivate and reduce the viability of pollen.[143][144] Some plants have special nectar exuding structures, extrafloral nectaries, that provide food for ants, which in turn protect the plant from more damaging herbivorous insects.[145] Species such as the bullhorn acacia (Acacia cornigera) in Central America have hollow thorns that house colonies of stinging ants (Pseudomyrmex ferruginea) who defend the tree against insects, browsing mammals, and epiphytic vines. Isotopic labelling studies suggest that plants also obtain nitrogen from the ants.[146] In return, the ants obtain food from protein- and lipid-rich Beltian bodies. In Fiji Philidris nagasau (Dolichoderinae) are known to selectively grow species of epiphytic Squamellaria (Rubiaceae) which produce large domatia inside which the ant colonies nest. The ants plant the seeds and the domatia of young seedling are immediately occupied and the ant faeces in them contribute to rapid growth.[147] Similar dispersal associations are found with other dolichoderines in the region as well.[148] Another example of this type of ectosymbiosis comes from the Macaranga tree, which has stems adapted to house colonies of Crematogaster ants.[149]
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Many plant species have seeds that are adapted for dispersal by ants.[150] Seed dispersal by ants or myrmecochory is widespread, and new estimates suggest that nearly 9% of all plant species may have such ant associations.[151][150] Often, seed-dispersing ants perform directed dispersal, depositing the seeds in locations that increase the likelihood of seed survival to reproduction.[152] Some plants in arid, fire-prone systems are particularly dependent on ants for their survival and dispersal as the seeds are transported to safety below the ground.[153] Many ant-dispersed seeds have special external structures, elaiosomes, that are sought after by ants as food.[154]
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A convergence, possibly a form of mimicry, is seen in the eggs of stick insects. They have an edible elaiosome-like structure and are taken into the ant nest where the young hatch.[155]
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Most ants are predatory and some prey on and obtain food from other social insects including other ants. Some species specialise in preying on termites (Megaponera and Termitopone) while a few Cerapachyinae prey on other ants.[105] Some termites, including Nasutitermes corniger, form associations with certain ant species to keep away predatory ant species.[156] The tropical wasp Mischocyttarus drewseni coats the pedicel of its nest with an ant-repellent chemical.[157] It is suggested that many tropical wasps may build their nests in trees and cover them to protect themselves from ants. Other wasps, such as A. multipicta, defend against ants by blasting them off the nest with bursts of wing buzzing.[158] Stingless bees (Trigona and Melipona) use chemical defences against ants.[105]
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Flies in the Old World genus Bengalia (Calliphoridae) prey on ants and are kleptoparasites, snatching prey or brood from the mandibles of adult ants.[159] Wingless and legless females of the Malaysian phorid fly (Vestigipoda myrmolarvoidea) live in the nests of ants of the genus Aenictus and are cared for by the ants.[159]
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Fungi in the genera Cordyceps and Ophiocordyceps infect ants. Ants react to their infection by climbing up plants and sinking their mandibles into plant tissue. The fungus kills the ants, grows on their remains, and produces a fruiting body. It appears that the fungus alters the behaviour of the ant to help disperse its spores [160] in a microhabitat that best suits the fungus.[161] Strepsipteran parasites also manipulate their ant host to climb grass stems, to help the parasite find mates.[162]
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A nematode (Myrmeconema neotropicum) that infects canopy ants (Cephalotes atratus) causes the black-coloured gasters of workers to turn red. The parasite also alters the behaviour of the ant, causing them to carry their gasters high. The conspicuous red gasters are mistaken by birds for ripe fruits, such as Hyeronima alchorneoides, and eaten. The droppings of the bird are collected by other ants and fed to their young, leading to further spread of the nematode.[163]
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South American poison dart frogs in the genus Dendrobates feed mainly on ants, and the toxins in their skin may come from the ants.[164]
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Army ants forage in a wide roving column, attacking any animals in that path that are unable to escape. In Central and South America, Eciton burchellii is the swarming ant most commonly attended by "ant-following" birds such as antbirds and woodcreepers.[165][166] This behaviour was once considered mutualistic, but later studies found the birds to be parasitic. Direct kleptoparasitism (birds stealing food from the ants' grasp) is rare and has been noted in Inca doves which pick seeds at nest entrances as they are being transported by species of Pogonomyrmex.[167] Birds that follow ants eat many prey insects and thus decrease the foraging success of ants.[168] Birds indulge in a peculiar behaviour called anting that, as yet, is not fully understood. Here birds rest on ant nests, or pick and drop ants onto their wings and feathers; this may be a means to remove ectoparasites from the birds.
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Anteaters, aardvarks, pangolins, echidnas and numbats have special adaptations for living on a diet of ants. These adaptations include long, sticky tongues to capture ants and strong claws to break into ant nests. Brown bears (Ursus arctos) have been found to feed on ants. About 12%, 16%, and 4% of their faecal volume in spring, summer, and autumn, respectively, is composed of ants.[169]
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Ants perform many ecological roles that are beneficial to humans, including the suppression of pest populations and aeration of the soil. The use of weaver ants in citrus cultivation in southern China is considered one of the oldest known applications of biological control.[170] On the other hand, ants may become nuisances when they invade buildings, or cause economic losses.
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In some parts of the world (mainly Africa and South America), large ants, especially army ants, are used as surgical sutures. The wound is pressed together and ants are applied along it. The ant seizes the edges of the wound in its mandibles and locks in place. The body is then cut off and the head and mandibles remain in place to close the wound.[171][172][173] The large heads of the dinergates (soldiers) of the leafcutting ant Atta cephalotes are also used by native surgeons in closing wounds.[174]
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Some ants have toxic venom and are of medical importance. The species include Paraponera clavata (tocandira) and Dinoponera spp. (false tocandiras) of South America [175] and the Myrmecia ants of Australia.[176]
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In South Africa, ants are used to help harvest the seeds of rooibos (Aspalathus linearis), a plant used to make a herbal tea. The plant disperses its seeds widely, making manual collection difficult. Black ants collect and store these and other seeds in their nest, where humans can gather them en masse. Up to half a pound (200 g) of seeds may be collected from one ant-heap.[177][178]
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Although most ants survive attempts by humans to eradicate them, a few are highly endangered. These tend to be island species that have evolved specialized traits and risk being displaced by introduced ant species. Examples include the critically endangered Sri Lankan relict ant (Aneuretus simoni) and Adetomyrma venatrix of Madagascar.[179]
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It has been estimated by E.O. Wilson that the total number of individual ants alive in the world at any one time is between one and ten quadrillion (short scale) (i.e., between 1015 and 1016). According to this estimate, the total biomass of all the ants in the world is approximately equal to the total biomass of the entire human race.[180] Also, according to this estimate, there are approximately 1 million ants for every human on Earth.[181]
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Ants and their larvae are eaten in different parts of the world. The eggs of two species of ants are used in Mexican escamoles. They are considered a form of insect caviar and can sell for as much as US$40 per pound ($90/kg) because they are seasonal and hard to find. In the Colombian department of Santander, hormigas culonas (roughly interpreted as "large-bottomed ants") Atta laevigata are toasted alive and eaten.[182]
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In areas of India, and throughout Burma and Thailand, a paste of the green weaver ant (Oecophylla smaragdina) is served as a condiment with curry.[183] Weaver ant eggs and larvae, as well as the ants, may be used in a Thai salad, yam (Thai: ยำ), in a dish called yam khai mot daeng (Thai: ยำไข่มดแดง) or red ant egg salad, a dish that comes from the Issan or north-eastern region of Thailand. Saville-Kent, in the Naturalist in Australia wrote "Beauty, in the case of the green ant, is more than skin-deep. Their attractive, almost sweetmeat-like translucency possibly invited the first essays at their consumption by the human species". Mashed up in water, after the manner of lemon squash, "these ants form a pleasant acid drink which is held in high favor by the natives of North Queensland, and is even appreciated by many European palates".[184]
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In his First Summer in the Sierra, John Muir notes that the Digger Indians of California ate the tickling, acid gasters of the large jet-black carpenter ants. The Mexican Indians eat the replete workers, or living honey-pots, of the honey ant (Myrmecocystus).[184]
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Some ant species are considered as pests, primarily those that occur in human habitations, where their presence is often problematic. For example, the presence of ants would be undesirable in sterile places such as hospitals or kitchens. Some species or genera commonly categorized as pests include the Argentine ant, pavement ant, yellow crazy ant, banded sugar ant, pharaoh ant, red ant, carpenter ant, odorous house ant, red imported fire ant, and European fire ant. Some ants will raid stored food, some will seek water sources, others may damage indoor structures, some may damage agricultural crops directly (or by aiding sucking pests). Some will sting or bite.[185] The adaptive nature of ant colonies make it nearly impossible to eliminate entire colonies and most pest management practices aim to control local populations and tend to be temporary solutions. Ant populations are managed by a combination of approaches that make use of chemical, biological, and physical methods. Chemical methods include the use of insecticidal bait which is gathered by ants as food and brought back to the nest where the poison is inadvertently spread to other colony members through trophallaxis. Management is based on the species and techniques may vary according to the location and circumstance.[185]
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Observed by humans since the dawn of history, the behaviour of ants has been documented and the subject of early writings and fables passed from one century to another. Those using scientific methods, myrmecologists, study ants in the laboratory and in their natural conditions. Their complex and variable social structures have made ants ideal model organisms. Ultraviolet vision was first discovered in ants by Sir John Lubbock in 1881.[186] Studies on ants have tested hypotheses in ecology and sociobiology, and have been particularly important in examining the predictions of theories of kin selection and evolutionarily stable strategies.[187] Ant colonies may be studied by rearing or temporarily maintaining them in formicaria, specially constructed glass framed enclosures.[188] Individuals may be tracked for study by marking them with dots of colours.[189]
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The successful techniques used by ant colonies have been studied in computer science and robotics to produce distributed and fault-tolerant systems for solving problems, for example Ant colony optimization and Ant robotics. This area of biomimetics has led to studies of ant locomotion, search engines that make use of "foraging trails", fault-tolerant storage, and networking algorithms.[190]
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From the late 1950s through the late 1970s, ant farms were popular educational children's toys in the United States. Some later commercial versions use transparent gel instead of soil, allowing greater visibility at the cost of stressing the ants with unnatural light.[191]
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Anthropomorphised ants have often been used in fables and children's stories to represent industriousness and cooperative effort. They also are mentioned in religious texts.[192][193] In the Book of Proverbs in the Bible, ants are held up as a good example for humans for their hard work and cooperation.[194] Aesop did the same in his fable The Ant and the Grasshopper. In the Quran, Sulayman is said to have heard and understood an ant warning other ants to return home to avoid being accidentally crushed by Sulayman and his marching army.[Quran 27:18][195] In parts of Africa, ants are considered to be the messengers of the deities. Some Native American mythology, such as the Hopi mythology, considers ants as the very first animals. Ant bites are often said to have curative properties. The sting of some species of Pseudomyrmex is claimed to give fever relief.[196] Ant bites are used in the initiation ceremonies of some Amazon Indian cultures as a test of endurance.[197][198]
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Ant society has always fascinated humans and has been written about both humorously and seriously. Mark Twain wrote about ants in his 1880 book A Tramp Abroad.[199] Some modern authors have used the example of the ants to comment on the relationship between society and the individual. Examples are Robert Frost in his poem "Departmental" and T. H. White in his fantasy novel The Once and Future King. The plot in French entomologist and writer Bernard Werber's Les Fourmis science-fiction trilogy is divided between the worlds of ants and humans; ants and their behaviour is described using contemporary scientific knowledge. H.G. Wells wrote about intelligent ants destroying human settlements in Brazil and threatening human civilization in his 1905 science-fiction short story, The Empire of the Ants. In more recent times, animated cartoons and 3-D animated films featuring ants have been produced including Antz, A Bug's Life, The Ant Bully, The Ant and the Aardvark, Ferdy the Ant and Atom Ant. Renowned myrmecologist E. O. Wilson wrote a short story, "Trailhead" in 2010 for The New Yorker magazine, which describes the life and death of an ant-queen and the rise and fall of her colony, from an ants' point of view.[200] The French neuroanatomist, psychiatrist and eugenicist Auguste Forel believed that ant societies were models for human society. He published a five volume work from 1921 to 1923 that examined ant biology and society.[201]
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In the early 1990s, the video game SimAnt, which simulated an ant colony, won the 1992 Codie award for "Best Simulation Program".[202]
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Ants also are quite popular inspiration for many science-fiction insectoids, such as the Formics of Ender's Game, the Bugs of Starship Troopers, the giant ants in the films Them! and Empire of the Ants, Marvel Comics' super hero Ant-Man, and ants mutated into super-intelligence in Phase IV. In computer strategy games, ant-based species often benefit from increased production rates due to their single-minded focus, such as the Klackons in the Master of Orion series of games or the ChCht in Deadlock II. These characters are often credited with a hive mind, a common misconception about ant colonies.[203]
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Earth is the third planet from the Sun and the only astronomical object known to harbor life. According to radiometric dating estimation and other evidence, Earth formed over 4.5 billion years ago. Earth's gravity interacts with other objects in space, especially the Sun and the Moon, which is Earth's only natural satellite. Earth orbits around the Sun in 365.256 solar days, a period known as an Earth sidereal year. During this time, Earth rotates about its axis 366.256 times, that is, a sidereal year has 366.256 sidereal days.[n 6]
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Earth's axis of rotation is tilted with respect to its orbital plane, producing seasons on Earth. The gravitational interaction between Earth and the Moon causes tides, stabilizes Earth's orientation on its axis, and gradually slows its rotation. Earth is the densest planet in the Solar System and the largest and most massive of the four rocky planets.
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Earth's outer layer (lithosphere) is divided into several rigid tectonic plates that migrate across the surface over many millions of years. About 29% of Earth's surface is land consisting of continents and islands. The remaining 71% is covered with water, mostly by oceans but also lakes, rivers and other fresh water, which all together constitute the hydrosphere. The majority of Earth's polar regions are covered in ice, including the Antarctic ice sheet and the sea ice of the Arctic ice pack. Earth's interior remains active with a solid iron inner core, a liquid outer core that generates Earth's magnetic field, and a convecting mantle that drives plate tectonics.
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Within the first billion years of Earth's history, life appeared in the oceans and began to affect Earth's atmosphere and surface, leading to the proliferation of anaerobic and, later, aerobic organisms. Some geological evidence indicates that life may have arisen as early as 4.1 billion years ago. Since then, the combination of Earth's distance from the Sun, physical properties and geological history have allowed life to evolve and thrive. In the history of life on Earth, biodiversity has gone through long periods of expansion, occasionally punctuated by mass extinctions. Over 99% of all species that ever lived on Earth are extinct. Estimates of the number of species on Earth today vary widely; most species have not been described. Over 7.7 billion humans live on Earth and depend on its biosphere and natural resources for their survival.[23]
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The modern English word Earth developed, via Middle English,[n 7] from an Old English noun most often spelled eorðe.[24] It has cognates in every Germanic language, and their ancestral root has been reconstructed as *erþō. In its earliest attestation, the word eorðe was already being used to translate the many senses of Latin terra and Greek γῆ gē: the ground,[n 8] its soil,[n 9] dry land,[n 10] the human world,[n 11] the surface of the world (including the sea),[n 12] and the globe itself.[n 13] As with Roman Terra/Tellūs and Greek Gaia, Earth may have been a personified goddess in Germanic paganism: late Norse mythology included Jörð ('Earth'), a giantess often given as the mother of Thor.[33]
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Originally, earth was written in lowercase, and from early Middle English, its definite sense as "the globe" was expressed as the earth. By Early Modern English, many nouns were capitalized, and the earth became (and often remained) the Earth, particularly when referenced along with other heavenly bodies. More recently, the name is sometimes simply given as Earth, by analogy with the names of the other planets.[24] House styles now vary: Oxford spelling recognizes the lowercase form as the most common, with the capitalized form an acceptable variant. Another convention capitalizes "Earth" when appearing as a name (e.g. "Earth's atmosphere") but writes it in lowercase when preceded by the (e.g. "the atmosphere of the earth"). It almost always appears in lowercase in colloquial expressions such as "what on earth are you doing?"[34]
|
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Occasionally, the name Terra /ˈtɛrə/ is used in scientific writing and especially in science fiction to distinguish our inhabited planet from others,[35] while in poetry Tellus /ˈtɛləs/ has been used to denote personification of the Earth.[36] The Greek poetic name Gaea (Gæa) /ˈdʒiːə/ is rare, though the alternative spelling Gaia has become common due to the Gaia hypothesis, in which case its pronunciation is /ˈɡaɪə/ rather than the more Classical /ˈɡeɪə/.[37]
|
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There are a number of adjectives for the planet Earth. From Earth itself comes earthly. From Latin Terra come Terran /ˈtɛrən/,[38] Terrestrial /təˈrɛstriəl/,[39] and (via French) Terrene /təˈriːn/,[40] and from Latin Tellus come Tellurian /tɛˈlʊəriən/[41] and, more rarely, Telluric and Tellural. From Greek Gaia and Gaea comes Gaian and Gaean.
|
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An inhabitant of the Earth is an Earthling, a Terran, a Terrestrial, a Tellurian or, rarely, an Earthian.
|
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The oldest material found in the Solar System is dated to 4.5672±0.0006 billion years ago (BYA).[42] By 4.54±0.04 BYA[43] the primordial Earth had formed. The bodies in the Solar System formed and evolved with the Sun. In theory, a solar nebula partitions a volume out of a molecular cloud by gravitational collapse, which begins to spin and flatten into a circumstellar disk, and then the planets grow out of that disk with the Sun. A nebula contains gas, ice grains, and dust (including primordial nuclides). According to nebular theory, planetesimals formed by accretion, with the primordial Earth taking 10–20 million years (Mys) to form.[44]
|
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A subject of research is the formation of the Moon, some 4.53 BYA.[45] A leading hypothesis is that it was formed by accretion from material loosed from Earth after a Mars-sized object, named Theia, hit Earth.[46] In this view, the mass of Theia was approximately 10 percent of Earth;[47] it hit Earth with a glancing blow and some of its mass merged with Earth.[48] Between approximately 4.1 and 3.8 BYA, numerous asteroid impacts during the Late Heavy Bombardment caused significant changes to the greater surface environment of the Moon and, by inference, to that of Earth.
|
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Earth's atmosphere and oceans were formed by volcanic activity and outgassing.[49] Water vapor from these sources condensed into the oceans, augmented by water and ice from asteroids, protoplanets, and comets.[50] In this model, atmospheric "greenhouse gases" kept the oceans from freezing when the newly forming Sun had only 70% of its current luminosity.[51] By 3.5 BYA, Earth's magnetic field was established, which helped prevent the atmosphere from being stripped away by the solar wind.[52]
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A crust formed when the molten outer layer of Earth cooled to form a solid. The two models[53] that explain land mass propose either a steady growth to the present-day forms[54] or, more likely, a rapid growth[55] early in Earth history[56] followed by a long-term steady continental area.[57][58][59] Continents formed by plate tectonics, a process ultimately driven by the continuous loss of heat from Earth's interior. Over the period of hundreds of millions of years, the supercontinents have assembled and broken apart. Roughly 750 million years ago (MYA), one of the earliest known supercontinents, Rodinia, began to break apart. The continents later recombined to form Pannotia 600–540 MYA, then finally Pangaea, which also broke apart 180 MYA.[60]
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The present pattern of ice ages began about 40 MYA,[61] and then intensified during the Pleistocene about 3 MYA.[62] High-latitude regions have since undergone repeated cycles of glaciation and thaw, repeating about every 40,000–100,000 years. The last continental glaciation ended 10,000 years ago.[63]
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Chemical reactions led to the first self-replicating molecules about four billion years ago. A half billion years later, the last common ancestor of all current life arose.[64] The evolution of photosynthesis allowed the Sun's energy to be harvested directly by life forms. The resultant molecular oxygen (O2) accumulated in the atmosphere and due to interaction with ultraviolet solar radiation, formed a protective ozone layer (O3) in the upper atmosphere.[65] The incorporation of smaller cells within larger ones resulted in the development of complex cells called eukaryotes.[66] True multicellular organisms formed as cells within colonies became increasingly specialized. Aided by the absorption of harmful ultraviolet radiation by the ozone layer, life colonized Earth's surface.[67] Among the earliest fossil evidence for life is microbial mat fossils found in 3.48 billion-year-old sandstone in Western Australia,[68] biogenic graphite found in 3.7 billion-year-old metasedimentary rocks in Western Greenland,[69] and remains of biotic material found in 4.1 billion-year-old rocks in Western Australia.[70][71] The earliest direct evidence of life on Earth is contained in 3.45 billion-year-old Australian rocks showing fossils of microorganisms.[72][73]
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During the Neoproterozoic, 750 to 580 MYA, much of Earth might have been covered in ice. This hypothesis has been termed "Snowball Earth", and it is of particular interest because it preceded the Cambrian explosion, when multicellular life forms significantly increased in complexity.[74] Following the Cambrian explosion, 535 MYA, there have been five mass extinctions.[75] The most recent such event was 66 MYA, when an asteroid impact triggered the extinction of the non-avian dinosaurs and other large reptiles, but spared some small animals such as mammals, which at the time resembled shrews. Mammalian life has diversified over the past 66 Mys, and several million years ago an African ape-like animal such as Orrorin tugenensis gained the ability to stand upright.[76] This facilitated tool use and encouraged communication that provided the nutrition and stimulation needed for a larger brain, which led to the evolution of humans. The development of agriculture, and then civilization, led to humans having an influence on Earth and the nature and quantity of other life forms that continues to this day.[77]
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Earth's expected long-term future is tied to that of the Sun. Over the next 1.1 billion years, solar luminosity will increase by 10%, and over the next 3.5 billion years by 40%.[78] Earth's increasing surface temperature will accelerate the inorganic carbon cycle, reducing CO2 concentration to levels lethally low for plants (10 ppm for C4 photosynthesis) in approximately 100–900 million years.[79][80] The lack of vegetation will result in the loss of oxygen in the atmosphere, making animal life impossible.[81] About a billion years from now, all surface water will have disappeared[82] and the mean global temperature will reach 70 °C (158 °F).[81] Earth is expected to be habitable until the end of photosynthesis about 500 million years from now,[79] but if nitrogen is removed from the atmosphere, life may continue until a runaway greenhouse effect occurs 2.3 billion years from now.[80] Anthropogenic emissions are "probably insufficient" to cause a runaway greenhouse at current solar luminosity.[83] Even if the Sun were eternal and stable, 27% of the water in the modern oceans will descend to the mantle in one billion years, due to reduced steam venting from mid-ocean ridges.[84]
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The Sun will evolve to become a red giant in about 5 billion years. Models predict that the Sun will expand to roughly 1 AU (150 million km; 93 million mi), about 250 times its present radius.[78][85] Earth's fate is less clear. As a red giant, the Sun will lose roughly 30% of its mass, so, without tidal effects, Earth will move to an orbit 1.7 AU (250 million km; 160 million mi) from the Sun when the star reaches its maximum radius. Most, if not all, remaining life will be destroyed by the Sun's increased luminosity (peaking at about 5,000 times its present level).[78] A 2008 simulation indicates that Earth's orbit will eventually decay due to tidal effects and drag, causing it to enter the Sun's atmosphere and be vaporized.[85]
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The shape of Earth is nearly spherical. There is a small flattening at the poles and bulging around the equator due to Earth's rotation.[89] To second order, Earth is approximately an oblate spheroid, whose equatorial diameter is 43 kilometres (27 mi) larger than the pole-to-pole diameter,[90] although the variation is less than 1% of the average radius of the Earth.
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The point on the surface farthest from Earth's center of mass is the summit of the equatorial Chimborazo volcano in Ecuador (6,384.4 km or 3,967.1 mi).[91][92][93][94] The average diameter of the reference spheroid is 12,742 kilometres (7,918 mi). Local topography deviates from this idealized spheroid, although on a global scale these deviations are small compared to Earth's radius: the maximum deviation of only 0.17% is at the Mariana Trench (10,911 metres or 35,797 feet below local sea level), whereas Mount Everest (8,848 metres or 29,029 feet above local sea level) represents a deviation of 0.14%.[n 14]
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In geodesy, the exact shape that Earth's oceans would adopt in the absence of land and perturbations such as tides and winds is called the geoid. More precisely, the geoid is the surface of gravitational equipotential at mean sea level.
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Earth's mass is approximately 5.97×1024 kg (5,970 Yg). It is composed mostly of iron (32.1%), oxygen (30.1%), silicon (15.1%), magnesium (13.9%), sulphur (2.9%), nickel (1.8%), calcium (1.5%), and aluminum (1.4%), with the remaining 1.2% consisting of trace amounts of other elements. Due to mass segregation, the core region is estimated to be primarily composed of iron (88.8%), with smaller amounts of nickel (5.8%), sulphur (4.5%), and less than 1% trace elements.[98]
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The most common rock constituents of the crust are nearly all oxides: chlorine, sulphur, and fluorine are the important exceptions to this and their total amount in any rock is usually much less than 1%. Over 99% of the crust is composed of 11 oxides, principally silica, alumina, iron oxides, lime, magnesia, potash and soda.[99][98][100]
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Earth's interior, like that of the other terrestrial planets, is divided into layers by their chemical or physical (rheological) properties. The outer layer is a chemically distinct silicate solid crust, which is underlain by a highly viscous solid mantle. The crust is separated from the mantle by the Mohorovičić discontinuity. The thickness of the crust varies from about 6 kilometres (3.7 mi) under the oceans to 30–50 km (19–31 mi) for the continents. The crust and the cold, rigid, top of the upper mantle are collectively known as the lithosphere, and it is of the lithosphere that the tectonic plates are composed. Beneath the lithosphere is the asthenosphere, a relatively low-viscosity layer on which the lithosphere rides. Important changes in crystal structure within the mantle occur at 410 and 660 km (250 and 410 mi) below the surface, spanning a transition zone that separates the upper and lower mantle. Beneath the mantle, an extremely low viscosity liquid outer core lies above a solid inner core.[101] Earth's inner core might rotate at a slightly higher angular velocity than the remainder of the planet, advancing by 0.1–0.5° per year.[102] The radius of the inner core is about one fifth of that of Earth.
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Earth's internal heat comes from a combination of residual heat from planetary accretion (about 20%) and heat produced through radioactive decay (80%).[105] The major heat-producing isotopes within Earth are potassium-40, uranium-238, and thorium-232.[106] At the center, the temperature may be up to 6,000 °C (10,830 °F),[107] and the pressure could reach 360 GPa (52 million psi).[108] Because much of the heat is provided by radioactive decay, scientists postulate that early in Earth's history, before isotopes with short half-lives were depleted, Earth's heat production was much higher. At approximately 3 Gyr, twice the present-day heat would have been produced, increasing the rates of mantle convection and plate tectonics, and allowing the production of uncommon igneous rocks such as komatiites that are rarely formed today.[105][109]
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The mean heat loss from Earth is 87 mW m−2, for a global heat loss of 4.42×1013 W.[111] A portion of the core's thermal energy is transported toward the crust by mantle plumes, a form of convection consisting of upwellings of higher-temperature rock. These plumes can produce hotspots and flood basalts.[112] More of the heat in Earth is lost through plate tectonics, by mantle upwelling associated with mid-ocean ridges. The final major mode of heat loss is through conduction through the lithosphere, the majority of which occurs under the oceans because the crust there is much thinner than that of the continents.[113]
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Earth's mechanically rigid outer layer, the lithosphere, is divided into tectonic plates. These plates are rigid segments that move relative to each other at one of three boundaries types: At convergent boundaries, two plates come together; at divergent boundaries, two plates are pulled apart; and at transform boundaries, two plates slide past one another laterally. Along these plate boundaries, earthquakes, volcanic activity, mountain-building, and oceanic trench formation can occur.[115] The tectonic plates ride on top of the asthenosphere, the solid but less-viscous part of the upper mantle that can flow and move along with the plates.[116]
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As the tectonic plates migrate, oceanic crust is subducted under the leading edges of the plates at convergent boundaries. At the same time, the upwelling of mantle material at divergent boundaries creates mid-ocean ridges. The combination of these processes recycles the oceanic crust back into the mantle. Due to this recycling, most of the ocean floor is less than 100 Ma old. The oldest oceanic crust is located in the Western Pacific and is estimated to be 200 Ma old.[117][118] By comparison, the oldest dated continental crust is 4,030 Ma.[119]
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The seven major plates are the Pacific, North American, Eurasian, African, Antarctic, Indo-Australian, and South American. Other notable plates include the Arabian Plate, the Caribbean Plate, the Nazca Plate off the west coast of South America and the Scotia Plate in the southern Atlantic Ocean. The Australian Plate fused with the Indian Plate between 50 and 55 Mya. The fastest-moving plates are the oceanic plates, with the Cocos Plate advancing at a rate of 75 mm/a (3.0 in/year)[120] and the Pacific Plate moving 52–69 mm/a (2.0–2.7 in/year). At the other extreme, the slowest-moving plate is the Eurasian Plate, progressing at a typical rate of 21 mm/a (0.83 in/year).[121]
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The total surface area of Earth is about 510 million km2 (197 million sq mi).[12] Of this, 70.8%,[12] or 361.13 million km2 (139.43 million sq mi), is below sea level and covered by ocean water.[122] Below the ocean's surface are much of the continental shelf, mountains, volcanoes,[90] oceanic trenches, submarine canyons, oceanic plateaus, abyssal plains, and a globe-spanning mid-ocean ridge system. The remaining 29.2%, or 148.94 million km2 (57.51 million sq mi), not covered by water has terrain that varies greatly from place to place and consists of mountains, deserts, plains, plateaus, and other landforms. Tectonics and erosion, volcanic eruptions, flooding, weathering, glaciation, the growth of coral reefs, and meteorite impacts are among the processes that constantly reshape Earth's surface over geological time.[123][124]
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The continental crust consists of lower density material such as the igneous rocks granite and andesite. Less common is basalt, a denser volcanic rock that is the primary constituent of the ocean floors.[125] Sedimentary rock is formed from the accumulation of sediment that becomes buried and compacted together. Nearly 75% of the continental surfaces are covered by sedimentary rocks, although they form about 5% of the crust.[126] The third form of rock material found on Earth is metamorphic rock, which is created from the transformation of pre-existing rock types through high pressures, high temperatures, or both. The most abundant silicate minerals on Earth's surface include quartz, feldspars, amphibole, mica, pyroxene and olivine.[127] Common carbonate minerals include calcite (found in limestone) and dolomite.[128]
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The elevation of the land surface varies from the low point of −418 m (−1,371 ft) at the Dead Sea, to a maximum altitude of 8,848 m (29,029 ft) at the top of Mount Everest. The mean height of land above sea level is about 797 m (2,615 ft).[129]
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The pedosphere is the outermost layer of Earth's continental surface and is composed of soil and subject to soil formation processes. The total arable land is 10.9% of the land surface, with 1.3% being permanent cropland.[130][131] Close to 40% of Earth's land surface is used for agriculture, or an estimated 16.7 million km2 (6.4 million sq mi) of cropland and 33.5 million km2 (12.9 million sq mi) of pastureland.[132]
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The abundance of water on Earth's surface is a unique feature that distinguishes the "Blue Planet" from other planets in the Solar System. Earth's hydrosphere consists chiefly of the oceans, but technically includes all water surfaces in the world, including inland seas, lakes, rivers, and underground waters down to a depth of 2,000 m (6,600 ft). The deepest underwater location is Challenger Deep of the Mariana Trench in the Pacific Ocean with a depth of 10,911.4 m (35,799 ft).[n 18][133]
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The mass of the oceans is approximately 1.35×1018 metric tons or about 1/4400 of Earth's total mass. The oceans cover an area of 361.8 million km2 (139.7 million sq mi) with a mean depth of 3,682 m (12,080 ft), resulting in an estimated volume of 1.332 billion km3 (320 million cu mi).[134] If all of Earth's crustal surface were at the same elevation as a smooth sphere, the depth of the resulting world ocean would be 2.7 to 2.8 km (1.68 to 1.74 mi).[135][136]
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About 97.5% of the water is saline; the remaining 2.5% is fresh water. Most fresh water, about 68.7%, is present as ice in ice caps and glaciers.[137]
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The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5% salt).[138] Most of this salt was released from volcanic activity or extracted from cool igneous rocks.[139] The oceans are also a reservoir of dissolved atmospheric gases, which are essential for the survival of many aquatic life forms.[140] Sea water has an important influence on the world's climate, with the oceans acting as a large heat reservoir.[141] Shifts in the oceanic temperature distribution can cause significant weather shifts, such as the El Niño–Southern Oscillation.[142]
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The atmospheric pressure at Earth's sea level averages 101.325 kPa (14.696 psi),[143] with a scale height of about 8.5 km (5.3 mi).[3] A dry atmosphere is composed of 78.084% nitrogen, 20.946% oxygen, 0.934% argon, and trace amounts of carbon dioxide and other gaseous molecules.[143] Water vapor content varies between 0.01% and 4%[143] but averages about 1%.[3] The height of the troposphere varies with latitude, ranging between 8 km (5 mi) at the poles to 17 km (11 mi) at the equator, with some variation resulting from weather and seasonal factors.[144]
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Earth's biosphere has significantly altered its atmosphere. Oxygenic photosynthesis evolved 2.7 Gya, forming the primarily nitrogen–oxygen atmosphere of today.[65] This change enabled the proliferation of aerobic organisms and, indirectly, the formation of the ozone layer due to the subsequent conversion of atmospheric O2 into O3. The ozone layer blocks ultraviolet solar radiation, permitting life on land.[145] Other atmospheric functions important to life include transporting water vapor, providing useful gases, causing small meteors to burn up before they strike the surface, and moderating temperature.[146] This last phenomenon is known as the greenhouse effect: trace molecules within the atmosphere serve to capture thermal energy emitted from the ground, thereby raising the average temperature. Water vapor, carbon dioxide, methane, nitrous oxide, and ozone are the primary greenhouse gases in the atmosphere. Without this heat-retention effect, the average surface temperature would be −18 °C (0 °F), in contrast to the current +15 °C (59 °F),[147] and life on Earth probably would not exist in its current form.[148] In May 2017, glints of light, seen as twinkling from an orbiting satellite a million miles away, were found to be reflected light from ice crystals in the atmosphere.[149][150]
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Earth's atmosphere has no definite boundary, slowly becoming thinner and fading into outer space. Three-quarters of the atmosphere's mass is contained within the first 11 km (6.8 mi) of the surface. This lowest layer is called the troposphere. Energy from the Sun heats this layer, and the surface below, causing expansion of the air. This lower-density air then rises and is replaced by cooler, higher-density air. The result is atmospheric circulation that drives the weather and climate through redistribution of thermal energy.[151]
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The primary atmospheric circulation bands consist of the trade winds in the equatorial region below 30° latitude and the westerlies in the mid-latitudes between 30° and 60°.[152] Ocean currents are also important factors in determining climate, particularly the thermohaline circulation that distributes thermal energy from the equatorial oceans to the polar regions.[153]
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Water vapor generated through surface evaporation is transported by circulatory patterns in the atmosphere. When atmospheric conditions permit an uplift of warm, humid air, this water condenses and falls to the surface as precipitation.[151] Most of the water is then transported to lower elevations by river systems and usually returned to the oceans or deposited into lakes. This water cycle is a vital mechanism for supporting life on land and is a primary factor in the erosion of surface features over geological periods. Precipitation patterns vary widely, ranging from several meters of water per year to less than a millimeter. Atmospheric circulation, topographic features, and temperature differences determine the average precipitation that falls in each region.[154]
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The amount of solar energy reaching Earth's surface decreases with increasing latitude. At higher latitudes, the sunlight reaches the surface at lower angles, and it must pass through thicker columns of the atmosphere. As a result, the mean annual air temperature at sea level decreases by about 0.4 °C (0.7 °F) per degree of latitude from the equator.[155] Earth's surface can be subdivided into specific latitudinal belts of approximately homogeneous climate. Ranging from the equator to the polar regions, these are the tropical (or equatorial), subtropical, temperate and polar climates.[156]
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This latitudinal rule has several anomalies:
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The commonly used Köppen climate classification system has five broad groups (humid tropics, arid, humid middle latitudes, continental and cold polar), which are further divided into more specific subtypes.[152] The Köppen system rates regions of terrain based on observed temperature and precipitation.
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The highest air temperature ever measured on Earth was 56.7 °C (134.1 °F) in Furnace Creek, California, in Death Valley, in 1913.[159] The lowest air temperature ever directly measured on Earth was −89.2 °C (−128.6 °F) at Vostok Station in 1983,[160] but satellites have used remote sensing to measure temperatures as low as −94.7 °C (−138.5 °F) in East Antarctica.[161] These temperature records are only measurements made with modern instruments from the 20th century onwards and likely do not reflect the full range of temperature on Earth.
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Above the troposphere, the atmosphere is usually divided into the stratosphere, mesosphere, and thermosphere.[146] Each layer has a different lapse rate, defining the rate of change in temperature with height. Beyond these, the exosphere thins out into the magnetosphere, where the geomagnetic fields interact with the solar wind.[162] Within the stratosphere is the ozone layer, a component that partially shields the surface from ultraviolet light and thus is important for life on Earth. The Kármán line, defined as 100 km above Earth's surface, is a working definition for the boundary between the atmosphere and outer space.[163]
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Thermal energy causes some of the molecules at the outer edge of the atmosphere to increase their velocity to the point where they can escape from Earth's gravity. This causes a slow but steady loss of the atmosphere into space. Because unfixed hydrogen has a low molecular mass, it can achieve escape velocity more readily, and it leaks into outer space at a greater rate than other gases.[164] The leakage of hydrogen into space contributes to the shifting of Earth's atmosphere and surface from an initially reducing state to its current oxidizing one. Photosynthesis provided a source of free oxygen, but the loss of reducing agents such as hydrogen is thought to have been a necessary precondition for the widespread accumulation of oxygen in the atmosphere.[165] Hence the ability of hydrogen to escape from the atmosphere may have influenced the nature of life that developed on Earth.[166] In the current, oxygen-rich atmosphere most hydrogen is converted into water before it has an opportunity to escape. Instead, most of the hydrogen loss comes from the destruction of methane in the upper atmosphere.[167]
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The gravity of Earth is the acceleration that is imparted to objects due to the distribution of mass within Earth. Near Earth's surface, gravitational acceleration is approximately 9.8 m/s2 (32 ft/s2). Local differences in topography, geology, and deeper tectonic structure cause local and broad, regional differences in Earth's gravitational field, known as gravity anomalies.[168]
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The main part of Earth's magnetic field is generated in the core, the site of a dynamo process that converts the kinetic energy of thermally and compositionally driven convection into electrical and magnetic field energy. The field extends outwards from the core, through the mantle, and up to Earth's surface, where it is, approximately, a dipole. The poles of the dipole are located close to Earth's geographic poles. At the equator of the magnetic field, the magnetic-field strength at the surface is 3.05×10−5 T, with a magnetic dipole moment of 7.79×1022 Am2 at epoch 2000, decreasing nearly 6% per century.[169] The convection movements in the core are chaotic; the magnetic poles drift and periodically change alignment. This causes secular variation of the main field and field reversals at irregular intervals averaging a few times every million years. The most recent reversal occurred approximately 700,000 years ago.[170][171]
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The extent of Earth's magnetic field in space defines the magnetosphere. Ions and electrons of the solar wind are deflected by the magnetosphere; solar wind pressure compresses the dayside of the magnetosphere, to about 10 Earth radii, and extends the nightside magnetosphere into a long tail.[172] Because the velocity of the solar wind is greater than the speed at which waves propagate through the solar wind, a supersonic bow shock precedes the dayside magnetosphere within the solar wind.[173] Charged particles are contained within the magnetosphere; the plasmasphere is defined by low-energy particles that essentially follow magnetic field lines as Earth rotates;[174][175] the ring current is defined by medium-energy particles that drift relative to the geomagnetic field, but with paths that are still dominated by the magnetic field,[176] and the Van Allen radiation belt are formed by high-energy particles whose motion is essentially random, but otherwise contained by the magnetosphere.[172][177]
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During magnetic storms and substorms, charged particles can be deflected from the outer magnetosphere and especially the magnetotail, directed along field lines into Earth's ionosphere, where atmospheric atoms can be excited and ionized, causing the aurora.[178]
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Earth's rotation period relative to the Sun—its mean solar day—is 86,400 seconds of mean solar time (86,400.0025 SI seconds).[179] Because Earth's solar day is now slightly longer than it was during the 19th century due to tidal deceleration, each day varies between 0 and 2 SI ms longer.[180][181]
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Earth's rotation period relative to the fixed stars, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is 86,164.0989 seconds of mean solar time (UT1), or 23h 56m 4.0989s.[2][n 19] Earth's rotation period relative to the precessing or moving mean March equinox, misnamed its sidereal day, is 86,164.0905 seconds of mean solar time (UT1) (23h 56m 4.0905s).[2] Thus the sidereal day is shorter than the stellar day by about 8.4 ms.[182] The length of the mean solar day in SI seconds is available from the IERS for the periods 1623–2005[183] and 1962–2005.[184]
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Apart from meteors within the atmosphere and low-orbiting satellites, the main apparent motion of celestial bodies in Earth's sky is to the west at a rate of 15°/h = 15'/min. For bodies near the celestial equator, this is equivalent to an apparent diameter of the Sun or the Moon every two minutes; from Earth's surface, the apparent sizes of the Sun and the Moon are approximately the same.[185][186]
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Earth orbits the Sun at an average distance of about 150 million km (93 million mi) every 365.2564 mean solar days, or one sidereal year. This gives an apparent movement of the Sun eastward with respect to the stars at a rate of about 1°/day, which is one apparent Sun or Moon diameter every 12 hours. Due to this motion, on average it takes 24 hours—a solar day—for Earth to complete a full rotation about its axis so that the Sun returns to the meridian. The orbital speed of Earth averages about 29.78 km/s (107,200 km/h; 66,600 mph), which is fast enough to travel a distance equal to Earth's diameter, about 12,742 km (7,918 mi), in seven minutes, and the distance to the Moon, 384,000 km (239,000 mi), in about 3.5 hours.[3]
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The Moon and Earth orbit a common barycenter every 27.32 days relative to the background stars. When combined with the Earth–Moon system's common orbit around the Sun, the period of the synodic month, from new moon to new moon, is 29.53 days. Viewed from the celestial north pole, the motion of Earth, the Moon, and their axial rotations are all counterclockwise. Viewed from a vantage point above the north poles of both the Sun and Earth, Earth orbits in a counterclockwise direction about the Sun. The orbital and axial planes are not precisely aligned: Earth's axis is tilted some 23.44 degrees from the perpendicular to the Earth–Sun plane (the ecliptic), and the Earth–Moon plane is tilted up to ±5.1 degrees against the Earth–Sun plane. Without this tilt, there would be an eclipse every two weeks, alternating between lunar eclipses and solar eclipses.[3][188]
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The Hill sphere, or the sphere of gravitational influence, of Earth is about 1.5 million km (930,000 mi) in radius.[189][n 20] This is the maximum distance at which Earth's gravitational influence is stronger than the more distant Sun and planets. Objects must orbit Earth within this radius, or they can become unbound by the gravitational perturbation of the Sun.
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Earth, along with the Solar System, is situated in the Milky Way and orbits about 28,000 light-years from its center. It is about 20 light-years above the galactic plane in the Orion Arm.[190]
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The axial tilt of Earth is approximately 23.439281°[2] with the axis of its orbit plane, always pointing towards the Celestial Poles. Due to Earth's axial tilt, the amount of sunlight reaching any given point on the surface varies over the course of the year. This causes the seasonal change in climate, with summer in the Northern Hemisphere occurring when the Tropic of Cancer is facing the Sun, and winter taking place when the Tropic of Capricorn in the Southern Hemisphere faces the Sun. During the summer, the day lasts longer, and the Sun climbs higher in the sky. In winter, the climate becomes cooler and the days shorter. In northern temperate latitudes, the Sun rises north of true east during the summer solstice, and sets north of true west, reversing in the winter. The Sun rises south of true east in the summer for the southern temperate zone and sets south of true west.
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Above the Arctic Circle, an extreme case is reached where there is no daylight at all for part of the year, up to six months at the North Pole itself, a polar night. In the Southern Hemisphere, the situation is exactly reversed, with the South Pole oriented opposite the direction of the North Pole. Six months later, this pole will experience a midnight sun, a day of 24 hours, again reversing with the South Pole.
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By astronomical convention, the four seasons can be determined by the solstices—the points in the orbit of maximum axial tilt toward or away from the Sun—and the equinoxes, when Earth's rotational axis is aligned with its orbital axis. In the Northern Hemisphere, winter solstice currently occurs around 21 December; summer solstice is near 21 June, spring equinox is around 20 March and autumnal equinox is about 22 or 23 September. In the Southern Hemisphere, the situation is reversed, with the summer and winter solstices exchanged and the spring and autumnal equinox dates swapped.[191]
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The angle of Earth's axial tilt is relatively stable over long periods of time. Its axial tilt does undergo nutation; a slight, irregular motion with a main period of 18.6 years.[192] The orientation (rather than the angle) of Earth's axis also changes over time, precessing around in a complete circle over each 25,800 year cycle; this precession is the reason for the difference between a sidereal year and a tropical year. Both of these motions are caused by the varying attraction of the Sun and the Moon on Earth's equatorial bulge. The poles also migrate a few meters across Earth's surface. This polar motion has multiple, cyclical components, which collectively are termed quasiperiodic motion. In addition to an annual component to this motion, there is a 14-month cycle called the Chandler wobble. Earth's rotational velocity also varies in a phenomenon known as length-of-day variation.[193]
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In modern times, Earth's perihelion occurs around 3 January, and its aphelion around 4 July. These dates change over time due to precession and other orbital factors, which follow cyclical patterns known as Milankovitch cycles. The changing Earth–Sun distance causes an increase of about 6.9%[n 21] in solar energy reaching Earth at perihelion relative to aphelion. Because the Southern Hemisphere is tilted toward the Sun at about the same time that Earth reaches the closest approach to the Sun, the Southern Hemisphere receives slightly more energy from the Sun than does the northern over the course of a year. This effect is much less significant than the total energy change due to the axial tilt, and most of the excess energy is absorbed by the higher proportion of water in the Southern Hemisphere.[194]
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A study from 2016 suggested that Planet Nine tilted all the planets of the Solar System, including Earth, by about six degrees.[195]
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A planet that can sustain life is termed habitable, even if life did not originate there. Earth provides liquid water—an environment where complex organic molecules can assemble and interact, and sufficient energy to sustain metabolism.[196] The distance of Earth from the Sun, as well as its orbital eccentricity, rate of rotation, axial tilt, geological history, sustaining atmosphere, and magnetic field all contribute to the current climatic conditions at the surface.[197]
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A planet's life forms inhabit ecosystems, whose total is sometimes said to form a "biosphere".[198] Earth's biosphere is thought to have begun evolving about 3.5 Gya.[65] The biosphere is divided into a number of biomes, inhabited by broadly similar plants and animals.[199] On land, biomes are separated primarily by differences in latitude, height above sea level and humidity. Terrestrial biomes lying within the Arctic or Antarctic Circles, at high altitudes or in extremely arid areas are relatively barren of plant and animal life; species diversity reaches a peak in humid lowlands at equatorial latitudes.[200]
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In July 2016, scientists reported identifying a set of 355 genes from the last universal common ancestor (LUCA) of all organisms living on Earth.[201]
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Earth has resources that have been exploited by humans.[203] Those termed non-renewable resources, such as fossil fuels, only renew over geological timescales.[204]
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Large deposits of fossil fuels are obtained from Earth's crust, consisting of coal, petroleum, and natural gas.[205] These deposits are used by humans both for energy production and as feedstock for chemical production.[206] Mineral ore bodies have also been formed within the crust through a process of ore genesis, resulting from actions of magmatism, erosion, and plate tectonics.[207] These bodies form concentrated sources for many metals and other useful elements.
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Earth's biosphere produces many useful biological products for humans, including food, wood, pharmaceuticals, oxygen, and the recycling of many organic wastes. The land-based ecosystem depends upon topsoil and fresh water, and the oceanic ecosystem depends upon dissolved nutrients washed down from the land.[208] In 1980, 50.53 million km2 (19.51 million sq mi) of Earth's land surface consisted of forest and woodlands, 67.88 million km2 (26.21 million sq mi) was grasslands and pasture, and 15.01 million km2 (5.80 million sq mi) was cultivated as croplands.[209] The estimated amount of irrigated land in 1993 was 2,481,250 km2 (958,020 sq mi).[13] Humans also live on the land by using building materials to construct shelters.
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Large areas of Earth's surface are subject to extreme weather such as tropical cyclones, hurricanes, or typhoons that dominate life in those areas. From 1980 to 2000, these events caused an average of 11,800 human deaths per year.[210] Many places are subject to earthquakes, landslides, tsunamis, volcanic eruptions, tornadoes, sinkholes, blizzards, floods, droughts, wildfires, and other calamities and disasters.
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Many localized areas are subject to human-made pollution of the air and water, acid rain and toxic substances, loss of vegetation (overgrazing, deforestation, desertification), loss of wildlife, species extinction, soil degradation, soil depletion and erosion.
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There is a scientific consensus linking human activities to global warming due to industrial carbon dioxide emissions. This is predicted to produce changes such as the melting of glaciers and ice sheets, more extreme temperature ranges, significant changes in weather and a global rise in average sea levels.[211]
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Cartography, the study and practice of map-making, and geography, the study of the lands, features, inhabitants and phenomena on Earth, have historically been the disciplines devoted to depicting Earth. Surveying, the determination of locations and distances, and to a lesser extent navigation, the determination of position and direction, have developed alongside cartography and geography, providing and suitably quantifying the requisite information.
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Earth's human population reached approximately seven billion on 31 October 2011.[213] Projections indicate that the world's human population will reach 9.2 billion in 2050.[214] Most of the growth is expected to take place in developing nations. Human population density varies widely around the world, but a majority live in Asia. By 2020, 60% of the world's population is expected to be living in urban, rather than rural, areas.[215]
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68% of the land mass of the world is in the northern hemisphere.[216] Partly due to the predominance of land mass, 90% of humans live in the northern hemisphere.[217]
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It is estimated that one-eighth of Earth's surface is suitable for humans to live on – three-quarters of Earth's surface is covered by oceans, leaving one-quarter as land. Half of that land area is desert (14%),[218] high mountains (27%),[219] or other unsuitable terrains. The northernmost permanent settlement in the world is Alert, on Ellesmere Island in Nunavut, Canada.[220] (82°28′N) The southernmost is the Amundsen–Scott South Pole Station, in Antarctica, almost exactly at the South Pole. (90°S)
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States claim the planet's entire land surface, except for parts of Antarctica and a few other unclaimed areas. Earth has never had a planetwide government, but the United Nations is the leading worldwide intergovernmental organization.
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The first human to orbit Earth was Yuri Gagarin on 12 April 1961.[221] In total, about 487 people have visited outer space and reached orbit as of 30 July 2010[update], and, of these, twelve have walked on the Moon.[222][223][224] Normally, the only humans in space are those on the International Space Station. The station's crew, made up of six people, is usually replaced every six months.[225] The farthest that humans have traveled from Earth is 400,171 km (248,655 mi), achieved during the Apollo 13 mission in 1970.[226]
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The Moon is a relatively large, terrestrial, planet-like natural satellite, with a diameter about one-quarter of Earth's. It is the largest moon in the Solar System relative to the size of its planet, although Charon is larger relative to the dwarf planet Pluto. The natural satellites of other planets are also referred to as "moons", after Earth's.
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The gravitational attraction between Earth and the Moon causes tides on Earth. The same effect on the Moon has led to its tidal locking: its rotation period is the same as the time it takes to orbit Earth. As a result, it always presents the same face to the planet. As the Moon orbits Earth, different parts of its face are illuminated by the Sun, leading to the lunar phases; the dark part of the face is separated from the light part by the solar terminator.
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Due to their tidal interaction, the Moon recedes from Earth at the rate of approximately 38 mm/a (1.5 in/year). Over millions of years, these tiny modifications—and the lengthening of Earth's day by about 23 µs/yr—add up to significant changes.[227] During the Devonian period, for example, (approximately 410 Mya) there were 400 days in a year, with each day lasting 21.8 hours.[228]
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The Moon may have dramatically affected the development of life by moderating the planet's climate. Paleontological evidence and computer simulations show that Earth's axial tilt is stabilized by tidal interactions with the Moon.[229] Some theorists think that without this stabilization against the torques applied by the Sun and planets to Earth's equatorial bulge, the rotational axis might be chaotically unstable, exhibiting chaotic changes over millions of years, as appears to be the case for Mars.[230]
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Viewed from Earth, the Moon is just far enough away to have almost the same apparent-sized disk as the Sun. The angular size (or solid angle) of these two bodies match because, although the Sun's diameter is about 400 times as large as the Moon's, it is also 400 times more distant.[186] This allows total and annular solar eclipses to occur on Earth.
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The most widely accepted theory of the Moon's origin, the giant-impact hypothesis, states that it formed from the collision of a Mars-size protoplanet called Theia with the early Earth. This hypothesis explains (among other things) the Moon's relative lack of iron and volatile elements and the fact that its composition is nearly identical to that of Earth's crust.[48]
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Earth has at least five co-orbital asteroids, including 3753 Cruithne and 2002 AA29.[231][232] A trojan asteroid companion, 2010 TK7, is librating around the leading Lagrange triangular point, L4, in Earth's orbit around the Sun.[233][234]
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The tiny near-Earth asteroid 2006 RH120 makes close approaches to the Earth–Moon system roughly every twenty years. During these approaches, it can orbit Earth for brief periods of time.[235]
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As of April 2018[update], there are 1,886 operational, human-made satellites orbiting Earth.[5] There are also inoperative satellites, including Vanguard 1, the oldest satellite currently in orbit, and over 16,000 pieces of tracked space debris.[n 3] Earth's largest artificial satellite is the International Space Station.
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The standard astronomical symbol of Earth consists of a cross circumscribed by a circle, ,[236] representing the four corners of the world.
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Human cultures have developed many views of the planet.[237] Earth is sometimes personified as a deity. In many cultures it is a mother goddess that is also the primary fertility deity,[238] and by the mid-20th century, the Gaia Principle compared Earth's environments and life as a single self-regulating organism leading to broad stabilization of the conditions of habitability.[239][240][241] Creation myths in many religions involve the creation of Earth by a supernatural deity or deities.[238]
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The Hindu Vedas (1500–900 BC) refer to the Earth as Bhūgola (भूगोल), which comes from Bhū (earth, ground) and Gola (ball, sphere, globe). It means the "globe of earth". There is no direct evidence that the Hindus of that time knew that the Earth was sphere-shaped, but this name has been used extensively since the inception of the Vedas.[citation needed]
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Scientific investigation has resulted in several culturally transformative shifts in people's view of the planet. Initial belief in a flat Earth was gradually displaced in the Greek colonies of southern Italy during the late 6th century BC by the idea of spherical Earth,[242][243][244] which was attributed to both the philosophers Pythagoras and Parmenides.[243][244] By the end of the 5th century BC, the sphericity of Earth was universally accepted among Greek intellectuals.[245] Earth was generally believed to be the center of the universe until the 16th century, when scientists first conclusively demonstrated that it was a moving object, comparable to the other planets in the Solar System.[246] Due to the efforts of influential Christian scholars and clerics such as James Ussher, who sought to determine the age of Earth through analysis of genealogies in Scripture, Westerners before the 19th century generally believed Earth to be a few thousand years old at most. It was only during the 19th century that geologists realized Earth's age was at least many millions of years.[247]
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Lord Kelvin used thermodynamics to estimate the age of Earth to be between 20 million and 400 million years in 1864, sparking a vigorous debate on the subject; it was only when radioactivity and radioactive dating were discovered in the late 19th and early 20th centuries that a reliable mechanism for determining Earth's age was established, proving the planet to be billions of years old.[248][249] The perception of Earth shifted again[further explanation needed] in the 20th century when humans first viewed it from orbit, and especially with photographs of Earth returned by the Apollo program.[250][251][252]
|
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Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
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1 |
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Earth is the third planet from the Sun and the only astronomical object known to harbor life. According to radiometric dating estimation and other evidence, Earth formed over 4.5 billion years ago. Earth's gravity interacts with other objects in space, especially the Sun and the Moon, which is Earth's only natural satellite. Earth orbits around the Sun in 365.256 solar days, a period known as an Earth sidereal year. During this time, Earth rotates about its axis 366.256 times, that is, a sidereal year has 366.256 sidereal days.[n 6]
|
6 |
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|
7 |
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Earth's axis of rotation is tilted with respect to its orbital plane, producing seasons on Earth. The gravitational interaction between Earth and the Moon causes tides, stabilizes Earth's orientation on its axis, and gradually slows its rotation. Earth is the densest planet in the Solar System and the largest and most massive of the four rocky planets.
|
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|
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Earth's outer layer (lithosphere) is divided into several rigid tectonic plates that migrate across the surface over many millions of years. About 29% of Earth's surface is land consisting of continents and islands. The remaining 71% is covered with water, mostly by oceans but also lakes, rivers and other fresh water, which all together constitute the hydrosphere. The majority of Earth's polar regions are covered in ice, including the Antarctic ice sheet and the sea ice of the Arctic ice pack. Earth's interior remains active with a solid iron inner core, a liquid outer core that generates Earth's magnetic field, and a convecting mantle that drives plate tectonics.
|
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|
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Within the first billion years of Earth's history, life appeared in the oceans and began to affect Earth's atmosphere and surface, leading to the proliferation of anaerobic and, later, aerobic organisms. Some geological evidence indicates that life may have arisen as early as 4.1 billion years ago. Since then, the combination of Earth's distance from the Sun, physical properties and geological history have allowed life to evolve and thrive. In the history of life on Earth, biodiversity has gone through long periods of expansion, occasionally punctuated by mass extinctions. Over 99% of all species that ever lived on Earth are extinct. Estimates of the number of species on Earth today vary widely; most species have not been described. Over 7.7 billion humans live on Earth and depend on its biosphere and natural resources for their survival.[23]
|
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|
13 |
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The modern English word Earth developed, via Middle English,[n 7] from an Old English noun most often spelled eorðe.[24] It has cognates in every Germanic language, and their ancestral root has been reconstructed as *erþō. In its earliest attestation, the word eorðe was already being used to translate the many senses of Latin terra and Greek γῆ gē: the ground,[n 8] its soil,[n 9] dry land,[n 10] the human world,[n 11] the surface of the world (including the sea),[n 12] and the globe itself.[n 13] As with Roman Terra/Tellūs and Greek Gaia, Earth may have been a personified goddess in Germanic paganism: late Norse mythology included Jörð ('Earth'), a giantess often given as the mother of Thor.[33]
|
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+
|
15 |
+
Originally, earth was written in lowercase, and from early Middle English, its definite sense as "the globe" was expressed as the earth. By Early Modern English, many nouns were capitalized, and the earth became (and often remained) the Earth, particularly when referenced along with other heavenly bodies. More recently, the name is sometimes simply given as Earth, by analogy with the names of the other planets.[24] House styles now vary: Oxford spelling recognizes the lowercase form as the most common, with the capitalized form an acceptable variant. Another convention capitalizes "Earth" when appearing as a name (e.g. "Earth's atmosphere") but writes it in lowercase when preceded by the (e.g. "the atmosphere of the earth"). It almost always appears in lowercase in colloquial expressions such as "what on earth are you doing?"[34]
|
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|
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+
Occasionally, the name Terra /ˈtɛrə/ is used in scientific writing and especially in science fiction to distinguish our inhabited planet from others,[35] while in poetry Tellus /ˈtɛləs/ has been used to denote personification of the Earth.[36] The Greek poetic name Gaea (Gæa) /ˈdʒiːə/ is rare, though the alternative spelling Gaia has become common due to the Gaia hypothesis, in which case its pronunciation is /ˈɡaɪə/ rather than the more Classical /ˈɡeɪə/.[37]
|
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+
|
19 |
+
There are a number of adjectives for the planet Earth. From Earth itself comes earthly. From Latin Terra come Terran /ˈtɛrən/,[38] Terrestrial /təˈrɛstriəl/,[39] and (via French) Terrene /təˈriːn/,[40] and from Latin Tellus come Tellurian /tɛˈlʊəriən/[41] and, more rarely, Telluric and Tellural. From Greek Gaia and Gaea comes Gaian and Gaean.
|
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|
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+
An inhabitant of the Earth is an Earthling, a Terran, a Terrestrial, a Tellurian or, rarely, an Earthian.
|
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+
|
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+
The oldest material found in the Solar System is dated to 4.5672±0.0006 billion years ago (BYA).[42] By 4.54±0.04 BYA[43] the primordial Earth had formed. The bodies in the Solar System formed and evolved with the Sun. In theory, a solar nebula partitions a volume out of a molecular cloud by gravitational collapse, which begins to spin and flatten into a circumstellar disk, and then the planets grow out of that disk with the Sun. A nebula contains gas, ice grains, and dust (including primordial nuclides). According to nebular theory, planetesimals formed by accretion, with the primordial Earth taking 10–20 million years (Mys) to form.[44]
|
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+
|
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+
A subject of research is the formation of the Moon, some 4.53 BYA.[45] A leading hypothesis is that it was formed by accretion from material loosed from Earth after a Mars-sized object, named Theia, hit Earth.[46] In this view, the mass of Theia was approximately 10 percent of Earth;[47] it hit Earth with a glancing blow and some of its mass merged with Earth.[48] Between approximately 4.1 and 3.8 BYA, numerous asteroid impacts during the Late Heavy Bombardment caused significant changes to the greater surface environment of the Moon and, by inference, to that of Earth.
|
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+
|
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+
Earth's atmosphere and oceans were formed by volcanic activity and outgassing.[49] Water vapor from these sources condensed into the oceans, augmented by water and ice from asteroids, protoplanets, and comets.[50] In this model, atmospheric "greenhouse gases" kept the oceans from freezing when the newly forming Sun had only 70% of its current luminosity.[51] By 3.5 BYA, Earth's magnetic field was established, which helped prevent the atmosphere from being stripped away by the solar wind.[52]
|
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+
A crust formed when the molten outer layer of Earth cooled to form a solid. The two models[53] that explain land mass propose either a steady growth to the present-day forms[54] or, more likely, a rapid growth[55] early in Earth history[56] followed by a long-term steady continental area.[57][58][59] Continents formed by plate tectonics, a process ultimately driven by the continuous loss of heat from Earth's interior. Over the period of hundreds of millions of years, the supercontinents have assembled and broken apart. Roughly 750 million years ago (MYA), one of the earliest known supercontinents, Rodinia, began to break apart. The continents later recombined to form Pannotia 600–540 MYA, then finally Pangaea, which also broke apart 180 MYA.[60]
|
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|
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The present pattern of ice ages began about 40 MYA,[61] and then intensified during the Pleistocene about 3 MYA.[62] High-latitude regions have since undergone repeated cycles of glaciation and thaw, repeating about every 40,000–100,000 years. The last continental glaciation ended 10,000 years ago.[63]
|
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|
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+
Chemical reactions led to the first self-replicating molecules about four billion years ago. A half billion years later, the last common ancestor of all current life arose.[64] The evolution of photosynthesis allowed the Sun's energy to be harvested directly by life forms. The resultant molecular oxygen (O2) accumulated in the atmosphere and due to interaction with ultraviolet solar radiation, formed a protective ozone layer (O3) in the upper atmosphere.[65] The incorporation of smaller cells within larger ones resulted in the development of complex cells called eukaryotes.[66] True multicellular organisms formed as cells within colonies became increasingly specialized. Aided by the absorption of harmful ultraviolet radiation by the ozone layer, life colonized Earth's surface.[67] Among the earliest fossil evidence for life is microbial mat fossils found in 3.48 billion-year-old sandstone in Western Australia,[68] biogenic graphite found in 3.7 billion-year-old metasedimentary rocks in Western Greenland,[69] and remains of biotic material found in 4.1 billion-year-old rocks in Western Australia.[70][71] The earliest direct evidence of life on Earth is contained in 3.45 billion-year-old Australian rocks showing fossils of microorganisms.[72][73]
|
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+
|
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+
During the Neoproterozoic, 750 to 580 MYA, much of Earth might have been covered in ice. This hypothesis has been termed "Snowball Earth", and it is of particular interest because it preceded the Cambrian explosion, when multicellular life forms significantly increased in complexity.[74] Following the Cambrian explosion, 535 MYA, there have been five mass extinctions.[75] The most recent such event was 66 MYA, when an asteroid impact triggered the extinction of the non-avian dinosaurs and other large reptiles, but spared some small animals such as mammals, which at the time resembled shrews. Mammalian life has diversified over the past 66 Mys, and several million years ago an African ape-like animal such as Orrorin tugenensis gained the ability to stand upright.[76] This facilitated tool use and encouraged communication that provided the nutrition and stimulation needed for a larger brain, which led to the evolution of humans. The development of agriculture, and then civilization, led to humans having an influence on Earth and the nature and quantity of other life forms that continues to this day.[77]
|
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|
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+
Earth's expected long-term future is tied to that of the Sun. Over the next 1.1 billion years, solar luminosity will increase by 10%, and over the next 3.5 billion years by 40%.[78] Earth's increasing surface temperature will accelerate the inorganic carbon cycle, reducing CO2 concentration to levels lethally low for plants (10 ppm for C4 photosynthesis) in approximately 100–900 million years.[79][80] The lack of vegetation will result in the loss of oxygen in the atmosphere, making animal life impossible.[81] About a billion years from now, all surface water will have disappeared[82] and the mean global temperature will reach 70 °C (158 °F).[81] Earth is expected to be habitable until the end of photosynthesis about 500 million years from now,[79] but if nitrogen is removed from the atmosphere, life may continue until a runaway greenhouse effect occurs 2.3 billion years from now.[80] Anthropogenic emissions are "probably insufficient" to cause a runaway greenhouse at current solar luminosity.[83] Even if the Sun were eternal and stable, 27% of the water in the modern oceans will descend to the mantle in one billion years, due to reduced steam venting from mid-ocean ridges.[84]
|
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+
|
39 |
+
The Sun will evolve to become a red giant in about 5 billion years. Models predict that the Sun will expand to roughly 1 AU (150 million km; 93 million mi), about 250 times its present radius.[78][85] Earth's fate is less clear. As a red giant, the Sun will lose roughly 30% of its mass, so, without tidal effects, Earth will move to an orbit 1.7 AU (250 million km; 160 million mi) from the Sun when the star reaches its maximum radius. Most, if not all, remaining life will be destroyed by the Sun's increased luminosity (peaking at about 5,000 times its present level).[78] A 2008 simulation indicates that Earth's orbit will eventually decay due to tidal effects and drag, causing it to enter the Sun's atmosphere and be vaporized.[85]
|
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+
|
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+
The shape of Earth is nearly spherical. There is a small flattening at the poles and bulging around the equator due to Earth's rotation.[89] To second order, Earth is approximately an oblate spheroid, whose equatorial diameter is 43 kilometres (27 mi) larger than the pole-to-pole diameter,[90] although the variation is less than 1% of the average radius of the Earth.
|
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+
|
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+
The point on the surface farthest from Earth's center of mass is the summit of the equatorial Chimborazo volcano in Ecuador (6,384.4 km or 3,967.1 mi).[91][92][93][94] The average diameter of the reference spheroid is 12,742 kilometres (7,918 mi). Local topography deviates from this idealized spheroid, although on a global scale these deviations are small compared to Earth's radius: the maximum deviation of only 0.17% is at the Mariana Trench (10,911 metres or 35,797 feet below local sea level), whereas Mount Everest (8,848 metres or 29,029 feet above local sea level) represents a deviation of 0.14%.[n 14]
|
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+
|
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+
In geodesy, the exact shape that Earth's oceans would adopt in the absence of land and perturbations such as tides and winds is called the geoid. More precisely, the geoid is the surface of gravitational equipotential at mean sea level.
|
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+
|
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+
Earth's mass is approximately 5.97×1024 kg (5,970 Yg). It is composed mostly of iron (32.1%), oxygen (30.1%), silicon (15.1%), magnesium (13.9%), sulphur (2.9%), nickel (1.8%), calcium (1.5%), and aluminum (1.4%), with the remaining 1.2% consisting of trace amounts of other elements. Due to mass segregation, the core region is estimated to be primarily composed of iron (88.8%), with smaller amounts of nickel (5.8%), sulphur (4.5%), and less than 1% trace elements.[98]
|
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|
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The most common rock constituents of the crust are nearly all oxides: chlorine, sulphur, and fluorine are the important exceptions to this and their total amount in any rock is usually much less than 1%. Over 99% of the crust is composed of 11 oxides, principally silica, alumina, iron oxides, lime, magnesia, potash and soda.[99][98][100]
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Earth's interior, like that of the other terrestrial planets, is divided into layers by their chemical or physical (rheological) properties. The outer layer is a chemically distinct silicate solid crust, which is underlain by a highly viscous solid mantle. The crust is separated from the mantle by the Mohorovičić discontinuity. The thickness of the crust varies from about 6 kilometres (3.7 mi) under the oceans to 30–50 km (19–31 mi) for the continents. The crust and the cold, rigid, top of the upper mantle are collectively known as the lithosphere, and it is of the lithosphere that the tectonic plates are composed. Beneath the lithosphere is the asthenosphere, a relatively low-viscosity layer on which the lithosphere rides. Important changes in crystal structure within the mantle occur at 410 and 660 km (250 and 410 mi) below the surface, spanning a transition zone that separates the upper and lower mantle. Beneath the mantle, an extremely low viscosity liquid outer core lies above a solid inner core.[101] Earth's inner core might rotate at a slightly higher angular velocity than the remainder of the planet, advancing by 0.1–0.5° per year.[102] The radius of the inner core is about one fifth of that of Earth.
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Earth's internal heat comes from a combination of residual heat from planetary accretion (about 20%) and heat produced through radioactive decay (80%).[105] The major heat-producing isotopes within Earth are potassium-40, uranium-238, and thorium-232.[106] At the center, the temperature may be up to 6,000 °C (10,830 °F),[107] and the pressure could reach 360 GPa (52 million psi).[108] Because much of the heat is provided by radioactive decay, scientists postulate that early in Earth's history, before isotopes with short half-lives were depleted, Earth's heat production was much higher. At approximately 3 Gyr, twice the present-day heat would have been produced, increasing the rates of mantle convection and plate tectonics, and allowing the production of uncommon igneous rocks such as komatiites that are rarely formed today.[105][109]
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The mean heat loss from Earth is 87 mW m−2, for a global heat loss of 4.42×1013 W.[111] A portion of the core's thermal energy is transported toward the crust by mantle plumes, a form of convection consisting of upwellings of higher-temperature rock. These plumes can produce hotspots and flood basalts.[112] More of the heat in Earth is lost through plate tectonics, by mantle upwelling associated with mid-ocean ridges. The final major mode of heat loss is through conduction through the lithosphere, the majority of which occurs under the oceans because the crust there is much thinner than that of the continents.[113]
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Earth's mechanically rigid outer layer, the lithosphere, is divided into tectonic plates. These plates are rigid segments that move relative to each other at one of three boundaries types: At convergent boundaries, two plates come together; at divergent boundaries, two plates are pulled apart; and at transform boundaries, two plates slide past one another laterally. Along these plate boundaries, earthquakes, volcanic activity, mountain-building, and oceanic trench formation can occur.[115] The tectonic plates ride on top of the asthenosphere, the solid but less-viscous part of the upper mantle that can flow and move along with the plates.[116]
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As the tectonic plates migrate, oceanic crust is subducted under the leading edges of the plates at convergent boundaries. At the same time, the upwelling of mantle material at divergent boundaries creates mid-ocean ridges. The combination of these processes recycles the oceanic crust back into the mantle. Due to this recycling, most of the ocean floor is less than 100 Ma old. The oldest oceanic crust is located in the Western Pacific and is estimated to be 200 Ma old.[117][118] By comparison, the oldest dated continental crust is 4,030 Ma.[119]
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The seven major plates are the Pacific, North American, Eurasian, African, Antarctic, Indo-Australian, and South American. Other notable plates include the Arabian Plate, the Caribbean Plate, the Nazca Plate off the west coast of South America and the Scotia Plate in the southern Atlantic Ocean. The Australian Plate fused with the Indian Plate between 50 and 55 Mya. The fastest-moving plates are the oceanic plates, with the Cocos Plate advancing at a rate of 75 mm/a (3.0 in/year)[120] and the Pacific Plate moving 52–69 mm/a (2.0–2.7 in/year). At the other extreme, the slowest-moving plate is the Eurasian Plate, progressing at a typical rate of 21 mm/a (0.83 in/year).[121]
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The total surface area of Earth is about 510 million km2 (197 million sq mi).[12] Of this, 70.8%,[12] or 361.13 million km2 (139.43 million sq mi), is below sea level and covered by ocean water.[122] Below the ocean's surface are much of the continental shelf, mountains, volcanoes,[90] oceanic trenches, submarine canyons, oceanic plateaus, abyssal plains, and a globe-spanning mid-ocean ridge system. The remaining 29.2%, or 148.94 million km2 (57.51 million sq mi), not covered by water has terrain that varies greatly from place to place and consists of mountains, deserts, plains, plateaus, and other landforms. Tectonics and erosion, volcanic eruptions, flooding, weathering, glaciation, the growth of coral reefs, and meteorite impacts are among the processes that constantly reshape Earth's surface over geological time.[123][124]
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The continental crust consists of lower density material such as the igneous rocks granite and andesite. Less common is basalt, a denser volcanic rock that is the primary constituent of the ocean floors.[125] Sedimentary rock is formed from the accumulation of sediment that becomes buried and compacted together. Nearly 75% of the continental surfaces are covered by sedimentary rocks, although they form about 5% of the crust.[126] The third form of rock material found on Earth is metamorphic rock, which is created from the transformation of pre-existing rock types through high pressures, high temperatures, or both. The most abundant silicate minerals on Earth's surface include quartz, feldspars, amphibole, mica, pyroxene and olivine.[127] Common carbonate minerals include calcite (found in limestone) and dolomite.[128]
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The elevation of the land surface varies from the low point of −418 m (−1,371 ft) at the Dead Sea, to a maximum altitude of 8,848 m (29,029 ft) at the top of Mount Everest. The mean height of land above sea level is about 797 m (2,615 ft).[129]
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The pedosphere is the outermost layer of Earth's continental surface and is composed of soil and subject to soil formation processes. The total arable land is 10.9% of the land surface, with 1.3% being permanent cropland.[130][131] Close to 40% of Earth's land surface is used for agriculture, or an estimated 16.7 million km2 (6.4 million sq mi) of cropland and 33.5 million km2 (12.9 million sq mi) of pastureland.[132]
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The abundance of water on Earth's surface is a unique feature that distinguishes the "Blue Planet" from other planets in the Solar System. Earth's hydrosphere consists chiefly of the oceans, but technically includes all water surfaces in the world, including inland seas, lakes, rivers, and underground waters down to a depth of 2,000 m (6,600 ft). The deepest underwater location is Challenger Deep of the Mariana Trench in the Pacific Ocean with a depth of 10,911.4 m (35,799 ft).[n 18][133]
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The mass of the oceans is approximately 1.35×1018 metric tons or about 1/4400 of Earth's total mass. The oceans cover an area of 361.8 million km2 (139.7 million sq mi) with a mean depth of 3,682 m (12,080 ft), resulting in an estimated volume of 1.332 billion km3 (320 million cu mi).[134] If all of Earth's crustal surface were at the same elevation as a smooth sphere, the depth of the resulting world ocean would be 2.7 to 2.8 km (1.68 to 1.74 mi).[135][136]
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About 97.5% of the water is saline; the remaining 2.5% is fresh water. Most fresh water, about 68.7%, is present as ice in ice caps and glaciers.[137]
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The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5% salt).[138] Most of this salt was released from volcanic activity or extracted from cool igneous rocks.[139] The oceans are also a reservoir of dissolved atmospheric gases, which are essential for the survival of many aquatic life forms.[140] Sea water has an important influence on the world's climate, with the oceans acting as a large heat reservoir.[141] Shifts in the oceanic temperature distribution can cause significant weather shifts, such as the El Niño–Southern Oscillation.[142]
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The atmospheric pressure at Earth's sea level averages 101.325 kPa (14.696 psi),[143] with a scale height of about 8.5 km (5.3 mi).[3] A dry atmosphere is composed of 78.084% nitrogen, 20.946% oxygen, 0.934% argon, and trace amounts of carbon dioxide and other gaseous molecules.[143] Water vapor content varies between 0.01% and 4%[143] but averages about 1%.[3] The height of the troposphere varies with latitude, ranging between 8 km (5 mi) at the poles to 17 km (11 mi) at the equator, with some variation resulting from weather and seasonal factors.[144]
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Earth's biosphere has significantly altered its atmosphere. Oxygenic photosynthesis evolved 2.7 Gya, forming the primarily nitrogen–oxygen atmosphere of today.[65] This change enabled the proliferation of aerobic organisms and, indirectly, the formation of the ozone layer due to the subsequent conversion of atmospheric O2 into O3. The ozone layer blocks ultraviolet solar radiation, permitting life on land.[145] Other atmospheric functions important to life include transporting water vapor, providing useful gases, causing small meteors to burn up before they strike the surface, and moderating temperature.[146] This last phenomenon is known as the greenhouse effect: trace molecules within the atmosphere serve to capture thermal energy emitted from the ground, thereby raising the average temperature. Water vapor, carbon dioxide, methane, nitrous oxide, and ozone are the primary greenhouse gases in the atmosphere. Without this heat-retention effect, the average surface temperature would be −18 °C (0 °F), in contrast to the current +15 °C (59 °F),[147] and life on Earth probably would not exist in its current form.[148] In May 2017, glints of light, seen as twinkling from an orbiting satellite a million miles away, were found to be reflected light from ice crystals in the atmosphere.[149][150]
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Earth's atmosphere has no definite boundary, slowly becoming thinner and fading into outer space. Three-quarters of the atmosphere's mass is contained within the first 11 km (6.8 mi) of the surface. This lowest layer is called the troposphere. Energy from the Sun heats this layer, and the surface below, causing expansion of the air. This lower-density air then rises and is replaced by cooler, higher-density air. The result is atmospheric circulation that drives the weather and climate through redistribution of thermal energy.[151]
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The primary atmospheric circulation bands consist of the trade winds in the equatorial region below 30° latitude and the westerlies in the mid-latitudes between 30° and 60°.[152] Ocean currents are also important factors in determining climate, particularly the thermohaline circulation that distributes thermal energy from the equatorial oceans to the polar regions.[153]
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Water vapor generated through surface evaporation is transported by circulatory patterns in the atmosphere. When atmospheric conditions permit an uplift of warm, humid air, this water condenses and falls to the surface as precipitation.[151] Most of the water is then transported to lower elevations by river systems and usually returned to the oceans or deposited into lakes. This water cycle is a vital mechanism for supporting life on land and is a primary factor in the erosion of surface features over geological periods. Precipitation patterns vary widely, ranging from several meters of water per year to less than a millimeter. Atmospheric circulation, topographic features, and temperature differences determine the average precipitation that falls in each region.[154]
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The amount of solar energy reaching Earth's surface decreases with increasing latitude. At higher latitudes, the sunlight reaches the surface at lower angles, and it must pass through thicker columns of the atmosphere. As a result, the mean annual air temperature at sea level decreases by about 0.4 °C (0.7 °F) per degree of latitude from the equator.[155] Earth's surface can be subdivided into specific latitudinal belts of approximately homogeneous climate. Ranging from the equator to the polar regions, these are the tropical (or equatorial), subtropical, temperate and polar climates.[156]
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This latitudinal rule has several anomalies:
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The commonly used Köppen climate classification system has five broad groups (humid tropics, arid, humid middle latitudes, continental and cold polar), which are further divided into more specific subtypes.[152] The Köppen system rates regions of terrain based on observed temperature and precipitation.
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The highest air temperature ever measured on Earth was 56.7 °C (134.1 °F) in Furnace Creek, California, in Death Valley, in 1913.[159] The lowest air temperature ever directly measured on Earth was −89.2 °C (−128.6 °F) at Vostok Station in 1983,[160] but satellites have used remote sensing to measure temperatures as low as −94.7 °C (−138.5 °F) in East Antarctica.[161] These temperature records are only measurements made with modern instruments from the 20th century onwards and likely do not reflect the full range of temperature on Earth.
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Above the troposphere, the atmosphere is usually divided into the stratosphere, mesosphere, and thermosphere.[146] Each layer has a different lapse rate, defining the rate of change in temperature with height. Beyond these, the exosphere thins out into the magnetosphere, where the geomagnetic fields interact with the solar wind.[162] Within the stratosphere is the ozone layer, a component that partially shields the surface from ultraviolet light and thus is important for life on Earth. The Kármán line, defined as 100 km above Earth's surface, is a working definition for the boundary between the atmosphere and outer space.[163]
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Thermal energy causes some of the molecules at the outer edge of the atmosphere to increase their velocity to the point where they can escape from Earth's gravity. This causes a slow but steady loss of the atmosphere into space. Because unfixed hydrogen has a low molecular mass, it can achieve escape velocity more readily, and it leaks into outer space at a greater rate than other gases.[164] The leakage of hydrogen into space contributes to the shifting of Earth's atmosphere and surface from an initially reducing state to its current oxidizing one. Photosynthesis provided a source of free oxygen, but the loss of reducing agents such as hydrogen is thought to have been a necessary precondition for the widespread accumulation of oxygen in the atmosphere.[165] Hence the ability of hydrogen to escape from the atmosphere may have influenced the nature of life that developed on Earth.[166] In the current, oxygen-rich atmosphere most hydrogen is converted into water before it has an opportunity to escape. Instead, most of the hydrogen loss comes from the destruction of methane in the upper atmosphere.[167]
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The gravity of Earth is the acceleration that is imparted to objects due to the distribution of mass within Earth. Near Earth's surface, gravitational acceleration is approximately 9.8 m/s2 (32 ft/s2). Local differences in topography, geology, and deeper tectonic structure cause local and broad, regional differences in Earth's gravitational field, known as gravity anomalies.[168]
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The main part of Earth's magnetic field is generated in the core, the site of a dynamo process that converts the kinetic energy of thermally and compositionally driven convection into electrical and magnetic field energy. The field extends outwards from the core, through the mantle, and up to Earth's surface, where it is, approximately, a dipole. The poles of the dipole are located close to Earth's geographic poles. At the equator of the magnetic field, the magnetic-field strength at the surface is 3.05×10−5 T, with a magnetic dipole moment of 7.79×1022 Am2 at epoch 2000, decreasing nearly 6% per century.[169] The convection movements in the core are chaotic; the magnetic poles drift and periodically change alignment. This causes secular variation of the main field and field reversals at irregular intervals averaging a few times every million years. The most recent reversal occurred approximately 700,000 years ago.[170][171]
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The extent of Earth's magnetic field in space defines the magnetosphere. Ions and electrons of the solar wind are deflected by the magnetosphere; solar wind pressure compresses the dayside of the magnetosphere, to about 10 Earth radii, and extends the nightside magnetosphere into a long tail.[172] Because the velocity of the solar wind is greater than the speed at which waves propagate through the solar wind, a supersonic bow shock precedes the dayside magnetosphere within the solar wind.[173] Charged particles are contained within the magnetosphere; the plasmasphere is defined by low-energy particles that essentially follow magnetic field lines as Earth rotates;[174][175] the ring current is defined by medium-energy particles that drift relative to the geomagnetic field, but with paths that are still dominated by the magnetic field,[176] and the Van Allen radiation belt are formed by high-energy particles whose motion is essentially random, but otherwise contained by the magnetosphere.[172][177]
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During magnetic storms and substorms, charged particles can be deflected from the outer magnetosphere and especially the magnetotail, directed along field lines into Earth's ionosphere, where atmospheric atoms can be excited and ionized, causing the aurora.[178]
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Earth's rotation period relative to the Sun—its mean solar day—is 86,400 seconds of mean solar time (86,400.0025 SI seconds).[179] Because Earth's solar day is now slightly longer than it was during the 19th century due to tidal deceleration, each day varies between 0 and 2 SI ms longer.[180][181]
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Earth's rotation period relative to the fixed stars, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is 86,164.0989 seconds of mean solar time (UT1), or 23h 56m 4.0989s.[2][n 19] Earth's rotation period relative to the precessing or moving mean March equinox, misnamed its sidereal day, is 86,164.0905 seconds of mean solar time (UT1) (23h 56m 4.0905s).[2] Thus the sidereal day is shorter than the stellar day by about 8.4 ms.[182] The length of the mean solar day in SI seconds is available from the IERS for the periods 1623–2005[183] and 1962–2005.[184]
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Apart from meteors within the atmosphere and low-orbiting satellites, the main apparent motion of celestial bodies in Earth's sky is to the west at a rate of 15°/h = 15'/min. For bodies near the celestial equator, this is equivalent to an apparent diameter of the Sun or the Moon every two minutes; from Earth's surface, the apparent sizes of the Sun and the Moon are approximately the same.[185][186]
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Earth orbits the Sun at an average distance of about 150 million km (93 million mi) every 365.2564 mean solar days, or one sidereal year. This gives an apparent movement of the Sun eastward with respect to the stars at a rate of about 1°/day, which is one apparent Sun or Moon diameter every 12 hours. Due to this motion, on average it takes 24 hours—a solar day—for Earth to complete a full rotation about its axis so that the Sun returns to the meridian. The orbital speed of Earth averages about 29.78 km/s (107,200 km/h; 66,600 mph), which is fast enough to travel a distance equal to Earth's diameter, about 12,742 km (7,918 mi), in seven minutes, and the distance to the Moon, 384,000 km (239,000 mi), in about 3.5 hours.[3]
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The Moon and Earth orbit a common barycenter every 27.32 days relative to the background stars. When combined with the Earth–Moon system's common orbit around the Sun, the period of the synodic month, from new moon to new moon, is 29.53 days. Viewed from the celestial north pole, the motion of Earth, the Moon, and their axial rotations are all counterclockwise. Viewed from a vantage point above the north poles of both the Sun and Earth, Earth orbits in a counterclockwise direction about the Sun. The orbital and axial planes are not precisely aligned: Earth's axis is tilted some 23.44 degrees from the perpendicular to the Earth–Sun plane (the ecliptic), and the Earth–Moon plane is tilted up to ±5.1 degrees against the Earth–Sun plane. Without this tilt, there would be an eclipse every two weeks, alternating between lunar eclipses and solar eclipses.[3][188]
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The Hill sphere, or the sphere of gravitational influence, of Earth is about 1.5 million km (930,000 mi) in radius.[189][n 20] This is the maximum distance at which Earth's gravitational influence is stronger than the more distant Sun and planets. Objects must orbit Earth within this radius, or they can become unbound by the gravitational perturbation of the Sun.
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Earth, along with the Solar System, is situated in the Milky Way and orbits about 28,000 light-years from its center. It is about 20 light-years above the galactic plane in the Orion Arm.[190]
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The axial tilt of Earth is approximately 23.439281°[2] with the axis of its orbit plane, always pointing towards the Celestial Poles. Due to Earth's axial tilt, the amount of sunlight reaching any given point on the surface varies over the course of the year. This causes the seasonal change in climate, with summer in the Northern Hemisphere occurring when the Tropic of Cancer is facing the Sun, and winter taking place when the Tropic of Capricorn in the Southern Hemisphere faces the Sun. During the summer, the day lasts longer, and the Sun climbs higher in the sky. In winter, the climate becomes cooler and the days shorter. In northern temperate latitudes, the Sun rises north of true east during the summer solstice, and sets north of true west, reversing in the winter. The Sun rises south of true east in the summer for the southern temperate zone and sets south of true west.
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Above the Arctic Circle, an extreme case is reached where there is no daylight at all for part of the year, up to six months at the North Pole itself, a polar night. In the Southern Hemisphere, the situation is exactly reversed, with the South Pole oriented opposite the direction of the North Pole. Six months later, this pole will experience a midnight sun, a day of 24 hours, again reversing with the South Pole.
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By astronomical convention, the four seasons can be determined by the solstices—the points in the orbit of maximum axial tilt toward or away from the Sun—and the equinoxes, when Earth's rotational axis is aligned with its orbital axis. In the Northern Hemisphere, winter solstice currently occurs around 21 December; summer solstice is near 21 June, spring equinox is around 20 March and autumnal equinox is about 22 or 23 September. In the Southern Hemisphere, the situation is reversed, with the summer and winter solstices exchanged and the spring and autumnal equinox dates swapped.[191]
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The angle of Earth's axial tilt is relatively stable over long periods of time. Its axial tilt does undergo nutation; a slight, irregular motion with a main period of 18.6 years.[192] The orientation (rather than the angle) of Earth's axis also changes over time, precessing around in a complete circle over each 25,800 year cycle; this precession is the reason for the difference between a sidereal year and a tropical year. Both of these motions are caused by the varying attraction of the Sun and the Moon on Earth's equatorial bulge. The poles also migrate a few meters across Earth's surface. This polar motion has multiple, cyclical components, which collectively are termed quasiperiodic motion. In addition to an annual component to this motion, there is a 14-month cycle called the Chandler wobble. Earth's rotational velocity also varies in a phenomenon known as length-of-day variation.[193]
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In modern times, Earth's perihelion occurs around 3 January, and its aphelion around 4 July. These dates change over time due to precession and other orbital factors, which follow cyclical patterns known as Milankovitch cycles. The changing Earth–Sun distance causes an increase of about 6.9%[n 21] in solar energy reaching Earth at perihelion relative to aphelion. Because the Southern Hemisphere is tilted toward the Sun at about the same time that Earth reaches the closest approach to the Sun, the Southern Hemisphere receives slightly more energy from the Sun than does the northern over the course of a year. This effect is much less significant than the total energy change due to the axial tilt, and most of the excess energy is absorbed by the higher proportion of water in the Southern Hemisphere.[194]
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A study from 2016 suggested that Planet Nine tilted all the planets of the Solar System, including Earth, by about six degrees.[195]
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A planet that can sustain life is termed habitable, even if life did not originate there. Earth provides liquid water—an environment where complex organic molecules can assemble and interact, and sufficient energy to sustain metabolism.[196] The distance of Earth from the Sun, as well as its orbital eccentricity, rate of rotation, axial tilt, geological history, sustaining atmosphere, and magnetic field all contribute to the current climatic conditions at the surface.[197]
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A planet's life forms inhabit ecosystems, whose total is sometimes said to form a "biosphere".[198] Earth's biosphere is thought to have begun evolving about 3.5 Gya.[65] The biosphere is divided into a number of biomes, inhabited by broadly similar plants and animals.[199] On land, biomes are separated primarily by differences in latitude, height above sea level and humidity. Terrestrial biomes lying within the Arctic or Antarctic Circles, at high altitudes or in extremely arid areas are relatively barren of plant and animal life; species diversity reaches a peak in humid lowlands at equatorial latitudes.[200]
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In July 2016, scientists reported identifying a set of 355 genes from the last universal common ancestor (LUCA) of all organisms living on Earth.[201]
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Earth has resources that have been exploited by humans.[203] Those termed non-renewable resources, such as fossil fuels, only renew over geological timescales.[204]
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Large deposits of fossil fuels are obtained from Earth's crust, consisting of coal, petroleum, and natural gas.[205] These deposits are used by humans both for energy production and as feedstock for chemical production.[206] Mineral ore bodies have also been formed within the crust through a process of ore genesis, resulting from actions of magmatism, erosion, and plate tectonics.[207] These bodies form concentrated sources for many metals and other useful elements.
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Earth's biosphere produces many useful biological products for humans, including food, wood, pharmaceuticals, oxygen, and the recycling of many organic wastes. The land-based ecosystem depends upon topsoil and fresh water, and the oceanic ecosystem depends upon dissolved nutrients washed down from the land.[208] In 1980, 50.53 million km2 (19.51 million sq mi) of Earth's land surface consisted of forest and woodlands, 67.88 million km2 (26.21 million sq mi) was grasslands and pasture, and 15.01 million km2 (5.80 million sq mi) was cultivated as croplands.[209] The estimated amount of irrigated land in 1993 was 2,481,250 km2 (958,020 sq mi).[13] Humans also live on the land by using building materials to construct shelters.
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Large areas of Earth's surface are subject to extreme weather such as tropical cyclones, hurricanes, or typhoons that dominate life in those areas. From 1980 to 2000, these events caused an average of 11,800 human deaths per year.[210] Many places are subject to earthquakes, landslides, tsunamis, volcanic eruptions, tornadoes, sinkholes, blizzards, floods, droughts, wildfires, and other calamities and disasters.
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Many localized areas are subject to human-made pollution of the air and water, acid rain and toxic substances, loss of vegetation (overgrazing, deforestation, desertification), loss of wildlife, species extinction, soil degradation, soil depletion and erosion.
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There is a scientific consensus linking human activities to global warming due to industrial carbon dioxide emissions. This is predicted to produce changes such as the melting of glaciers and ice sheets, more extreme temperature ranges, significant changes in weather and a global rise in average sea levels.[211]
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Cartography, the study and practice of map-making, and geography, the study of the lands, features, inhabitants and phenomena on Earth, have historically been the disciplines devoted to depicting Earth. Surveying, the determination of locations and distances, and to a lesser extent navigation, the determination of position and direction, have developed alongside cartography and geography, providing and suitably quantifying the requisite information.
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Earth's human population reached approximately seven billion on 31 October 2011.[213] Projections indicate that the world's human population will reach 9.2 billion in 2050.[214] Most of the growth is expected to take place in developing nations. Human population density varies widely around the world, but a majority live in Asia. By 2020, 60% of the world's population is expected to be living in urban, rather than rural, areas.[215]
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68% of the land mass of the world is in the northern hemisphere.[216] Partly due to the predominance of land mass, 90% of humans live in the northern hemisphere.[217]
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It is estimated that one-eighth of Earth's surface is suitable for humans to live on – three-quarters of Earth's surface is covered by oceans, leaving one-quarter as land. Half of that land area is desert (14%),[218] high mountains (27%),[219] or other unsuitable terrains. The northernmost permanent settlement in the world is Alert, on Ellesmere Island in Nunavut, Canada.[220] (82°28′N) The southernmost is the Amundsen–Scott South Pole Station, in Antarctica, almost exactly at the South Pole. (90°S)
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States claim the planet's entire land surface, except for parts of Antarctica and a few other unclaimed areas. Earth has never had a planetwide government, but the United Nations is the leading worldwide intergovernmental organization.
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The first human to orbit Earth was Yuri Gagarin on 12 April 1961.[221] In total, about 487 people have visited outer space and reached orbit as of 30 July 2010[update], and, of these, twelve have walked on the Moon.[222][223][224] Normally, the only humans in space are those on the International Space Station. The station's crew, made up of six people, is usually replaced every six months.[225] The farthest that humans have traveled from Earth is 400,171 km (248,655 mi), achieved during the Apollo 13 mission in 1970.[226]
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The Moon is a relatively large, terrestrial, planet-like natural satellite, with a diameter about one-quarter of Earth's. It is the largest moon in the Solar System relative to the size of its planet, although Charon is larger relative to the dwarf planet Pluto. The natural satellites of other planets are also referred to as "moons", after Earth's.
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The gravitational attraction between Earth and the Moon causes tides on Earth. The same effect on the Moon has led to its tidal locking: its rotation period is the same as the time it takes to orbit Earth. As a result, it always presents the same face to the planet. As the Moon orbits Earth, different parts of its face are illuminated by the Sun, leading to the lunar phases; the dark part of the face is separated from the light part by the solar terminator.
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Due to their tidal interaction, the Moon recedes from Earth at the rate of approximately 38 mm/a (1.5 in/year). Over millions of years, these tiny modifications—and the lengthening of Earth's day by about 23 µs/yr—add up to significant changes.[227] During the Devonian period, for example, (approximately 410 Mya) there were 400 days in a year, with each day lasting 21.8 hours.[228]
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The Moon may have dramatically affected the development of life by moderating the planet's climate. Paleontological evidence and computer simulations show that Earth's axial tilt is stabilized by tidal interactions with the Moon.[229] Some theorists think that without this stabilization against the torques applied by the Sun and planets to Earth's equatorial bulge, the rotational axis might be chaotically unstable, exhibiting chaotic changes over millions of years, as appears to be the case for Mars.[230]
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Viewed from Earth, the Moon is just far enough away to have almost the same apparent-sized disk as the Sun. The angular size (or solid angle) of these two bodies match because, although the Sun's diameter is about 400 times as large as the Moon's, it is also 400 times more distant.[186] This allows total and annular solar eclipses to occur on Earth.
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The most widely accepted theory of the Moon's origin, the giant-impact hypothesis, states that it formed from the collision of a Mars-size protoplanet called Theia with the early Earth. This hypothesis explains (among other things) the Moon's relative lack of iron and volatile elements and the fact that its composition is nearly identical to that of Earth's crust.[48]
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Earth has at least five co-orbital asteroids, including 3753 Cruithne and 2002 AA29.[231][232] A trojan asteroid companion, 2010 TK7, is librating around the leading Lagrange triangular point, L4, in Earth's orbit around the Sun.[233][234]
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The tiny near-Earth asteroid 2006 RH120 makes close approaches to the Earth–Moon system roughly every twenty years. During these approaches, it can orbit Earth for brief periods of time.[235]
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As of April 2018[update], there are 1,886 operational, human-made satellites orbiting Earth.[5] There are also inoperative satellites, including Vanguard 1, the oldest satellite currently in orbit, and over 16,000 pieces of tracked space debris.[n 3] Earth's largest artificial satellite is the International Space Station.
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The standard astronomical symbol of Earth consists of a cross circumscribed by a circle, ,[236] representing the four corners of the world.
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Human cultures have developed many views of the planet.[237] Earth is sometimes personified as a deity. In many cultures it is a mother goddess that is also the primary fertility deity,[238] and by the mid-20th century, the Gaia Principle compared Earth's environments and life as a single self-regulating organism leading to broad stabilization of the conditions of habitability.[239][240][241] Creation myths in many religions involve the creation of Earth by a supernatural deity or deities.[238]
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The Hindu Vedas (1500–900 BC) refer to the Earth as Bhūgola (भूगोल), which comes from Bhū (earth, ground) and Gola (ball, sphere, globe). It means the "globe of earth". There is no direct evidence that the Hindus of that time knew that the Earth was sphere-shaped, but this name has been used extensively since the inception of the Vedas.[citation needed]
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Scientific investigation has resulted in several culturally transformative shifts in people's view of the planet. Initial belief in a flat Earth was gradually displaced in the Greek colonies of southern Italy during the late 6th century BC by the idea of spherical Earth,[242][243][244] which was attributed to both the philosophers Pythagoras and Parmenides.[243][244] By the end of the 5th century BC, the sphericity of Earth was universally accepted among Greek intellectuals.[245] Earth was generally believed to be the center of the universe until the 16th century, when scientists first conclusively demonstrated that it was a moving object, comparable to the other planets in the Solar System.[246] Due to the efforts of influential Christian scholars and clerics such as James Ussher, who sought to determine the age of Earth through analysis of genealogies in Scripture, Westerners before the 19th century generally believed Earth to be a few thousand years old at most. It was only during the 19th century that geologists realized Earth's age was at least many millions of years.[247]
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Lord Kelvin used thermodynamics to estimate the age of Earth to be between 20 million and 400 million years in 1864, sparking a vigorous debate on the subject; it was only when radioactivity and radioactive dating were discovered in the late 19th and early 20th centuries that a reliable mechanism for determining Earth's age was established, proving the planet to be billions of years old.[248][249] The perception of Earth shifted again[further explanation needed] in the 20th century when humans first viewed it from orbit, and especially with photographs of Earth returned by the Apollo program.[250][251][252]
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Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
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Earth is the third planet from the Sun and the only astronomical object known to harbor life. According to radiometric dating estimation and other evidence, Earth formed over 4.5 billion years ago. Earth's gravity interacts with other objects in space, especially the Sun and the Moon, which is Earth's only natural satellite. Earth orbits around the Sun in 365.256 solar days, a period known as an Earth sidereal year. During this time, Earth rotates about its axis 366.256 times, that is, a sidereal year has 366.256 sidereal days.[n 6]
|
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|
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Earth's axis of rotation is tilted with respect to its orbital plane, producing seasons on Earth. The gravitational interaction between Earth and the Moon causes tides, stabilizes Earth's orientation on its axis, and gradually slows its rotation. Earth is the densest planet in the Solar System and the largest and most massive of the four rocky planets.
|
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|
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Earth's outer layer (lithosphere) is divided into several rigid tectonic plates that migrate across the surface over many millions of years. About 29% of Earth's surface is land consisting of continents and islands. The remaining 71% is covered with water, mostly by oceans but also lakes, rivers and other fresh water, which all together constitute the hydrosphere. The majority of Earth's polar regions are covered in ice, including the Antarctic ice sheet and the sea ice of the Arctic ice pack. Earth's interior remains active with a solid iron inner core, a liquid outer core that generates Earth's magnetic field, and a convecting mantle that drives plate tectonics.
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Within the first billion years of Earth's history, life appeared in the oceans and began to affect Earth's atmosphere and surface, leading to the proliferation of anaerobic and, later, aerobic organisms. Some geological evidence indicates that life may have arisen as early as 4.1 billion years ago. Since then, the combination of Earth's distance from the Sun, physical properties and geological history have allowed life to evolve and thrive. In the history of life on Earth, biodiversity has gone through long periods of expansion, occasionally punctuated by mass extinctions. Over 99% of all species that ever lived on Earth are extinct. Estimates of the number of species on Earth today vary widely; most species have not been described. Over 7.7 billion humans live on Earth and depend on its biosphere and natural resources for their survival.[23]
|
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|
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The modern English word Earth developed, via Middle English,[n 7] from an Old English noun most often spelled eorðe.[24] It has cognates in every Germanic language, and their ancestral root has been reconstructed as *erþō. In its earliest attestation, the word eorðe was already being used to translate the many senses of Latin terra and Greek γῆ gē: the ground,[n 8] its soil,[n 9] dry land,[n 10] the human world,[n 11] the surface of the world (including the sea),[n 12] and the globe itself.[n 13] As with Roman Terra/Tellūs and Greek Gaia, Earth may have been a personified goddess in Germanic paganism: late Norse mythology included Jörð ('Earth'), a giantess often given as the mother of Thor.[33]
|
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+
|
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+
Originally, earth was written in lowercase, and from early Middle English, its definite sense as "the globe" was expressed as the earth. By Early Modern English, many nouns were capitalized, and the earth became (and often remained) the Earth, particularly when referenced along with other heavenly bodies. More recently, the name is sometimes simply given as Earth, by analogy with the names of the other planets.[24] House styles now vary: Oxford spelling recognizes the lowercase form as the most common, with the capitalized form an acceptable variant. Another convention capitalizes "Earth" when appearing as a name (e.g. "Earth's atmosphere") but writes it in lowercase when preceded by the (e.g. "the atmosphere of the earth"). It almost always appears in lowercase in colloquial expressions such as "what on earth are you doing?"[34]
|
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|
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Occasionally, the name Terra /ˈtɛrə/ is used in scientific writing and especially in science fiction to distinguish our inhabited planet from others,[35] while in poetry Tellus /ˈtɛləs/ has been used to denote personification of the Earth.[36] The Greek poetic name Gaea (Gæa) /ˈdʒiːə/ is rare, though the alternative spelling Gaia has become common due to the Gaia hypothesis, in which case its pronunciation is /ˈɡaɪə/ rather than the more Classical /ˈɡeɪə/.[37]
|
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+
|
19 |
+
There are a number of adjectives for the planet Earth. From Earth itself comes earthly. From Latin Terra come Terran /ˈtɛrən/,[38] Terrestrial /təˈrɛstriəl/,[39] and (via French) Terrene /təˈriːn/,[40] and from Latin Tellus come Tellurian /tɛˈlʊəriən/[41] and, more rarely, Telluric and Tellural. From Greek Gaia and Gaea comes Gaian and Gaean.
|
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|
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An inhabitant of the Earth is an Earthling, a Terran, a Terrestrial, a Tellurian or, rarely, an Earthian.
|
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+
|
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+
The oldest material found in the Solar System is dated to 4.5672±0.0006 billion years ago (BYA).[42] By 4.54±0.04 BYA[43] the primordial Earth had formed. The bodies in the Solar System formed and evolved with the Sun. In theory, a solar nebula partitions a volume out of a molecular cloud by gravitational collapse, which begins to spin and flatten into a circumstellar disk, and then the planets grow out of that disk with the Sun. A nebula contains gas, ice grains, and dust (including primordial nuclides). According to nebular theory, planetesimals formed by accretion, with the primordial Earth taking 10–20 million years (Mys) to form.[44]
|
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+
|
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+
A subject of research is the formation of the Moon, some 4.53 BYA.[45] A leading hypothesis is that it was formed by accretion from material loosed from Earth after a Mars-sized object, named Theia, hit Earth.[46] In this view, the mass of Theia was approximately 10 percent of Earth;[47] it hit Earth with a glancing blow and some of its mass merged with Earth.[48] Between approximately 4.1 and 3.8 BYA, numerous asteroid impacts during the Late Heavy Bombardment caused significant changes to the greater surface environment of the Moon and, by inference, to that of Earth.
|
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|
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Earth's atmosphere and oceans were formed by volcanic activity and outgassing.[49] Water vapor from these sources condensed into the oceans, augmented by water and ice from asteroids, protoplanets, and comets.[50] In this model, atmospheric "greenhouse gases" kept the oceans from freezing when the newly forming Sun had only 70% of its current luminosity.[51] By 3.5 BYA, Earth's magnetic field was established, which helped prevent the atmosphere from being stripped away by the solar wind.[52]
|
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|
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A crust formed when the molten outer layer of Earth cooled to form a solid. The two models[53] that explain land mass propose either a steady growth to the present-day forms[54] or, more likely, a rapid growth[55] early in Earth history[56] followed by a long-term steady continental area.[57][58][59] Continents formed by plate tectonics, a process ultimately driven by the continuous loss of heat from Earth's interior. Over the period of hundreds of millions of years, the supercontinents have assembled and broken apart. Roughly 750 million years ago (MYA), one of the earliest known supercontinents, Rodinia, began to break apart. The continents later recombined to form Pannotia 600–540 MYA, then finally Pangaea, which also broke apart 180 MYA.[60]
|
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|
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The present pattern of ice ages began about 40 MYA,[61] and then intensified during the Pleistocene about 3 MYA.[62] High-latitude regions have since undergone repeated cycles of glaciation and thaw, repeating about every 40,000–100,000 years. The last continental glaciation ended 10,000 years ago.[63]
|
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|
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Chemical reactions led to the first self-replicating molecules about four billion years ago. A half billion years later, the last common ancestor of all current life arose.[64] The evolution of photosynthesis allowed the Sun's energy to be harvested directly by life forms. The resultant molecular oxygen (O2) accumulated in the atmosphere and due to interaction with ultraviolet solar radiation, formed a protective ozone layer (O3) in the upper atmosphere.[65] The incorporation of smaller cells within larger ones resulted in the development of complex cells called eukaryotes.[66] True multicellular organisms formed as cells within colonies became increasingly specialized. Aided by the absorption of harmful ultraviolet radiation by the ozone layer, life colonized Earth's surface.[67] Among the earliest fossil evidence for life is microbial mat fossils found in 3.48 billion-year-old sandstone in Western Australia,[68] biogenic graphite found in 3.7 billion-year-old metasedimentary rocks in Western Greenland,[69] and remains of biotic material found in 4.1 billion-year-old rocks in Western Australia.[70][71] The earliest direct evidence of life on Earth is contained in 3.45 billion-year-old Australian rocks showing fossils of microorganisms.[72][73]
|
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During the Neoproterozoic, 750 to 580 MYA, much of Earth might have been covered in ice. This hypothesis has been termed "Snowball Earth", and it is of particular interest because it preceded the Cambrian explosion, when multicellular life forms significantly increased in complexity.[74] Following the Cambrian explosion, 535 MYA, there have been five mass extinctions.[75] The most recent such event was 66 MYA, when an asteroid impact triggered the extinction of the non-avian dinosaurs and other large reptiles, but spared some small animals such as mammals, which at the time resembled shrews. Mammalian life has diversified over the past 66 Mys, and several million years ago an African ape-like animal such as Orrorin tugenensis gained the ability to stand upright.[76] This facilitated tool use and encouraged communication that provided the nutrition and stimulation needed for a larger brain, which led to the evolution of humans. The development of agriculture, and then civilization, led to humans having an influence on Earth and the nature and quantity of other life forms that continues to this day.[77]
|
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Earth's expected long-term future is tied to that of the Sun. Over the next 1.1 billion years, solar luminosity will increase by 10%, and over the next 3.5 billion years by 40%.[78] Earth's increasing surface temperature will accelerate the inorganic carbon cycle, reducing CO2 concentration to levels lethally low for plants (10 ppm for C4 photosynthesis) in approximately 100–900 million years.[79][80] The lack of vegetation will result in the loss of oxygen in the atmosphere, making animal life impossible.[81] About a billion years from now, all surface water will have disappeared[82] and the mean global temperature will reach 70 °C (158 °F).[81] Earth is expected to be habitable until the end of photosynthesis about 500 million years from now,[79] but if nitrogen is removed from the atmosphere, life may continue until a runaway greenhouse effect occurs 2.3 billion years from now.[80] Anthropogenic emissions are "probably insufficient" to cause a runaway greenhouse at current solar luminosity.[83] Even if the Sun were eternal and stable, 27% of the water in the modern oceans will descend to the mantle in one billion years, due to reduced steam venting from mid-ocean ridges.[84]
|
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|
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The Sun will evolve to become a red giant in about 5 billion years. Models predict that the Sun will expand to roughly 1 AU (150 million km; 93 million mi), about 250 times its present radius.[78][85] Earth's fate is less clear. As a red giant, the Sun will lose roughly 30% of its mass, so, without tidal effects, Earth will move to an orbit 1.7 AU (250 million km; 160 million mi) from the Sun when the star reaches its maximum radius. Most, if not all, remaining life will be destroyed by the Sun's increased luminosity (peaking at about 5,000 times its present level).[78] A 2008 simulation indicates that Earth's orbit will eventually decay due to tidal effects and drag, causing it to enter the Sun's atmosphere and be vaporized.[85]
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The shape of Earth is nearly spherical. There is a small flattening at the poles and bulging around the equator due to Earth's rotation.[89] To second order, Earth is approximately an oblate spheroid, whose equatorial diameter is 43 kilometres (27 mi) larger than the pole-to-pole diameter,[90] although the variation is less than 1% of the average radius of the Earth.
|
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The point on the surface farthest from Earth's center of mass is the summit of the equatorial Chimborazo volcano in Ecuador (6,384.4 km or 3,967.1 mi).[91][92][93][94] The average diameter of the reference spheroid is 12,742 kilometres (7,918 mi). Local topography deviates from this idealized spheroid, although on a global scale these deviations are small compared to Earth's radius: the maximum deviation of only 0.17% is at the Mariana Trench (10,911 metres or 35,797 feet below local sea level), whereas Mount Everest (8,848 metres or 29,029 feet above local sea level) represents a deviation of 0.14%.[n 14]
|
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In geodesy, the exact shape that Earth's oceans would adopt in the absence of land and perturbations such as tides and winds is called the geoid. More precisely, the geoid is the surface of gravitational equipotential at mean sea level.
|
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|
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Earth's mass is approximately 5.97×1024 kg (5,970 Yg). It is composed mostly of iron (32.1%), oxygen (30.1%), silicon (15.1%), magnesium (13.9%), sulphur (2.9%), nickel (1.8%), calcium (1.5%), and aluminum (1.4%), with the remaining 1.2% consisting of trace amounts of other elements. Due to mass segregation, the core region is estimated to be primarily composed of iron (88.8%), with smaller amounts of nickel (5.8%), sulphur (4.5%), and less than 1% trace elements.[98]
|
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The most common rock constituents of the crust are nearly all oxides: chlorine, sulphur, and fluorine are the important exceptions to this and their total amount in any rock is usually much less than 1%. Over 99% of the crust is composed of 11 oxides, principally silica, alumina, iron oxides, lime, magnesia, potash and soda.[99][98][100]
|
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Earth's interior, like that of the other terrestrial planets, is divided into layers by their chemical or physical (rheological) properties. The outer layer is a chemically distinct silicate solid crust, which is underlain by a highly viscous solid mantle. The crust is separated from the mantle by the Mohorovičić discontinuity. The thickness of the crust varies from about 6 kilometres (3.7 mi) under the oceans to 30–50 km (19–31 mi) for the continents. The crust and the cold, rigid, top of the upper mantle are collectively known as the lithosphere, and it is of the lithosphere that the tectonic plates are composed. Beneath the lithosphere is the asthenosphere, a relatively low-viscosity layer on which the lithosphere rides. Important changes in crystal structure within the mantle occur at 410 and 660 km (250 and 410 mi) below the surface, spanning a transition zone that separates the upper and lower mantle. Beneath the mantle, an extremely low viscosity liquid outer core lies above a solid inner core.[101] Earth's inner core might rotate at a slightly higher angular velocity than the remainder of the planet, advancing by 0.1–0.5° per year.[102] The radius of the inner core is about one fifth of that of Earth.
|
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Earth's internal heat comes from a combination of residual heat from planetary accretion (about 20%) and heat produced through radioactive decay (80%).[105] The major heat-producing isotopes within Earth are potassium-40, uranium-238, and thorium-232.[106] At the center, the temperature may be up to 6,000 °C (10,830 °F),[107] and the pressure could reach 360 GPa (52 million psi).[108] Because much of the heat is provided by radioactive decay, scientists postulate that early in Earth's history, before isotopes with short half-lives were depleted, Earth's heat production was much higher. At approximately 3 Gyr, twice the present-day heat would have been produced, increasing the rates of mantle convection and plate tectonics, and allowing the production of uncommon igneous rocks such as komatiites that are rarely formed today.[105][109]
|
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The mean heat loss from Earth is 87 mW m−2, for a global heat loss of 4.42×1013 W.[111] A portion of the core's thermal energy is transported toward the crust by mantle plumes, a form of convection consisting of upwellings of higher-temperature rock. These plumes can produce hotspots and flood basalts.[112] More of the heat in Earth is lost through plate tectonics, by mantle upwelling associated with mid-ocean ridges. The final major mode of heat loss is through conduction through the lithosphere, the majority of which occurs under the oceans because the crust there is much thinner than that of the continents.[113]
|
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Earth's mechanically rigid outer layer, the lithosphere, is divided into tectonic plates. These plates are rigid segments that move relative to each other at one of three boundaries types: At convergent boundaries, two plates come together; at divergent boundaries, two plates are pulled apart; and at transform boundaries, two plates slide past one another laterally. Along these plate boundaries, earthquakes, volcanic activity, mountain-building, and oceanic trench formation can occur.[115] The tectonic plates ride on top of the asthenosphere, the solid but less-viscous part of the upper mantle that can flow and move along with the plates.[116]
|
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|
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As the tectonic plates migrate, oceanic crust is subducted under the leading edges of the plates at convergent boundaries. At the same time, the upwelling of mantle material at divergent boundaries creates mid-ocean ridges. The combination of these processes recycles the oceanic crust back into the mantle. Due to this recycling, most of the ocean floor is less than 100 Ma old. The oldest oceanic crust is located in the Western Pacific and is estimated to be 200 Ma old.[117][118] By comparison, the oldest dated continental crust is 4,030 Ma.[119]
|
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The seven major plates are the Pacific, North American, Eurasian, African, Antarctic, Indo-Australian, and South American. Other notable plates include the Arabian Plate, the Caribbean Plate, the Nazca Plate off the west coast of South America and the Scotia Plate in the southern Atlantic Ocean. The Australian Plate fused with the Indian Plate between 50 and 55 Mya. The fastest-moving plates are the oceanic plates, with the Cocos Plate advancing at a rate of 75 mm/a (3.0 in/year)[120] and the Pacific Plate moving 52–69 mm/a (2.0–2.7 in/year). At the other extreme, the slowest-moving plate is the Eurasian Plate, progressing at a typical rate of 21 mm/a (0.83 in/year).[121]
|
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The total surface area of Earth is about 510 million km2 (197 million sq mi).[12] Of this, 70.8%,[12] or 361.13 million km2 (139.43 million sq mi), is below sea level and covered by ocean water.[122] Below the ocean's surface are much of the continental shelf, mountains, volcanoes,[90] oceanic trenches, submarine canyons, oceanic plateaus, abyssal plains, and a globe-spanning mid-ocean ridge system. The remaining 29.2%, or 148.94 million km2 (57.51 million sq mi), not covered by water has terrain that varies greatly from place to place and consists of mountains, deserts, plains, plateaus, and other landforms. Tectonics and erosion, volcanic eruptions, flooding, weathering, glaciation, the growth of coral reefs, and meteorite impacts are among the processes that constantly reshape Earth's surface over geological time.[123][124]
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The continental crust consists of lower density material such as the igneous rocks granite and andesite. Less common is basalt, a denser volcanic rock that is the primary constituent of the ocean floors.[125] Sedimentary rock is formed from the accumulation of sediment that becomes buried and compacted together. Nearly 75% of the continental surfaces are covered by sedimentary rocks, although they form about 5% of the crust.[126] The third form of rock material found on Earth is metamorphic rock, which is created from the transformation of pre-existing rock types through high pressures, high temperatures, or both. The most abundant silicate minerals on Earth's surface include quartz, feldspars, amphibole, mica, pyroxene and olivine.[127] Common carbonate minerals include calcite (found in limestone) and dolomite.[128]
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The elevation of the land surface varies from the low point of −418 m (−1,371 ft) at the Dead Sea, to a maximum altitude of 8,848 m (29,029 ft) at the top of Mount Everest. The mean height of land above sea level is about 797 m (2,615 ft).[129]
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The pedosphere is the outermost layer of Earth's continental surface and is composed of soil and subject to soil formation processes. The total arable land is 10.9% of the land surface, with 1.3% being permanent cropland.[130][131] Close to 40% of Earth's land surface is used for agriculture, or an estimated 16.7 million km2 (6.4 million sq mi) of cropland and 33.5 million km2 (12.9 million sq mi) of pastureland.[132]
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The abundance of water on Earth's surface is a unique feature that distinguishes the "Blue Planet" from other planets in the Solar System. Earth's hydrosphere consists chiefly of the oceans, but technically includes all water surfaces in the world, including inland seas, lakes, rivers, and underground waters down to a depth of 2,000 m (6,600 ft). The deepest underwater location is Challenger Deep of the Mariana Trench in the Pacific Ocean with a depth of 10,911.4 m (35,799 ft).[n 18][133]
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The mass of the oceans is approximately 1.35×1018 metric tons or about 1/4400 of Earth's total mass. The oceans cover an area of 361.8 million km2 (139.7 million sq mi) with a mean depth of 3,682 m (12,080 ft), resulting in an estimated volume of 1.332 billion km3 (320 million cu mi).[134] If all of Earth's crustal surface were at the same elevation as a smooth sphere, the depth of the resulting world ocean would be 2.7 to 2.8 km (1.68 to 1.74 mi).[135][136]
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About 97.5% of the water is saline; the remaining 2.5% is fresh water. Most fresh water, about 68.7%, is present as ice in ice caps and glaciers.[137]
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The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5% salt).[138] Most of this salt was released from volcanic activity or extracted from cool igneous rocks.[139] The oceans are also a reservoir of dissolved atmospheric gases, which are essential for the survival of many aquatic life forms.[140] Sea water has an important influence on the world's climate, with the oceans acting as a large heat reservoir.[141] Shifts in the oceanic temperature distribution can cause significant weather shifts, such as the El Niño–Southern Oscillation.[142]
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The atmospheric pressure at Earth's sea level averages 101.325 kPa (14.696 psi),[143] with a scale height of about 8.5 km (5.3 mi).[3] A dry atmosphere is composed of 78.084% nitrogen, 20.946% oxygen, 0.934% argon, and trace amounts of carbon dioxide and other gaseous molecules.[143] Water vapor content varies between 0.01% and 4%[143] but averages about 1%.[3] The height of the troposphere varies with latitude, ranging between 8 km (5 mi) at the poles to 17 km (11 mi) at the equator, with some variation resulting from weather and seasonal factors.[144]
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Earth's biosphere has significantly altered its atmosphere. Oxygenic photosynthesis evolved 2.7 Gya, forming the primarily nitrogen–oxygen atmosphere of today.[65] This change enabled the proliferation of aerobic organisms and, indirectly, the formation of the ozone layer due to the subsequent conversion of atmospheric O2 into O3. The ozone layer blocks ultraviolet solar radiation, permitting life on land.[145] Other atmospheric functions important to life include transporting water vapor, providing useful gases, causing small meteors to burn up before they strike the surface, and moderating temperature.[146] This last phenomenon is known as the greenhouse effect: trace molecules within the atmosphere serve to capture thermal energy emitted from the ground, thereby raising the average temperature. Water vapor, carbon dioxide, methane, nitrous oxide, and ozone are the primary greenhouse gases in the atmosphere. Without this heat-retention effect, the average surface temperature would be −18 °C (0 °F), in contrast to the current +15 °C (59 °F),[147] and life on Earth probably would not exist in its current form.[148] In May 2017, glints of light, seen as twinkling from an orbiting satellite a million miles away, were found to be reflected light from ice crystals in the atmosphere.[149][150]
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Earth's atmosphere has no definite boundary, slowly becoming thinner and fading into outer space. Three-quarters of the atmosphere's mass is contained within the first 11 km (6.8 mi) of the surface. This lowest layer is called the troposphere. Energy from the Sun heats this layer, and the surface below, causing expansion of the air. This lower-density air then rises and is replaced by cooler, higher-density air. The result is atmospheric circulation that drives the weather and climate through redistribution of thermal energy.[151]
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The primary atmospheric circulation bands consist of the trade winds in the equatorial region below 30° latitude and the westerlies in the mid-latitudes between 30° and 60°.[152] Ocean currents are also important factors in determining climate, particularly the thermohaline circulation that distributes thermal energy from the equatorial oceans to the polar regions.[153]
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Water vapor generated through surface evaporation is transported by circulatory patterns in the atmosphere. When atmospheric conditions permit an uplift of warm, humid air, this water condenses and falls to the surface as precipitation.[151] Most of the water is then transported to lower elevations by river systems and usually returned to the oceans or deposited into lakes. This water cycle is a vital mechanism for supporting life on land and is a primary factor in the erosion of surface features over geological periods. Precipitation patterns vary widely, ranging from several meters of water per year to less than a millimeter. Atmospheric circulation, topographic features, and temperature differences determine the average precipitation that falls in each region.[154]
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The amount of solar energy reaching Earth's surface decreases with increasing latitude. At higher latitudes, the sunlight reaches the surface at lower angles, and it must pass through thicker columns of the atmosphere. As a result, the mean annual air temperature at sea level decreases by about 0.4 °C (0.7 °F) per degree of latitude from the equator.[155] Earth's surface can be subdivided into specific latitudinal belts of approximately homogeneous climate. Ranging from the equator to the polar regions, these are the tropical (or equatorial), subtropical, temperate and polar climates.[156]
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This latitudinal rule has several anomalies:
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The commonly used Köppen climate classification system has five broad groups (humid tropics, arid, humid middle latitudes, continental and cold polar), which are further divided into more specific subtypes.[152] The Köppen system rates regions of terrain based on observed temperature and precipitation.
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The highest air temperature ever measured on Earth was 56.7 °C (134.1 °F) in Furnace Creek, California, in Death Valley, in 1913.[159] The lowest air temperature ever directly measured on Earth was −89.2 °C (−128.6 °F) at Vostok Station in 1983,[160] but satellites have used remote sensing to measure temperatures as low as −94.7 °C (−138.5 °F) in East Antarctica.[161] These temperature records are only measurements made with modern instruments from the 20th century onwards and likely do not reflect the full range of temperature on Earth.
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Above the troposphere, the atmosphere is usually divided into the stratosphere, mesosphere, and thermosphere.[146] Each layer has a different lapse rate, defining the rate of change in temperature with height. Beyond these, the exosphere thins out into the magnetosphere, where the geomagnetic fields interact with the solar wind.[162] Within the stratosphere is the ozone layer, a component that partially shields the surface from ultraviolet light and thus is important for life on Earth. The Kármán line, defined as 100 km above Earth's surface, is a working definition for the boundary between the atmosphere and outer space.[163]
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Thermal energy causes some of the molecules at the outer edge of the atmosphere to increase their velocity to the point where they can escape from Earth's gravity. This causes a slow but steady loss of the atmosphere into space. Because unfixed hydrogen has a low molecular mass, it can achieve escape velocity more readily, and it leaks into outer space at a greater rate than other gases.[164] The leakage of hydrogen into space contributes to the shifting of Earth's atmosphere and surface from an initially reducing state to its current oxidizing one. Photosynthesis provided a source of free oxygen, but the loss of reducing agents such as hydrogen is thought to have been a necessary precondition for the widespread accumulation of oxygen in the atmosphere.[165] Hence the ability of hydrogen to escape from the atmosphere may have influenced the nature of life that developed on Earth.[166] In the current, oxygen-rich atmosphere most hydrogen is converted into water before it has an opportunity to escape. Instead, most of the hydrogen loss comes from the destruction of methane in the upper atmosphere.[167]
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The gravity of Earth is the acceleration that is imparted to objects due to the distribution of mass within Earth. Near Earth's surface, gravitational acceleration is approximately 9.8 m/s2 (32 ft/s2). Local differences in topography, geology, and deeper tectonic structure cause local and broad, regional differences in Earth's gravitational field, known as gravity anomalies.[168]
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The main part of Earth's magnetic field is generated in the core, the site of a dynamo process that converts the kinetic energy of thermally and compositionally driven convection into electrical and magnetic field energy. The field extends outwards from the core, through the mantle, and up to Earth's surface, where it is, approximately, a dipole. The poles of the dipole are located close to Earth's geographic poles. At the equator of the magnetic field, the magnetic-field strength at the surface is 3.05×10−5 T, with a magnetic dipole moment of 7.79×1022 Am2 at epoch 2000, decreasing nearly 6% per century.[169] The convection movements in the core are chaotic; the magnetic poles drift and periodically change alignment. This causes secular variation of the main field and field reversals at irregular intervals averaging a few times every million years. The most recent reversal occurred approximately 700,000 years ago.[170][171]
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The extent of Earth's magnetic field in space defines the magnetosphere. Ions and electrons of the solar wind are deflected by the magnetosphere; solar wind pressure compresses the dayside of the magnetosphere, to about 10 Earth radii, and extends the nightside magnetosphere into a long tail.[172] Because the velocity of the solar wind is greater than the speed at which waves propagate through the solar wind, a supersonic bow shock precedes the dayside magnetosphere within the solar wind.[173] Charged particles are contained within the magnetosphere; the plasmasphere is defined by low-energy particles that essentially follow magnetic field lines as Earth rotates;[174][175] the ring current is defined by medium-energy particles that drift relative to the geomagnetic field, but with paths that are still dominated by the magnetic field,[176] and the Van Allen radiation belt are formed by high-energy particles whose motion is essentially random, but otherwise contained by the magnetosphere.[172][177]
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During magnetic storms and substorms, charged particles can be deflected from the outer magnetosphere and especially the magnetotail, directed along field lines into Earth's ionosphere, where atmospheric atoms can be excited and ionized, causing the aurora.[178]
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Earth's rotation period relative to the Sun—its mean solar day—is 86,400 seconds of mean solar time (86,400.0025 SI seconds).[179] Because Earth's solar day is now slightly longer than it was during the 19th century due to tidal deceleration, each day varies between 0 and 2 SI ms longer.[180][181]
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Earth's rotation period relative to the fixed stars, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is 86,164.0989 seconds of mean solar time (UT1), or 23h 56m 4.0989s.[2][n 19] Earth's rotation period relative to the precessing or moving mean March equinox, misnamed its sidereal day, is 86,164.0905 seconds of mean solar time (UT1) (23h 56m 4.0905s).[2] Thus the sidereal day is shorter than the stellar day by about 8.4 ms.[182] The length of the mean solar day in SI seconds is available from the IERS for the periods 1623–2005[183] and 1962–2005.[184]
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Apart from meteors within the atmosphere and low-orbiting satellites, the main apparent motion of celestial bodies in Earth's sky is to the west at a rate of 15°/h = 15'/min. For bodies near the celestial equator, this is equivalent to an apparent diameter of the Sun or the Moon every two minutes; from Earth's surface, the apparent sizes of the Sun and the Moon are approximately the same.[185][186]
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Earth orbits the Sun at an average distance of about 150 million km (93 million mi) every 365.2564 mean solar days, or one sidereal year. This gives an apparent movement of the Sun eastward with respect to the stars at a rate of about 1°/day, which is one apparent Sun or Moon diameter every 12 hours. Due to this motion, on average it takes 24 hours—a solar day—for Earth to complete a full rotation about its axis so that the Sun returns to the meridian. The orbital speed of Earth averages about 29.78 km/s (107,200 km/h; 66,600 mph), which is fast enough to travel a distance equal to Earth's diameter, about 12,742 km (7,918 mi), in seven minutes, and the distance to the Moon, 384,000 km (239,000 mi), in about 3.5 hours.[3]
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The Moon and Earth orbit a common barycenter every 27.32 days relative to the background stars. When combined with the Earth–Moon system's common orbit around the Sun, the period of the synodic month, from new moon to new moon, is 29.53 days. Viewed from the celestial north pole, the motion of Earth, the Moon, and their axial rotations are all counterclockwise. Viewed from a vantage point above the north poles of both the Sun and Earth, Earth orbits in a counterclockwise direction about the Sun. The orbital and axial planes are not precisely aligned: Earth's axis is tilted some 23.44 degrees from the perpendicular to the Earth–Sun plane (the ecliptic), and the Earth–Moon plane is tilted up to ±5.1 degrees against the Earth–Sun plane. Without this tilt, there would be an eclipse every two weeks, alternating between lunar eclipses and solar eclipses.[3][188]
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The Hill sphere, or the sphere of gravitational influence, of Earth is about 1.5 million km (930,000 mi) in radius.[189][n 20] This is the maximum distance at which Earth's gravitational influence is stronger than the more distant Sun and planets. Objects must orbit Earth within this radius, or they can become unbound by the gravitational perturbation of the Sun.
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Earth, along with the Solar System, is situated in the Milky Way and orbits about 28,000 light-years from its center. It is about 20 light-years above the galactic plane in the Orion Arm.[190]
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The axial tilt of Earth is approximately 23.439281°[2] with the axis of its orbit plane, always pointing towards the Celestial Poles. Due to Earth's axial tilt, the amount of sunlight reaching any given point on the surface varies over the course of the year. This causes the seasonal change in climate, with summer in the Northern Hemisphere occurring when the Tropic of Cancer is facing the Sun, and winter taking place when the Tropic of Capricorn in the Southern Hemisphere faces the Sun. During the summer, the day lasts longer, and the Sun climbs higher in the sky. In winter, the climate becomes cooler and the days shorter. In northern temperate latitudes, the Sun rises north of true east during the summer solstice, and sets north of true west, reversing in the winter. The Sun rises south of true east in the summer for the southern temperate zone and sets south of true west.
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Above the Arctic Circle, an extreme case is reached where there is no daylight at all for part of the year, up to six months at the North Pole itself, a polar night. In the Southern Hemisphere, the situation is exactly reversed, with the South Pole oriented opposite the direction of the North Pole. Six months later, this pole will experience a midnight sun, a day of 24 hours, again reversing with the South Pole.
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By astronomical convention, the four seasons can be determined by the solstices—the points in the orbit of maximum axial tilt toward or away from the Sun—and the equinoxes, when Earth's rotational axis is aligned with its orbital axis. In the Northern Hemisphere, winter solstice currently occurs around 21 December; summer solstice is near 21 June, spring equinox is around 20 March and autumnal equinox is about 22 or 23 September. In the Southern Hemisphere, the situation is reversed, with the summer and winter solstices exchanged and the spring and autumnal equinox dates swapped.[191]
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The angle of Earth's axial tilt is relatively stable over long periods of time. Its axial tilt does undergo nutation; a slight, irregular motion with a main period of 18.6 years.[192] The orientation (rather than the angle) of Earth's axis also changes over time, precessing around in a complete circle over each 25,800 year cycle; this precession is the reason for the difference between a sidereal year and a tropical year. Both of these motions are caused by the varying attraction of the Sun and the Moon on Earth's equatorial bulge. The poles also migrate a few meters across Earth's surface. This polar motion has multiple, cyclical components, which collectively are termed quasiperiodic motion. In addition to an annual component to this motion, there is a 14-month cycle called the Chandler wobble. Earth's rotational velocity also varies in a phenomenon known as length-of-day variation.[193]
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In modern times, Earth's perihelion occurs around 3 January, and its aphelion around 4 July. These dates change over time due to precession and other orbital factors, which follow cyclical patterns known as Milankovitch cycles. The changing Earth–Sun distance causes an increase of about 6.9%[n 21] in solar energy reaching Earth at perihelion relative to aphelion. Because the Southern Hemisphere is tilted toward the Sun at about the same time that Earth reaches the closest approach to the Sun, the Southern Hemisphere receives slightly more energy from the Sun than does the northern over the course of a year. This effect is much less significant than the total energy change due to the axial tilt, and most of the excess energy is absorbed by the higher proportion of water in the Southern Hemisphere.[194]
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A study from 2016 suggested that Planet Nine tilted all the planets of the Solar System, including Earth, by about six degrees.[195]
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A planet that can sustain life is termed habitable, even if life did not originate there. Earth provides liquid water—an environment where complex organic molecules can assemble and interact, and sufficient energy to sustain metabolism.[196] The distance of Earth from the Sun, as well as its orbital eccentricity, rate of rotation, axial tilt, geological history, sustaining atmosphere, and magnetic field all contribute to the current climatic conditions at the surface.[197]
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A planet's life forms inhabit ecosystems, whose total is sometimes said to form a "biosphere".[198] Earth's biosphere is thought to have begun evolving about 3.5 Gya.[65] The biosphere is divided into a number of biomes, inhabited by broadly similar plants and animals.[199] On land, biomes are separated primarily by differences in latitude, height above sea level and humidity. Terrestrial biomes lying within the Arctic or Antarctic Circles, at high altitudes or in extremely arid areas are relatively barren of plant and animal life; species diversity reaches a peak in humid lowlands at equatorial latitudes.[200]
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In July 2016, scientists reported identifying a set of 355 genes from the last universal common ancestor (LUCA) of all organisms living on Earth.[201]
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Earth has resources that have been exploited by humans.[203] Those termed non-renewable resources, such as fossil fuels, only renew over geological timescales.[204]
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Large deposits of fossil fuels are obtained from Earth's crust, consisting of coal, petroleum, and natural gas.[205] These deposits are used by humans both for energy production and as feedstock for chemical production.[206] Mineral ore bodies have also been formed within the crust through a process of ore genesis, resulting from actions of magmatism, erosion, and plate tectonics.[207] These bodies form concentrated sources for many metals and other useful elements.
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Earth's biosphere produces many useful biological products for humans, including food, wood, pharmaceuticals, oxygen, and the recycling of many organic wastes. The land-based ecosystem depends upon topsoil and fresh water, and the oceanic ecosystem depends upon dissolved nutrients washed down from the land.[208] In 1980, 50.53 million km2 (19.51 million sq mi) of Earth's land surface consisted of forest and woodlands, 67.88 million km2 (26.21 million sq mi) was grasslands and pasture, and 15.01 million km2 (5.80 million sq mi) was cultivated as croplands.[209] The estimated amount of irrigated land in 1993 was 2,481,250 km2 (958,020 sq mi).[13] Humans also live on the land by using building materials to construct shelters.
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Large areas of Earth's surface are subject to extreme weather such as tropical cyclones, hurricanes, or typhoons that dominate life in those areas. From 1980 to 2000, these events caused an average of 11,800 human deaths per year.[210] Many places are subject to earthquakes, landslides, tsunamis, volcanic eruptions, tornadoes, sinkholes, blizzards, floods, droughts, wildfires, and other calamities and disasters.
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Many localized areas are subject to human-made pollution of the air and water, acid rain and toxic substances, loss of vegetation (overgrazing, deforestation, desertification), loss of wildlife, species extinction, soil degradation, soil depletion and erosion.
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There is a scientific consensus linking human activities to global warming due to industrial carbon dioxide emissions. This is predicted to produce changes such as the melting of glaciers and ice sheets, more extreme temperature ranges, significant changes in weather and a global rise in average sea levels.[211]
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Cartography, the study and practice of map-making, and geography, the study of the lands, features, inhabitants and phenomena on Earth, have historically been the disciplines devoted to depicting Earth. Surveying, the determination of locations and distances, and to a lesser extent navigation, the determination of position and direction, have developed alongside cartography and geography, providing and suitably quantifying the requisite information.
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Earth's human population reached approximately seven billion on 31 October 2011.[213] Projections indicate that the world's human population will reach 9.2 billion in 2050.[214] Most of the growth is expected to take place in developing nations. Human population density varies widely around the world, but a majority live in Asia. By 2020, 60% of the world's population is expected to be living in urban, rather than rural, areas.[215]
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68% of the land mass of the world is in the northern hemisphere.[216] Partly due to the predominance of land mass, 90% of humans live in the northern hemisphere.[217]
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It is estimated that one-eighth of Earth's surface is suitable for humans to live on – three-quarters of Earth's surface is covered by oceans, leaving one-quarter as land. Half of that land area is desert (14%),[218] high mountains (27%),[219] or other unsuitable terrains. The northernmost permanent settlement in the world is Alert, on Ellesmere Island in Nunavut, Canada.[220] (82°28′N) The southernmost is the Amundsen–Scott South Pole Station, in Antarctica, almost exactly at the South Pole. (90°S)
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States claim the planet's entire land surface, except for parts of Antarctica and a few other unclaimed areas. Earth has never had a planetwide government, but the United Nations is the leading worldwide intergovernmental organization.
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The first human to orbit Earth was Yuri Gagarin on 12 April 1961.[221] In total, about 487 people have visited outer space and reached orbit as of 30 July 2010[update], and, of these, twelve have walked on the Moon.[222][223][224] Normally, the only humans in space are those on the International Space Station. The station's crew, made up of six people, is usually replaced every six months.[225] The farthest that humans have traveled from Earth is 400,171 km (248,655 mi), achieved during the Apollo 13 mission in 1970.[226]
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The Moon is a relatively large, terrestrial, planet-like natural satellite, with a diameter about one-quarter of Earth's. It is the largest moon in the Solar System relative to the size of its planet, although Charon is larger relative to the dwarf planet Pluto. The natural satellites of other planets are also referred to as "moons", after Earth's.
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The gravitational attraction between Earth and the Moon causes tides on Earth. The same effect on the Moon has led to its tidal locking: its rotation period is the same as the time it takes to orbit Earth. As a result, it always presents the same face to the planet. As the Moon orbits Earth, different parts of its face are illuminated by the Sun, leading to the lunar phases; the dark part of the face is separated from the light part by the solar terminator.
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Due to their tidal interaction, the Moon recedes from Earth at the rate of approximately 38 mm/a (1.5 in/year). Over millions of years, these tiny modifications—and the lengthening of Earth's day by about 23 µs/yr—add up to significant changes.[227] During the Devonian period, for example, (approximately 410 Mya) there were 400 days in a year, with each day lasting 21.8 hours.[228]
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The Moon may have dramatically affected the development of life by moderating the planet's climate. Paleontological evidence and computer simulations show that Earth's axial tilt is stabilized by tidal interactions with the Moon.[229] Some theorists think that without this stabilization against the torques applied by the Sun and planets to Earth's equatorial bulge, the rotational axis might be chaotically unstable, exhibiting chaotic changes over millions of years, as appears to be the case for Mars.[230]
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Viewed from Earth, the Moon is just far enough away to have almost the same apparent-sized disk as the Sun. The angular size (or solid angle) of these two bodies match because, although the Sun's diameter is about 400 times as large as the Moon's, it is also 400 times more distant.[186] This allows total and annular solar eclipses to occur on Earth.
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The most widely accepted theory of the Moon's origin, the giant-impact hypothesis, states that it formed from the collision of a Mars-size protoplanet called Theia with the early Earth. This hypothesis explains (among other things) the Moon's relative lack of iron and volatile elements and the fact that its composition is nearly identical to that of Earth's crust.[48]
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Earth has at least five co-orbital asteroids, including 3753 Cruithne and 2002 AA29.[231][232] A trojan asteroid companion, 2010 TK7, is librating around the leading Lagrange triangular point, L4, in Earth's orbit around the Sun.[233][234]
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The tiny near-Earth asteroid 2006 RH120 makes close approaches to the Earth–Moon system roughly every twenty years. During these approaches, it can orbit Earth for brief periods of time.[235]
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As of April 2018[update], there are 1,886 operational, human-made satellites orbiting Earth.[5] There are also inoperative satellites, including Vanguard 1, the oldest satellite currently in orbit, and over 16,000 pieces of tracked space debris.[n 3] Earth's largest artificial satellite is the International Space Station.
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The standard astronomical symbol of Earth consists of a cross circumscribed by a circle, ,[236] representing the four corners of the world.
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Human cultures have developed many views of the planet.[237] Earth is sometimes personified as a deity. In many cultures it is a mother goddess that is also the primary fertility deity,[238] and by the mid-20th century, the Gaia Principle compared Earth's environments and life as a single self-regulating organism leading to broad stabilization of the conditions of habitability.[239][240][241] Creation myths in many religions involve the creation of Earth by a supernatural deity or deities.[238]
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The Hindu Vedas (1500–900 BC) refer to the Earth as Bhūgola (भूगोल), which comes from Bhū (earth, ground) and Gola (ball, sphere, globe). It means the "globe of earth". There is no direct evidence that the Hindus of that time knew that the Earth was sphere-shaped, but this name has been used extensively since the inception of the Vedas.[citation needed]
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Scientific investigation has resulted in several culturally transformative shifts in people's view of the planet. Initial belief in a flat Earth was gradually displaced in the Greek colonies of southern Italy during the late 6th century BC by the idea of spherical Earth,[242][243][244] which was attributed to both the philosophers Pythagoras and Parmenides.[243][244] By the end of the 5th century BC, the sphericity of Earth was universally accepted among Greek intellectuals.[245] Earth was generally believed to be the center of the universe until the 16th century, when scientists first conclusively demonstrated that it was a moving object, comparable to the other planets in the Solar System.[246] Due to the efforts of influential Christian scholars and clerics such as James Ussher, who sought to determine the age of Earth through analysis of genealogies in Scripture, Westerners before the 19th century generally believed Earth to be a few thousand years old at most. It was only during the 19th century that geologists realized Earth's age was at least many millions of years.[247]
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Lord Kelvin used thermodynamics to estimate the age of Earth to be between 20 million and 400 million years in 1864, sparking a vigorous debate on the subject; it was only when radioactivity and radioactive dating were discovered in the late 19th and early 20th centuries that a reliable mechanism for determining Earth's age was established, proving the planet to be billions of years old.[248][249] The perception of Earth shifted again[further explanation needed] in the 20th century when humans first viewed it from orbit, and especially with photographs of Earth returned by the Apollo program.[250][251][252]
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Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
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Baseball is a bat-and-ball game played between two opposing teams who take turns batting and fielding. The game proceeds when a player on the fielding team, called the pitcher, throws a ball which a player on the batting team tries to hit with a bat. The objective of the offensive team (batting team) is to hit the ball into the field of play, allowing its players to run the bases, having them advance counter-clockwise around four bases to score what are called "runs". The objective of the defensive team (fielding team) is to prevent batters from becoming runners, and to prevent runners' advance around the bases.[2] A run is scored when a runner legally advances around the bases in order and touches home plate (the place where the player started as a batter). The team that scores the most runs by the end of the game is the winner.
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The first objective of the batting team is to have a player reach first base safely. A player on the batting team who reaches first base without being called "out" can attempt to advance to subsequent bases as a runner, either immediately or during teammates' turns batting. The fielding team tries to prevent runs by getting batters or runners "out", which forces them out of the field of play. Both the pitcher and fielders have methods of getting the batting team's players out. The opposing teams switch back and forth between batting and fielding; the batting team's turn to bat is over once the fielding team records three outs. One turn batting for each team constitutes an inning. A game is usually composed of nine innings, and the team with the greater number of runs at the end of the game wins. If scores are tied at the end of nine innings, extra innings are usually played. Baseball has no game clock, although most games end in the ninth inning.
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Baseball evolved from older bat-and-ball games already being played in England by the mid-18th century. This game was brought by immigrants to North America, where the modern version developed. By the late 19th century, baseball was widely recognized as the national sport of the United States. Baseball is popular in North America and parts of Central and South America, the Caribbean, and East Asia, particularly in Japan, South Korea and Taiwan.
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In the United States and Canada, professional Major League Baseball (MLB) teams are divided into the National League (NL) and American League (AL), each with three divisions: East, West, and Central. The MLB champion is determined by playoffs that culminate in the World Series. The top level of play is similarly split in Japan between the Central and Pacific Leagues and in Cuba between the West League and East League. The World Baseball Classic, organized by the World Baseball Softball Confederation, is the major international competition of the sport and attracts the top national teams from around the world.
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A baseball game is played between two teams, each composed of nine players, that take turns playing offense (batting and baserunning) and defense (pitching and fielding). A pair of turns, one at bat and one in the field, by each team constitutes an inning. A game consists of nine innings (seven innings at the high school level and in doubleheaders in college and minor leagues, and six innings at the Little League level).[3] One team—customarily the visiting team—bats in the top, or first half, of every inning. The other team—customarily the home team—bats in the bottom, or second half, of every inning. The goal of the game is to score more points (runs) than the other team. The players on the team at bat attempt to score runs by touching all four bases, in order, set at the corners of the square-shaped baseball diamond. A player bats at home plate and must attempt to safely reach a base before proceeding, counterclockwise, from first base, to second base, third base, and back home to score a run. The team in the field attempts to prevent runs from scoring by recording outs, which remove opposing players from offensive action, until their next turn at bat comes up again. When three outs are recorded, the teams switch roles for the next half-inning. If the score of the game is tied after nine innings, extra innings are played to resolve the contest. Many amateur games, particularly unorganized ones, involve different numbers of players and innings.[4]
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The game is played on a field whose primary boundaries, the foul lines, extend forward from home plate at 45-degree angles. The 90-degree area within the foul lines is referred to as fair territory; the 270-degree area outside them is foul territory. The part of the field enclosed by the bases and several yards beyond them is the infield; the area farther beyond the infield is the outfield. In the middle of the infield is a raised pitcher's mound, with a rectangular rubber plate (the rubber) at its center. The outer boundary of the outfield is typically demarcated by a raised fence, which may be of any material and height. The fair territory between home plate and the outfield boundary is baseball's field of play, though significant events can take place in foul territory, as well.[5]
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There are three basic tools of baseball: the ball, the bat, and the glove or mitt:
|
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Protective helmets are also standard equipment for all batters.[9]
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At the beginning of each half-inning, the nine players of the fielding team arrange themselves around the field. One of them, the pitcher, stands on the pitcher's mound. The pitcher begins the pitching delivery with one foot on the rubber, pushing off it to gain velocity when throwing toward home plate. Another fielding team player, the catcher, squats on the far side of home plate, facing the pitcher. The rest of the fielding team faces home plate, typically arranged as four infielders—who set up along or within a few yards outside the imaginary lines (basepaths) between first, second, and third base—and three outfielders. In the standard arrangement, there is a first baseman positioned several steps to the left of first base, a second baseman to the right of second base, a shortstop to the left of second base, and a third baseman to the right of third base. The basic outfield positions are left fielder, center fielder, and right fielder. With the exception of the catcher, all fielders are required to be in fair territory when the pitch is delivered. A neutral umpire sets up behind the catcher.[10] Other umpires will be distributed around the field as well.[11]
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Play starts with a member of the batting team, the batter, standing at home plate, holding a bat.[12] The batter waits for the pitcher to throw a pitch (the ball) toward home plate, and attempts to hit the ball[13] with the bat.[12] The catcher catches pitches that the batter does not hit—as a result of either electing not to swing or failing to connect—and returns them to the pitcher. A batter who hits the ball into the field of play must drop the bat and begin running toward first base, at which point the player is referred to as a runner (or, until the play is over, a batter-runner). A batter-runner who reaches first base without being put out is said to be safe and is on base. A batter-runner may choose to remain at first base or attempt to advance to second base or even beyond—however far the player believes can be reached safely. A player who reaches base despite proper play by the fielders has recorded a hit. A player who reaches first base safely on a hit is credited with a single. If a player makes it to second base safely as a direct result of a hit, it is a double; third base, a triple. If the ball is hit in the air within the foul lines over the entire outfield (and outfield fence, if there is one), or otherwise safely circles all the bases, it is a home run: the batter and any runners on base may all freely circle the bases, each scoring a run. This is the most desirable result for the batter. A player who reaches base due to a fielding mistake is not credited with a hit—instead, the responsible fielder is charged with an error.[12]
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Any runners already on base may attempt to advance on batted balls that land, or contact the ground, in fair territory, before or after the ball lands. A runner on first base must attempt to advance if a ball lands in play. If a ball hit into play rolls foul before passing through the infield, it becomes dead and any runners must return to the base they occupied when the play began. If the ball is hit in the air and caught before it lands, the batter has flied out and any runners on base may attempt to advance only if they tag up (contact the base they occupied when the play began, as or after the ball is caught). Runners may also attempt to advance to the next base while the pitcher is in the process of delivering the ball to home plate; a successful effort is a stolen base.[14]
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A pitch that is not hit into the field of play is called either a strike or a ball. A batter against whom three strikes are recorded strikes out. A batter against whom four balls are recorded is awarded a base on balls or walk, a free advance to first base. (A batter may also freely advance to first base if the batter's body or uniform is struck by a pitch outside the strike zone, provided the batter does not swing and attempts to avoid being hit.)[15] Crucial to determining balls and strikes is the umpire's judgment as to whether a pitch has passed through the strike zone, a conceptual area above home plate extending from the midpoint between the batter's shoulders and belt down to the hollow of the knee.[16]
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While the team at bat is trying to score runs, the team in the field is attempting to record outs. In addition to the strikeout, common ways a member of the batting team may be put out include the flyout, ground out, force out, and tag out. It is possible to record two outs in the course of the same play. This is called a double play. Three outs in one play, a triple play, is possible, though rare. Players put out or retired must leave the field, returning to their team's dugout or bench. A runner may be stranded on base when a third out is recorded against another player on the team. Stranded runners do not benefit the team in its next turn at bat as every half-inning begins with the bases empty.[17]
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An individual player's turn batting or plate appearance is complete when the player reaches base, hits a home run, makes an out, or hits a ball that results in the team's third out, even if it is recorded against a teammate. On rare occasions, a batter may be at the plate when, without the batter's hitting the ball, a third out is recorded against a teammate—for instance, a runner getting caught stealing (tagged out attempting to steal a base). A batter with this sort of incomplete plate appearance starts off the team's next turn batting; any balls or strikes recorded against the batter the previous inning are erased. A runner may circle the bases only once per plate appearance and thus can score at most a single run per batting turn. Once a player has completed a plate appearance, that player may not bat again until the eight other members of the player's team have all taken their turn at bat in the batting order. The batting order is set before the game begins, and may not be altered except for substitutions. Once a player has been removed for a substitute, that player may not reenter the game. Children's games often have more lenient rules, such as Little League rules, which allow players to be substituted back into the same game.[18][3]
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If the designated hitter (DH) rule is in effect, each team has a tenth player whose sole responsibility is to bat (and run). The DH takes the place of another player—almost invariably the pitcher—in the batting order, but does not field. Thus, even with the DH, each team still has a batting order of nine players and a fielding arrangement of nine players.[19]
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The number of players on a baseball roster, or squad, varies by league and by the level of organized play. A Major League Baseball (MLB) team has a roster of 25 players with specific roles. A typical roster features the following players:[20]
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Most baseball leagues worldwide have the DH rule, including MLB's American League, Japan's Pacific League, and Caribbean professional leagues, along with major American amateur organizations.[21] The Central League in Japan and the National League do not have the rule, and high-level minor league clubs connected to National League teams are not required to field a DH.[22] In leagues that apply the designated hitter rule, a typical team has nine offensive regulars (including the DH), five starting pitchers,[23] seven or eight relievers, a backup catcher, and two or three other reserve players.[24][25]
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The manager, or head coach, oversees the team's major strategic decisions, such as establishing the starting rotation, setting the lineup, or batting order, before each game, and making substitutions during games—in particular, bringing in relief pitchers. Managers are typically assisted by two or more coaches; they may have specialized responsibilities, such as working with players on hitting, fielding, pitching, or strength and conditioning. At most levels of organized play, two coaches are stationed on the field when the team is at bat: the first base coach and third base coach, whom occupy designated coaches' boxes, just outside the foul lines. These coaches assist in the direction of baserunners, when the ball is in play, and relay tactical signals from the manager to batters and runners, during pauses in play.[26] In contrast to many other team sports, baseball managers and coaches generally wear their team's uniforms; coaches must be in uniform to be allowed on the field to confer with players during a game.[27]
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Any baseball game involves one or more umpires, who make rulings on the outcome of each play. At a minimum, one umpire will stand behind the catcher, to have a good view of the strike zone, and call balls and strikes. Additional umpires may be stationed near the other bases, thus making it easier to judge plays such as attempted force outs and tag outs. In MLB, four umpires are used for each game, one near each base. In the playoffs, six umpires are used: one at each base and two in the outfield along the foul lines.[28]
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Many of the pre-game and in-game strategic decisions in baseball revolve around a fundamental fact: in general, right-handed batters tend to be more successful against left-handed pitchers and, to an even greater degree, left-handed batters tend to be more successful against right-handed pitchers.[29] A manager with several left-handed batters in the regular lineup, who knows the team will be facing a left-handed starting pitcher, may respond by starting one or more of the right-handed backups on the team's roster. During the late innings of a game, as relief pitchers and pinch hitters are brought in, the opposing managers will often go back and forth trying to create favorable matchups with their substitutions. The manager of the fielding team trying to arrange same-handed pitcher-batter matchups and the manager of the batting team trying to arrange opposite-handed matchups. With a team that has the lead in the late innings, a manager may remove a starting position player—especially one whose turn at bat is not likely to come up again—for a more skillful fielder (known as a defensive substitution).[30]
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The tactical decision that precedes almost every play in a baseball game involves pitch selection.[31] By gripping and then releasing the baseball in a certain manner, and by throwing it at a certain speed, pitchers can cause the baseball to break to either side, or downward, as it approaches the batter; thus creating differing pitches that can be selected.[32] Among the resulting wide variety of pitches that may be thrown, the four basic types are the fastball, the changeup (or off-speed pitch), and two breaking balls—the curveball and the slider.[33] Pitchers have different repertoires of pitches they are skillful at throwing. Conventionally, before each pitch, the catcher signals the pitcher what type of pitch to throw, as well as its general vertical and/or horizontal location.[34] If there is disagreement on the selection, the pitcher may shake off the sign and the catcher will call for a different pitch.
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With a runner on base and taking a lead, the pitcher may attempt a pickoff, a quick throw to a fielder covering the base to keep the runner's lead in check or, optimally, effect a tag out.[35] Pickoff attempts, however, are subject to rules that severely restrict the pitcher's movements before and during the pickoff attempt. Violation of any one of these rules could result in the umpire calling a balk against the pitcher, which permits any runners on base to advance one base with impunity.[36] If an attempted stolen base is anticipated, the catcher may call for a pitchout, a ball thrown deliberately off the plate, allowing the catcher to catch it while standing and throw quickly to a base.[37] Facing a batter with a strong tendency to hit to one side of the field, the fielding team may employ a shift, with most or all of the fielders moving to the left or right of their usual positions. With a runner on third base, the infielders may play in, moving closer to home plate to improve the odds of throwing out the runner on a ground ball, though a sharply hit grounder is more likely to carry through a drawn-in infield.[38]
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Several basic offensive tactics come into play with a runner on first base, including the fundamental choice of whether to attempt a steal of second base. The hit and run is sometimes employed, with a skillful contact hitter, the runner takes off with the pitch, drawing the shortstop or second baseman over to second base, creating a gap in the infield for the batter to poke the ball through.[39] The sacrifice bunt, calls for the batter to focus on making soft contact with the ball, so that it rolls a short distance into the infield, allowing the runner to advance into scoring position as the batter is thrown out at first. A batter, particularly one who is a fast runner, may also attempt to bunt for a hit. A sacrifice bunt employed with a runner on third base, aimed at bringing that runner home, is known as a squeeze play.[40] With a runner on third and fewer than two outs, a batter may instead concentrate on hitting a fly ball that, even if it is caught, will be deep enough to allow the runner to tag up and score—a successful batter, in this case, gets credit for a sacrifice fly.[38] The manager will sometimes signal a batter who is ahead in the count (i.e., has more balls than strikes) to take, or not swing at, the next pitch.[41]
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The evolution of baseball from older bat-and-ball games is difficult to trace with precision. Consensus once held that today's baseball is a North American development from the older game rounders, popular among children in Great Britain and Ireland.[42][43][44] Baseball Before We Knew It: A Search for the Roots of the Game (2005), by American baseball historian David Block, suggests that the game originated in England; recently uncovered historical evidence supports this position. Block argues that rounders and early baseball were actually regional variants of each other, and that the game's most direct antecedents are the English games of stoolball and "tut-ball".[42] The earliest known reference to baseball is in a 1744 British publication, A Little Pretty Pocket-Book, by John Newbery.[45] Block discovered that the first recorded game of "Bass-Ball" took place in 1749 in Surrey, and featured the Prince of Wales as a player.[46] This early form of the game was apparently brought to Canada by English immigrants.[47]
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By the early 1830s, there were reports of a variety of uncodified bat-and-ball games recognizable as early forms of baseball being played around North America.[48] The first officially recorded baseball game on this continent was played in Beachville, Ontario, Canada, on June 4, 1838.[49] In 1845, Alexander Cartwright, a member of New York City's Knickerbocker Club, led the codification of the so-called Knickerbocker Rules,[50] which in turn were based on rules developed in 1837 by William R. Wheaton of the Gotham Club.[51] While there are reports that the New York Knickerbockers played games in 1845, the contest long recognized as the first officially recorded baseball game in U.S. history took place on June 19, 1846, in Hoboken, New Jersey: the "New York Nine" defeated the Knickerbockers, 23–1, in four innings.[52] With the Knickerbocker code as the basis, the rules of modern baseball continued to evolve over the next half-century.[53]
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In the mid-1850s, a baseball craze hit the New York metropolitan area,[54] and by 1856, local journals were referring to baseball as the "national pastime" or "national game".[55] A year later, the sport's first governing body, the National Association of Base Ball Players, was formed. In 1867, it barred participation by African Americans.[56] The more formally structured National League was founded in 1876.[57] Professional Negro leagues formed, but quickly folded.[58] In 1887, softball, under the name of indoor baseball or indoor-outdoor, was invented as a winter version of the parent game.[59] The National League's first successful counterpart, the American League, which evolved from the minor Western League, was established in 1893, and virtually all of the modern baseball rules were in place by then.[60][61]
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The National Agreement of 1903 formalized relations both between the two major leagues and between them and the National Association of Professional Base Ball Leagues, representing most of the country's minor professional leagues.[62] The World Series, pitting the two major league champions against each other, was inaugurated that fall.[63] The Black Sox Scandal of the 1919 World Series led to the formation of a new National Commission of baseball that drew the two major leagues closer together.[64] The first major league baseball commissioner, Kenesaw Mountain Landis, was elected in 1920. That year also saw the founding of the Negro National League; the first significant Negro league, it would operate until 1931. For part of the 1920s, it was joined by the Eastern Colored League.[65]
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Compared with the present, professional baseball in the early 20th century was lower-scoring, and pitchers were more dominant.[66] The so-called dead-ball era ended in the early 1920s with several changes in rule and circumstance that were advantageous to hitters. Strict new regulations governed the ball's size, shape and composition, along with a new rule officially banning the spitball and other pitches that depended on the ball being treated or roughed-up with foreign substances, resulted in a ball that traveled farther when hit.[67] The rise of the legendary player Babe Ruth, the first great power hitter of the new era, helped permanently alter the nature of the game.[68] In the late 1920s and early 1930s, St. Louis Cardinals general manager Branch Rickey invested in several minor league clubs and developed the first modern farm system.[69] A new Negro National League was organized in 1933; four years later, it was joined by the Negro American League. The first elections to the National Baseball Hall of Fame took place in 1936. In 1939, Little League Baseball was founded in Pennsylvania.[70]
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A large number of minor league teams disbanded when World War II led to a player shortage. Chicago Cubs owner Philip K. Wrigley led the formation of the All-American Girls Professional Baseball League to help keep the game in the public eye.[71] The first crack in the unwritten agreement barring blacks from white-controlled professional ball occurred in 1945: Jackie Robinson was signed by the National League's Brooklyn Dodgers and began playing for their minor league team in Montreal.[72] In 1947, Robinson broke the major leagues' color barrier when he debuted with the Dodgers.[73] Latin American players, largely overlooked before, also started entering the majors in greater numbers. In 1951, two Chicago White Sox, Venezuelan-born Chico Carrasquel and black Cuban-born Minnie Miñoso, became the first Hispanic All-Stars.[74][75] Integration proceeded slowly: by 1953, only six of the 16 major league teams had a black player on the roster.[74]
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In 1975, the union's power—and players' salaries—began to increase greatly when the reserve clause was effectively struck down, leading to the free agency system.[76] Significant work stoppages occurred in 1981 and 1994, the latter forcing the cancellation of the World Series for the first time in 90 years.[77] Attendance had been growing steadily since the mid-1970s and in 1994, before the stoppage, the majors were setting their all-time record for per-game attendance.[78][79] After play resumed in 1995, non-division-winning wild card teams became a permanent fixture of the post-season. Regular-season interleague play was introduced in 1997 and the second-highest attendance mark for a full season was set.[80] In 2000, the National and American Leagues were dissolved as legal entities. While their identities were maintained for scheduling purposes (and the designated hitter distinction), the regulations and other functions—such as player discipline and umpire supervision—they had administered separately were consolidated under the rubric of MLB.[81]
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In 2001, Barry Bonds established the current record of 73 home runs in a single season. There had long been suspicions that the dramatic increase in power hitting was fueled in large part by the abuse of illegal steroids (as well as by the dilution of pitching talent due to expansion), but the issue only began attracting significant media attention in 2002 and there was no penalty for the use of performance-enhancing drugs before 2004.[82] In 2007, Bonds became MLB's all-time home run leader, surpassing Hank Aaron, as total major league and minor league attendance both reached all-time highs.[83][84]
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Widely known as America's pastime, baseball is well established in several other countries as well. As early as 1877, a professional league, the International Association, featured teams from both Canada and the US.[85] While baseball is widely played in Canada and many minor league teams have been based in the country,[86][87] the American major leagues did not include a Canadian club until 1969, when the Montreal Expos joined the National League as an expansion team. In 1977, the expansion Toronto Blue Jays joined the American League.[88]
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In 1847, American soldiers played what may have been the first baseball game in Mexico at Parque Los Berros in Xalapa, Veracruz.[89] The first formal baseball league outside of the United States and Canada was founded in 1878 in Cuba, which maintains a rich baseball tradition. The Dominican Republic held its first islandwide championship tournament in 1912.[90] Professional baseball tournaments and leagues began to form in other countries between the world wars, including the Netherlands (formed in 1922), Australia (1934), Japan (1936), Mexico (1937), and Puerto Rico (1938).[91] The Japanese major leagues have long been considered the highest quality professional circuits outside of the United States.[92]
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After World War II, professional leagues were founded in many Latin American countries, most prominently Venezuela (1946) and the Dominican Republic (1955).[93] Since the early 1970s, the annual Caribbean Series has matched the championship clubs from the four leading Latin American winter leagues: the Dominican Professional Baseball League, Mexican Pacific League, Puerto Rican Professional Baseball League, and Venezuelan Professional Baseball League. In Asia, South Korea (1982), Taiwan (1990) and China (2003) all have professional leagues.[94]
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Many European countries have professional leagues as well; the most successful, other than the Dutch league, is the Italian league, founded in 1948.[95] In 2004, Australia won a surprise silver medal at the Olympic Games.[96] The Confédération Européene de Baseball (European Baseball Confederation), founded in 1953, organizes a number of competitions between clubs from different countries. Other competitions between national teams, such as the Baseball World Cup and the Olympic baseball tournament, were administered by the International Baseball Federation (IBAF) from its formation in 1938 until its 2013 merger with the International Softball Federation to create the current joint governing body for both sports, the World Baseball Softball Confederation (WBSC).[97] Women's baseball is played on an organized amateur basis in numerous countries.[98]
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After being admitted to the Olympics as a medal sport beginning with the 1992 Games, baseball was dropped from the 2012 Summer Olympic Games at the 2005 International Olympic Committee meeting. It remained part of the 2008 Games.[99] While the sport's lack of a following in much of the world was a factor,[100] more important was MLB's reluctance to allow its players to participate during the major league season.[101] MLB initiated the World Baseball Classic, scheduled to precede its season, partly as a replacement, high-profile international tournament. The inaugural Classic, held in March 2006, was the first tournament involving national teams to feature a significant number of MLB participants.[102][103] The Baseball World Cup was discontinued after its 2011 edition in favor of an expanded World Baseball Classic.[104]
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Baseball has certain attributes that set it apart from the other popular team sports in the countries where it has a following. All of these sports use a clock,[105] play is less individual,[106] and the variation between playing fields is not as substantial or important.[107] The comparison between cricket and baseball demonstrates that many of baseball's distinctive elements are shared in various ways with its cousin sports.[108]
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In clock-limited sports, games often end with a team that holds the lead killing the clock rather than competing aggressively against the opposing team. In contrast, baseball has no clock, thus a team cannot win without getting the last batter out and rallies are not constrained by time. At almost any turn in any baseball game, the most advantageous strategy is some form of aggressive strategy.[109] Whereas, in the case of multi-day Test and first-class cricket, the possibility of a draw often encourages a team that is batting last and well behind, to bat defensively and run out the clock, giving up any faint chance at a win, to avoid an overall loss.[110]
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While nine innings has been the standard since the beginning of professional baseball, the duration of the average major league game has increased steadily through the years. At the turn of the 20th century, games typically took an hour and a half to play. In the 1920s, they averaged just less than two hours, which eventually ballooned to 2:38 in 1960.[111] By 1997, the average American League game lasted 2:57 (National League games were about 10 minutes shorter—pitchers at the plate making for quicker outs than designated hitters).[112] In 2004, Major League Baseball declared that its goal was an average game of 2:45.[111] By 2014, though, the average MLB game took over three hours to complete.[113] The lengthening of games is attributed to longer breaks between half-innings for television commercials, increased offense, more pitching changes, and a slower pace of play with pitchers taking more time between each delivery, and batters stepping out of the box more frequently.[111][112] Other leagues have experienced similar issues. In 2008, Nippon Professional Baseball took steps aimed at shortening games by 12 minutes from the preceding decade's average of 3:18.[114]
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In 2016, the average nine-inning playoff game in Major League baseball was 3 hours and 35 minutes. This was up 10 minutes from 2015 and 21 minutes from 2014.[115]
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Although baseball is a team sport, individual players are often placed under scrutiny and pressure. In 1915, a baseball instructional manual pointed out that every single pitch, of which there are often more than two hundred in a game, involves an individual, one-on-one contest: "the pitcher and the batter in a battle of wits".[116] Pitcher, batter, and fielder all act essentially independent of each other. While coaching staffs can signal pitcher or batter to pursue certain tactics, the execution of the play itself is a series of solitary acts. If the batter hits a line drive, the outfielder is solely responsible for deciding to try to catch it or play it on the bounce and for succeeding or failing. The statistical precision of baseball is both facilitated by this isolation and reinforces it.
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Cricket is more similar to baseball than many other team sports in this regard: while the individual focus in cricket is mitigated by the importance of the batting partnership and the practicalities of tandem running, it is enhanced by the fact that a batsman may occupy the wicket for an hour or much more. There is no statistical equivalent in cricket for the fielding error and thus less emphasis on personal responsibility in this area of play.[117]
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Unlike those of most sports, baseball playing fields can vary significantly in size and shape. While the dimensions of the infield are specifically regulated, the only constraint on outfield size and shape for professional teams, following the rules of MLB and Minor League Baseball, is that fields built or remodeled since June 1, 1958, must have a minimum distance of 325 feet (99 m) from home plate to the fences in left and right field and 400 feet (122 m) to center.[118] Major league teams often skirt even this rule. For example, at Minute Maid Park, which became the home of the Houston Astros in 2000, the Crawford Boxes in left field are only 315 feet (96 m) from home plate.[119] There are no rules at all that address the height of fences or other structures at the edge of the outfield. The most famously idiosyncratic outfield boundary is the left-field wall at Boston's Fenway Park, in use since 1912: the Green Monster is 310 feet (94 m) from home plate down the line and 37 feet (11 m) tall.[120]
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Similarly, there are no regulations at all concerning the dimensions of foul territory. Thus a foul fly ball may be entirely out of play in a park with little space between the foul lines and the stands, but a foulout in a park with more expansive foul ground.[121] A fence in foul territory that is close to the outfield line will tend to direct balls that strike it back toward the fielders, while one that is farther away may actually prompt more collisions, as outfielders run full speed to field balls deep in the corner. These variations can make the difference between a double and a triple or inside-the-park home run.[122] The surface of the field is also unregulated. While the adjacent image shows a traditional field surfacing arrangement (and the one used by virtually all MLB teams with naturally surfaced fields), teams are free to decide what areas will be grassed or bare.[123] Some fields—including several in MLB—use an artificial surface, such as AstroTurf. Surface variations can have a significant effect on how ground balls behave and are fielded as well as on baserunning. Similarly, the presence of a roof (seven major league teams play in stadiums with permanent or retractable roofs) can greatly affect how fly balls are played.[124] While football and soccer players deal with similar variations of field surface and stadium covering, the size and shape of their fields are much more standardized. The area out-of-bounds on a football or soccer field does not affect play the way foul territory in baseball does, so variations in that regard are largely insignificant.[125]
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These physical variations create a distinctive set of playing conditions at each ballpark. Other local factors, such as altitude and climate, can also significantly affect play. A given stadium may acquire a reputation as a pitcher's park or a hitter's park, if one or the other discipline notably benefits from its unique mix of elements. The most exceptional park in this regard is Coors Field, home of the Colorado Rockies. Its high altitude—5,282 feet (1,610 m) above sea level—is partly responsible for giving it the strongest hitter's park effect in the major leagues due to the low air pressure.[126] Wrigley Field, home of the Chicago Cubs, is known for its fickle disposition: a hitter's park when the strong winds off Lake Michigan are blowing out, it becomes more of a pitcher's park when they are blowing in.[127] The absence of a standardized field affects not only how particular games play out, but the nature of team rosters and players' statistical records. For example, hitting a fly ball 330 feet (100 m) into right field might result in an easy catch on the warning track at one park, and a home run at another. A team that plays in a park with a relatively short right field, such as the New York Yankees, will tend to stock its roster with left-handed pull hitters, who can best exploit it. On the individual level, a player who spends most of his career with a team that plays in a hitter's park will gain an advantage in batting statistics over time—even more so if his talents are especially suited to the park.[128]
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Organized baseball lends itself to statistics to a greater degree than many other sports. Each play is discrete and has a relatively small number of possible outcomes. In the late 19th century, a former cricket player, English-born Henry Chadwick of Brooklyn, was responsible for the "development of the box score, tabular standings, the annual baseball guide, the batting average, and most of the common statistics and tables used to describe baseball."[129] The statistical record is so central to the game's "historical essence" that Chadwick came to be known as Father Baseball.[129] In the 1920s, American newspapers began devoting more and more attention to baseball statistics, initiating what journalist and historian Alan Schwarz describes as a "tectonic shift in sports, as intrigue that once focused mostly on teams began to go to individual players and their statistics lines."[130]
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The Official Baseball Rules administered by MLB require the official scorer to categorize each baseball play unambiguously. The rules provide detailed criteria to promote consistency. The score report is the official basis for both the box score of the game and the relevant statistical records.[131] General managers, managers, and baseball scouts use statistics to evaluate players and make strategic decisions.
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Certain traditional statistics are familiar to most baseball fans. The basic batting statistics include:[132]
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The basic baserunning statistics include:[133]
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The basic pitching statistics include:[134]
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The basic fielding statistics include:[135]
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Among the many other statistics that are kept are those collectively known as situational statistics. For example, statistics can indicate which specific pitchers a certain batter performs best against. If a given situation statistically favors a certain batter, the manager of the fielding team may be more likely to change pitchers or have the pitcher intentionally walk the batter in order to face one who is less likely to succeed.[136]
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Sabermetrics refers to the field of baseball statistical study and the development of new statistics and analytical tools. The term is also used to refer directly to new statistics themselves. The term was coined around 1980 by one of the field's leading proponents, Bill James, and derives from the Society for American Baseball Research (SABR).[137]
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The growing popularity of sabermetrics since the early 1980s has brought more attention to two batting statistics that sabermetricians argue are much better gauges of a batter's skill than batting average:[138]
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Some of the new statistics devised by sabermetricians have gained wide use:
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Writing in 1919, philosopher Morris Raphael Cohen described baseball as America's national religion.[143] In the words of sports columnist Jayson Stark, baseball has long been "a unique paragon of American culture"—a status he sees as devastated by the steroid abuse scandal.[144] Baseball has an important place in other national cultures as well: Scholar Peter Bjarkman describes "how deeply the sport is ingrained in the history and culture of a nation such as Cuba, [and] how thoroughly it was radically reshaped and nativized in Japan."[145] Since the early 1980s, the Dominican Republic, in particular the city of San Pedro de Macorís, has been the major leagues' primary source of foreign talent.[146] In 2017, 83 of the 868 players on MLB Opening Day rosters (and disabled lists) were from the country. Among other Caribbean countries and territories, a combined 97 MLB players were born in Venezuela, Cuba, and Puerto Rico.[147] Hall-of-Famer Roberto Clemente remains one of the greatest national heroes in Puerto Rico's history.[148] While baseball has long been the island's primary athletic pastime, its once well-attended professional winter league has declined in popularity since 1990, when young Puerto Rican players began to be included in the major leagues' annual first-year player draft.[149] In Asia, baseball is among the most popular sports in Japan and South Korea.[150]
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The major league game in the United States was originally targeted toward a middle-class, white-collar audience: relative to other spectator pastimes, the National League's set ticket price of 50 cents in 1876 was high, while the location of playing fields outside the inner city and the workweek daytime scheduling of games were also obstacles to a blue-collar audience.[151] A century later, the situation was very different. With the rise in popularity of other team sports with much higher average ticket prices—football, basketball, and hockey—professional baseball had become among the most blue-collar-oriented of leading American spectator sports.[152]
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Overall, baseball has a large following in the United States; a 2006 poll found that nearly half of Americans are fans.[153] In the late 1900s and early 2000s, baseball's position compared to football in the United States moved in contradictory directions. In 2008, MLB set a revenue record of $6.5 billion, matching the NFL's revenue for the first time in decades.[154] A new MLB revenue record of more than $10 billion was set in 2017.[155] On the other hand, the percentage of American sports fans polled who named baseball as their favorite sport was 9%, compared to pro football at 37%.[156] In 1985, the respective figures were pro football 24%, baseball 23%.[157] Because there are so many more major league games played, there is no comparison in overall attendance.[158] In 2008, total attendance at major league games was the second-highest in history: 78.6 million, 0.7% off the record set the previous year.[83] The following year, amid the U.S. recession, attendance fell by 6.6% to 73.4 million.[159] Eight years later, it dropped under 73 million.[160] Attendance at games held under the Minor League Baseball umbrella set a record in 2008, with 43.3 million.[161] While MLB games have not drawn the same national TV viewership as football games, MLB games are dominant in teams' local markets and regularly lead all programs in primetime in their markets during the summer.[162]
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In Japan, where baseball is inarguably the leading spectator team sport, combined revenue for the twelve teams in Nippon Professional Baseball (NPB), the body that oversees both the Central and Pacific Leagues, was estimated at $1 billion in 2007. Total NPB attendance for the year was approximately 20 million. While in the preceding two decades, MLB attendance grew by 50 percent and revenue nearly tripled, the comparable NPB figures were stagnant. There are concerns that MLB's growing interest in acquiring star Japanese players will hurt the game in their home country.[163] In Cuba, where baseball is by every reckoning the national sport,[164] the national team overshadows the city and provincial teams that play in the top-level domestic leagues.[165] Revenue figures are not released for the country's amateur system. Similarly, according to one official pronouncement, the sport's governing authority "has never taken into account attendance ... because its greatest interest has always been the development of athletes".[166]
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As of 2018[update], Little League Baseball oversees leagues with close to 2.4 million participants in over 80 countries.[167] The number of players has fallen since the 1990s, when 3 million children took part in Little League Baseball annually.[168] Babe Ruth League teams have over 1 million participants.[169] According to the president of the International Baseball Federation, between 300,000 and 500,000 women and girls play baseball around the world, including Little League and the introductory game of Tee Ball.[170]
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A varsity baseball team is an established part of physical education departments at most high schools and colleges in the United States.[171] In 2015, nearly half a million high schoolers and over 34,000 collegians played on their schools' baseball teams.[172] By early in the 20th century, intercollegiate baseball was Japan's leading sport. Today, high school baseball in particular is immensely popular there.[173] The final rounds of the two annual tournaments—the National High School Baseball Invitational Tournament in the spring, and the even more important National High School Baseball Championship in the summer—are broadcast around the country. The tournaments are known, respectively, as Spring Koshien and Summer Koshien after the 55,000-capacity stadium where they are played.[174] In Cuba, baseball is a mandatory part of the state system of physical education, which begins at age six. Talented children as young as seven are sent to special district schools for more intensive training—the first step on a ladder whose acme is the national baseball team.[165]
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Baseball has had a broad impact on popular culture, both in the United States and elsewhere. Dozens of English-language idioms have been derived from baseball; in particular, the game is the source of a number of widely used sexual euphemisms.[176] The first networked radio broadcasts in North America were of the 1922 World Series: famed sportswriter Grantland Rice announced play-by-play from New York City's Polo Grounds on WJZ–Newark, New Jersey, which was connected by wire to WGY–Schenectady, New York, and WBZ–Springfield, Massachusetts.[177] The baseball cap has become a ubiquitous fashion item not only in the United States and Japan, but also in countries where the sport itself is not particularly popular, such as the United Kingdom.[178]
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Baseball has inspired many works of art and entertainment. One of the first major examples, Ernest Thayer's poem "Casey at the Bat", appeared in 1888. A wry description of the failure of a star player in what would now be called a "clutch situation", the poem became the source of vaudeville and other staged performances, audio recordings, film adaptations, and an opera, as well as a host of sequels and parodies in various media. There have been many baseball movies, including the Academy Award–winning The Pride of the Yankees (1942) and the Oscar nominees The Natural (1984) and Field of Dreams (1989). The American Film Institute's selection of the ten best sports movies includes The Pride of the Yankees at number 3 and Bull Durham (1988) at number 5.[179] Baseball has provided thematic material for hits on both stage—the Adler–Ross musical Damn Yankees—and record—George J. Gaskin's "Slide, Kelly, Slide", Simon and Garfunkel's "Mrs. Robinson", and John Fogerty's "Centerfield".[180] The baseball-inspired comedic sketch "Who's on First", popularized by Abbott and Costello in 1938, quickly became famous. Six decades later, Time named it the best comedy routine of the 20th century.[181]
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Literary works connected to the game include the short fiction of Ring Lardner and novels such as Bernard Malamud's The Natural (the source for the movie), Robert Coover's The Universal Baseball Association, Inc., J. Henry Waugh, Prop., and W. P. Kinsella's Shoeless Joe (the source for Field of Dreams). Baseball's literary canon also includes the beat reportage of Damon Runyon; the columns of Grantland Rice, Red Smith, Dick Young, and Peter Gammons; and the essays of Roger Angell. Among the celebrated nonfiction books in the field are Lawrence S. Ritter's The Glory of Their Times, Roger Kahn's The Boys of Summer, and Michael Lewis's Moneyball. The 1970 publication of major league pitcher Jim Bouton's tell-all chronicle Ball Four is considered a turning point in the reporting of professional sports.[182]
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Baseball has also inspired the creation of new cultural forms. Baseball cards were introduced in the late 19th century as trade cards. A typical example featured an image of a baseball player on one side and advertising for a business on the other. In the early 1900s they were produced widely as promotional items by tobacco and confectionery companies. The 1930s saw the popularization of the modern style of baseball card, with a player photograph accompanied on the rear by statistics and biographical data. Baseball cards—many of which are now prized collectibles—are the source of the much broader trading card industry, involving similar products for different sports and non-sports-related fields.[183]
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Modern fantasy sports began in 1980 with the invention of Rotisserie League Baseball by New York writer Daniel Okrent and several friends. Participants in a Rotisserie league draft notional teams from the list of active MLB players and play out an entire imaginary season with game outcomes based on the players' latest real-world statistics. Rotisserie-style play quickly became a phenomenon. Now known more generically as fantasy baseball, it has inspired similar games based on an array of different sports.[184] The field boomed with increasing Internet access and new fantasy sports-related websites. By 2008, 29.9 million people in the United States and Canada were playing fantasy sports, spending $800 million on the hobby.[185] The burgeoning popularity of fantasy baseball is also credited with the increasing attention paid to sabermetrics—first among fans, only later among baseball professionals.[186]
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Earth is the third planet from the Sun and the only astronomical object known to harbor life. According to radiometric dating estimation and other evidence, Earth formed over 4.5 billion years ago. Earth's gravity interacts with other objects in space, especially the Sun and the Moon, which is Earth's only natural satellite. Earth orbits around the Sun in 365.256 solar days, a period known as an Earth sidereal year. During this time, Earth rotates about its axis 366.256 times, that is, a sidereal year has 366.256 sidereal days.[n 6]
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Earth's axis of rotation is tilted with respect to its orbital plane, producing seasons on Earth. The gravitational interaction between Earth and the Moon causes tides, stabilizes Earth's orientation on its axis, and gradually slows its rotation. Earth is the densest planet in the Solar System and the largest and most massive of the four rocky planets.
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Earth's outer layer (lithosphere) is divided into several rigid tectonic plates that migrate across the surface over many millions of years. About 29% of Earth's surface is land consisting of continents and islands. The remaining 71% is covered with water, mostly by oceans but also lakes, rivers and other fresh water, which all together constitute the hydrosphere. The majority of Earth's polar regions are covered in ice, including the Antarctic ice sheet and the sea ice of the Arctic ice pack. Earth's interior remains active with a solid iron inner core, a liquid outer core that generates Earth's magnetic field, and a convecting mantle that drives plate tectonics.
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Within the first billion years of Earth's history, life appeared in the oceans and began to affect Earth's atmosphere and surface, leading to the proliferation of anaerobic and, later, aerobic organisms. Some geological evidence indicates that life may have arisen as early as 4.1 billion years ago. Since then, the combination of Earth's distance from the Sun, physical properties and geological history have allowed life to evolve and thrive. In the history of life on Earth, biodiversity has gone through long periods of expansion, occasionally punctuated by mass extinctions. Over 99% of all species that ever lived on Earth are extinct. Estimates of the number of species on Earth today vary widely; most species have not been described. Over 7.7 billion humans live on Earth and depend on its biosphere and natural resources for their survival.[23]
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The modern English word Earth developed, via Middle English,[n 7] from an Old English noun most often spelled eorðe.[24] It has cognates in every Germanic language, and their ancestral root has been reconstructed as *erþō. In its earliest attestation, the word eorðe was already being used to translate the many senses of Latin terra and Greek γῆ gē: the ground,[n 8] its soil,[n 9] dry land,[n 10] the human world,[n 11] the surface of the world (including the sea),[n 12] and the globe itself.[n 13] As with Roman Terra/Tellūs and Greek Gaia, Earth may have been a personified goddess in Germanic paganism: late Norse mythology included Jörð ('Earth'), a giantess often given as the mother of Thor.[33]
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Originally, earth was written in lowercase, and from early Middle English, its definite sense as "the globe" was expressed as the earth. By Early Modern English, many nouns were capitalized, and the earth became (and often remained) the Earth, particularly when referenced along with other heavenly bodies. More recently, the name is sometimes simply given as Earth, by analogy with the names of the other planets.[24] House styles now vary: Oxford spelling recognizes the lowercase form as the most common, with the capitalized form an acceptable variant. Another convention capitalizes "Earth" when appearing as a name (e.g. "Earth's atmosphere") but writes it in lowercase when preceded by the (e.g. "the atmosphere of the earth"). It almost always appears in lowercase in colloquial expressions such as "what on earth are you doing?"[34]
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Occasionally, the name Terra /ˈtɛrə/ is used in scientific writing and especially in science fiction to distinguish our inhabited planet from others,[35] while in poetry Tellus /ˈtɛləs/ has been used to denote personification of the Earth.[36] The Greek poetic name Gaea (Gæa) /ˈdʒiːə/ is rare, though the alternative spelling Gaia has become common due to the Gaia hypothesis, in which case its pronunciation is /ˈɡaɪə/ rather than the more Classical /ˈɡeɪə/.[37]
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There are a number of adjectives for the planet Earth. From Earth itself comes earthly. From Latin Terra come Terran /ˈtɛrən/,[38] Terrestrial /təˈrɛstriəl/,[39] and (via French) Terrene /təˈriːn/,[40] and from Latin Tellus come Tellurian /tɛˈlʊəriən/[41] and, more rarely, Telluric and Tellural. From Greek Gaia and Gaea comes Gaian and Gaean.
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An inhabitant of the Earth is an Earthling, a Terran, a Terrestrial, a Tellurian or, rarely, an Earthian.
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The oldest material found in the Solar System is dated to 4.5672±0.0006 billion years ago (BYA).[42] By 4.54±0.04 BYA[43] the primordial Earth had formed. The bodies in the Solar System formed and evolved with the Sun. In theory, a solar nebula partitions a volume out of a molecular cloud by gravitational collapse, which begins to spin and flatten into a circumstellar disk, and then the planets grow out of that disk with the Sun. A nebula contains gas, ice grains, and dust (including primordial nuclides). According to nebular theory, planetesimals formed by accretion, with the primordial Earth taking 10–20 million years (Mys) to form.[44]
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A subject of research is the formation of the Moon, some 4.53 BYA.[45] A leading hypothesis is that it was formed by accretion from material loosed from Earth after a Mars-sized object, named Theia, hit Earth.[46] In this view, the mass of Theia was approximately 10 percent of Earth;[47] it hit Earth with a glancing blow and some of its mass merged with Earth.[48] Between approximately 4.1 and 3.8 BYA, numerous asteroid impacts during the Late Heavy Bombardment caused significant changes to the greater surface environment of the Moon and, by inference, to that of Earth.
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Earth's atmosphere and oceans were formed by volcanic activity and outgassing.[49] Water vapor from these sources condensed into the oceans, augmented by water and ice from asteroids, protoplanets, and comets.[50] In this model, atmospheric "greenhouse gases" kept the oceans from freezing when the newly forming Sun had only 70% of its current luminosity.[51] By 3.5 BYA, Earth's magnetic field was established, which helped prevent the atmosphere from being stripped away by the solar wind.[52]
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A crust formed when the molten outer layer of Earth cooled to form a solid. The two models[53] that explain land mass propose either a steady growth to the present-day forms[54] or, more likely, a rapid growth[55] early in Earth history[56] followed by a long-term steady continental area.[57][58][59] Continents formed by plate tectonics, a process ultimately driven by the continuous loss of heat from Earth's interior. Over the period of hundreds of millions of years, the supercontinents have assembled and broken apart. Roughly 750 million years ago (MYA), one of the earliest known supercontinents, Rodinia, began to break apart. The continents later recombined to form Pannotia 600–540 MYA, then finally Pangaea, which also broke apart 180 MYA.[60]
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The present pattern of ice ages began about 40 MYA,[61] and then intensified during the Pleistocene about 3 MYA.[62] High-latitude regions have since undergone repeated cycles of glaciation and thaw, repeating about every 40,000–100,000 years. The last continental glaciation ended 10,000 years ago.[63]
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Chemical reactions led to the first self-replicating molecules about four billion years ago. A half billion years later, the last common ancestor of all current life arose.[64] The evolution of photosynthesis allowed the Sun's energy to be harvested directly by life forms. The resultant molecular oxygen (O2) accumulated in the atmosphere and due to interaction with ultraviolet solar radiation, formed a protective ozone layer (O3) in the upper atmosphere.[65] The incorporation of smaller cells within larger ones resulted in the development of complex cells called eukaryotes.[66] True multicellular organisms formed as cells within colonies became increasingly specialized. Aided by the absorption of harmful ultraviolet radiation by the ozone layer, life colonized Earth's surface.[67] Among the earliest fossil evidence for life is microbial mat fossils found in 3.48 billion-year-old sandstone in Western Australia,[68] biogenic graphite found in 3.7 billion-year-old metasedimentary rocks in Western Greenland,[69] and remains of biotic material found in 4.1 billion-year-old rocks in Western Australia.[70][71] The earliest direct evidence of life on Earth is contained in 3.45 billion-year-old Australian rocks showing fossils of microorganisms.[72][73]
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During the Neoproterozoic, 750 to 580 MYA, much of Earth might have been covered in ice. This hypothesis has been termed "Snowball Earth", and it is of particular interest because it preceded the Cambrian explosion, when multicellular life forms significantly increased in complexity.[74] Following the Cambrian explosion, 535 MYA, there have been five mass extinctions.[75] The most recent such event was 66 MYA, when an asteroid impact triggered the extinction of the non-avian dinosaurs and other large reptiles, but spared some small animals such as mammals, which at the time resembled shrews. Mammalian life has diversified over the past 66 Mys, and several million years ago an African ape-like animal such as Orrorin tugenensis gained the ability to stand upright.[76] This facilitated tool use and encouraged communication that provided the nutrition and stimulation needed for a larger brain, which led to the evolution of humans. The development of agriculture, and then civilization, led to humans having an influence on Earth and the nature and quantity of other life forms that continues to this day.[77]
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Earth's expected long-term future is tied to that of the Sun. Over the next 1.1 billion years, solar luminosity will increase by 10%, and over the next 3.5 billion years by 40%.[78] Earth's increasing surface temperature will accelerate the inorganic carbon cycle, reducing CO2 concentration to levels lethally low for plants (10 ppm for C4 photosynthesis) in approximately 100–900 million years.[79][80] The lack of vegetation will result in the loss of oxygen in the atmosphere, making animal life impossible.[81] About a billion years from now, all surface water will have disappeared[82] and the mean global temperature will reach 70 °C (158 °F).[81] Earth is expected to be habitable until the end of photosynthesis about 500 million years from now,[79] but if nitrogen is removed from the atmosphere, life may continue until a runaway greenhouse effect occurs 2.3 billion years from now.[80] Anthropogenic emissions are "probably insufficient" to cause a runaway greenhouse at current solar luminosity.[83] Even if the Sun were eternal and stable, 27% of the water in the modern oceans will descend to the mantle in one billion years, due to reduced steam venting from mid-ocean ridges.[84]
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The Sun will evolve to become a red giant in about 5 billion years. Models predict that the Sun will expand to roughly 1 AU (150 million km; 93 million mi), about 250 times its present radius.[78][85] Earth's fate is less clear. As a red giant, the Sun will lose roughly 30% of its mass, so, without tidal effects, Earth will move to an orbit 1.7 AU (250 million km; 160 million mi) from the Sun when the star reaches its maximum radius. Most, if not all, remaining life will be destroyed by the Sun's increased luminosity (peaking at about 5,000 times its present level).[78] A 2008 simulation indicates that Earth's orbit will eventually decay due to tidal effects and drag, causing it to enter the Sun's atmosphere and be vaporized.[85]
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The shape of Earth is nearly spherical. There is a small flattening at the poles and bulging around the equator due to Earth's rotation.[89] To second order, Earth is approximately an oblate spheroid, whose equatorial diameter is 43 kilometres (27 mi) larger than the pole-to-pole diameter,[90] although the variation is less than 1% of the average radius of the Earth.
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The point on the surface farthest from Earth's center of mass is the summit of the equatorial Chimborazo volcano in Ecuador (6,384.4 km or 3,967.1 mi).[91][92][93][94] The average diameter of the reference spheroid is 12,742 kilometres (7,918 mi). Local topography deviates from this idealized spheroid, although on a global scale these deviations are small compared to Earth's radius: the maximum deviation of only 0.17% is at the Mariana Trench (10,911 metres or 35,797 feet below local sea level), whereas Mount Everest (8,848 metres or 29,029 feet above local sea level) represents a deviation of 0.14%.[n 14]
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In geodesy, the exact shape that Earth's oceans would adopt in the absence of land and perturbations such as tides and winds is called the geoid. More precisely, the geoid is the surface of gravitational equipotential at mean sea level.
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Earth's mass is approximately 5.97×1024 kg (5,970 Yg). It is composed mostly of iron (32.1%), oxygen (30.1%), silicon (15.1%), magnesium (13.9%), sulphur (2.9%), nickel (1.8%), calcium (1.5%), and aluminum (1.4%), with the remaining 1.2% consisting of trace amounts of other elements. Due to mass segregation, the core region is estimated to be primarily composed of iron (88.8%), with smaller amounts of nickel (5.8%), sulphur (4.5%), and less than 1% trace elements.[98]
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The most common rock constituents of the crust are nearly all oxides: chlorine, sulphur, and fluorine are the important exceptions to this and their total amount in any rock is usually much less than 1%. Over 99% of the crust is composed of 11 oxides, principally silica, alumina, iron oxides, lime, magnesia, potash and soda.[99][98][100]
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Earth's interior, like that of the other terrestrial planets, is divided into layers by their chemical or physical (rheological) properties. The outer layer is a chemically distinct silicate solid crust, which is underlain by a highly viscous solid mantle. The crust is separated from the mantle by the Mohorovičić discontinuity. The thickness of the crust varies from about 6 kilometres (3.7 mi) under the oceans to 30–50 km (19–31 mi) for the continents. The crust and the cold, rigid, top of the upper mantle are collectively known as the lithosphere, and it is of the lithosphere that the tectonic plates are composed. Beneath the lithosphere is the asthenosphere, a relatively low-viscosity layer on which the lithosphere rides. Important changes in crystal structure within the mantle occur at 410 and 660 km (250 and 410 mi) below the surface, spanning a transition zone that separates the upper and lower mantle. Beneath the mantle, an extremely low viscosity liquid outer core lies above a solid inner core.[101] Earth's inner core might rotate at a slightly higher angular velocity than the remainder of the planet, advancing by 0.1–0.5° per year.[102] The radius of the inner core is about one fifth of that of Earth.
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Earth's internal heat comes from a combination of residual heat from planetary accretion (about 20%) and heat produced through radioactive decay (80%).[105] The major heat-producing isotopes within Earth are potassium-40, uranium-238, and thorium-232.[106] At the center, the temperature may be up to 6,000 °C (10,830 °F),[107] and the pressure could reach 360 GPa (52 million psi).[108] Because much of the heat is provided by radioactive decay, scientists postulate that early in Earth's history, before isotopes with short half-lives were depleted, Earth's heat production was much higher. At approximately 3 Gyr, twice the present-day heat would have been produced, increasing the rates of mantle convection and plate tectonics, and allowing the production of uncommon igneous rocks such as komatiites that are rarely formed today.[105][109]
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The mean heat loss from Earth is 87 mW m−2, for a global heat loss of 4.42×1013 W.[111] A portion of the core's thermal energy is transported toward the crust by mantle plumes, a form of convection consisting of upwellings of higher-temperature rock. These plumes can produce hotspots and flood basalts.[112] More of the heat in Earth is lost through plate tectonics, by mantle upwelling associated with mid-ocean ridges. The final major mode of heat loss is through conduction through the lithosphere, the majority of which occurs under the oceans because the crust there is much thinner than that of the continents.[113]
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Earth's mechanically rigid outer layer, the lithosphere, is divided into tectonic plates. These plates are rigid segments that move relative to each other at one of three boundaries types: At convergent boundaries, two plates come together; at divergent boundaries, two plates are pulled apart; and at transform boundaries, two plates slide past one another laterally. Along these plate boundaries, earthquakes, volcanic activity, mountain-building, and oceanic trench formation can occur.[115] The tectonic plates ride on top of the asthenosphere, the solid but less-viscous part of the upper mantle that can flow and move along with the plates.[116]
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As the tectonic plates migrate, oceanic crust is subducted under the leading edges of the plates at convergent boundaries. At the same time, the upwelling of mantle material at divergent boundaries creates mid-ocean ridges. The combination of these processes recycles the oceanic crust back into the mantle. Due to this recycling, most of the ocean floor is less than 100 Ma old. The oldest oceanic crust is located in the Western Pacific and is estimated to be 200 Ma old.[117][118] By comparison, the oldest dated continental crust is 4,030 Ma.[119]
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The seven major plates are the Pacific, North American, Eurasian, African, Antarctic, Indo-Australian, and South American. Other notable plates include the Arabian Plate, the Caribbean Plate, the Nazca Plate off the west coast of South America and the Scotia Plate in the southern Atlantic Ocean. The Australian Plate fused with the Indian Plate between 50 and 55 Mya. The fastest-moving plates are the oceanic plates, with the Cocos Plate advancing at a rate of 75 mm/a (3.0 in/year)[120] and the Pacific Plate moving 52–69 mm/a (2.0–2.7 in/year). At the other extreme, the slowest-moving plate is the Eurasian Plate, progressing at a typical rate of 21 mm/a (0.83 in/year).[121]
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The total surface area of Earth is about 510 million km2 (197 million sq mi).[12] Of this, 70.8%,[12] or 361.13 million km2 (139.43 million sq mi), is below sea level and covered by ocean water.[122] Below the ocean's surface are much of the continental shelf, mountains, volcanoes,[90] oceanic trenches, submarine canyons, oceanic plateaus, abyssal plains, and a globe-spanning mid-ocean ridge system. The remaining 29.2%, or 148.94 million km2 (57.51 million sq mi), not covered by water has terrain that varies greatly from place to place and consists of mountains, deserts, plains, plateaus, and other landforms. Tectonics and erosion, volcanic eruptions, flooding, weathering, glaciation, the growth of coral reefs, and meteorite impacts are among the processes that constantly reshape Earth's surface over geological time.[123][124]
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The continental crust consists of lower density material such as the igneous rocks granite and andesite. Less common is basalt, a denser volcanic rock that is the primary constituent of the ocean floors.[125] Sedimentary rock is formed from the accumulation of sediment that becomes buried and compacted together. Nearly 75% of the continental surfaces are covered by sedimentary rocks, although they form about 5% of the crust.[126] The third form of rock material found on Earth is metamorphic rock, which is created from the transformation of pre-existing rock types through high pressures, high temperatures, or both. The most abundant silicate minerals on Earth's surface include quartz, feldspars, amphibole, mica, pyroxene and olivine.[127] Common carbonate minerals include calcite (found in limestone) and dolomite.[128]
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The elevation of the land surface varies from the low point of −418 m (−1,371 ft) at the Dead Sea, to a maximum altitude of 8,848 m (29,029 ft) at the top of Mount Everest. The mean height of land above sea level is about 797 m (2,615 ft).[129]
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The pedosphere is the outermost layer of Earth's continental surface and is composed of soil and subject to soil formation processes. The total arable land is 10.9% of the land surface, with 1.3% being permanent cropland.[130][131] Close to 40% of Earth's land surface is used for agriculture, or an estimated 16.7 million km2 (6.4 million sq mi) of cropland and 33.5 million km2 (12.9 million sq mi) of pastureland.[132]
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The abundance of water on Earth's surface is a unique feature that distinguishes the "Blue Planet" from other planets in the Solar System. Earth's hydrosphere consists chiefly of the oceans, but technically includes all water surfaces in the world, including inland seas, lakes, rivers, and underground waters down to a depth of 2,000 m (6,600 ft). The deepest underwater location is Challenger Deep of the Mariana Trench in the Pacific Ocean with a depth of 10,911.4 m (35,799 ft).[n 18][133]
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The mass of the oceans is approximately 1.35×1018 metric tons or about 1/4400 of Earth's total mass. The oceans cover an area of 361.8 million km2 (139.7 million sq mi) with a mean depth of 3,682 m (12,080 ft), resulting in an estimated volume of 1.332 billion km3 (320 million cu mi).[134] If all of Earth's crustal surface were at the same elevation as a smooth sphere, the depth of the resulting world ocean would be 2.7 to 2.8 km (1.68 to 1.74 mi).[135][136]
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About 97.5% of the water is saline; the remaining 2.5% is fresh water. Most fresh water, about 68.7%, is present as ice in ice caps and glaciers.[137]
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The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5% salt).[138] Most of this salt was released from volcanic activity or extracted from cool igneous rocks.[139] The oceans are also a reservoir of dissolved atmospheric gases, which are essential for the survival of many aquatic life forms.[140] Sea water has an important influence on the world's climate, with the oceans acting as a large heat reservoir.[141] Shifts in the oceanic temperature distribution can cause significant weather shifts, such as the El Niño–Southern Oscillation.[142]
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The atmospheric pressure at Earth's sea level averages 101.325 kPa (14.696 psi),[143] with a scale height of about 8.5 km (5.3 mi).[3] A dry atmosphere is composed of 78.084% nitrogen, 20.946% oxygen, 0.934% argon, and trace amounts of carbon dioxide and other gaseous molecules.[143] Water vapor content varies between 0.01% and 4%[143] but averages about 1%.[3] The height of the troposphere varies with latitude, ranging between 8 km (5 mi) at the poles to 17 km (11 mi) at the equator, with some variation resulting from weather and seasonal factors.[144]
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Earth's biosphere has significantly altered its atmosphere. Oxygenic photosynthesis evolved 2.7 Gya, forming the primarily nitrogen–oxygen atmosphere of today.[65] This change enabled the proliferation of aerobic organisms and, indirectly, the formation of the ozone layer due to the subsequent conversion of atmospheric O2 into O3. The ozone layer blocks ultraviolet solar radiation, permitting life on land.[145] Other atmospheric functions important to life include transporting water vapor, providing useful gases, causing small meteors to burn up before they strike the surface, and moderating temperature.[146] This last phenomenon is known as the greenhouse effect: trace molecules within the atmosphere serve to capture thermal energy emitted from the ground, thereby raising the average temperature. Water vapor, carbon dioxide, methane, nitrous oxide, and ozone are the primary greenhouse gases in the atmosphere. Without this heat-retention effect, the average surface temperature would be −18 °C (0 °F), in contrast to the current +15 °C (59 °F),[147] and life on Earth probably would not exist in its current form.[148] In May 2017, glints of light, seen as twinkling from an orbiting satellite a million miles away, were found to be reflected light from ice crystals in the atmosphere.[149][150]
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Earth's atmosphere has no definite boundary, slowly becoming thinner and fading into outer space. Three-quarters of the atmosphere's mass is contained within the first 11 km (6.8 mi) of the surface. This lowest layer is called the troposphere. Energy from the Sun heats this layer, and the surface below, causing expansion of the air. This lower-density air then rises and is replaced by cooler, higher-density air. The result is atmospheric circulation that drives the weather and climate through redistribution of thermal energy.[151]
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The primary atmospheric circulation bands consist of the trade winds in the equatorial region below 30° latitude and the westerlies in the mid-latitudes between 30° and 60°.[152] Ocean currents are also important factors in determining climate, particularly the thermohaline circulation that distributes thermal energy from the equatorial oceans to the polar regions.[153]
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Water vapor generated through surface evaporation is transported by circulatory patterns in the atmosphere. When atmospheric conditions permit an uplift of warm, humid air, this water condenses and falls to the surface as precipitation.[151] Most of the water is then transported to lower elevations by river systems and usually returned to the oceans or deposited into lakes. This water cycle is a vital mechanism for supporting life on land and is a primary factor in the erosion of surface features over geological periods. Precipitation patterns vary widely, ranging from several meters of water per year to less than a millimeter. Atmospheric circulation, topographic features, and temperature differences determine the average precipitation that falls in each region.[154]
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The amount of solar energy reaching Earth's surface decreases with increasing latitude. At higher latitudes, the sunlight reaches the surface at lower angles, and it must pass through thicker columns of the atmosphere. As a result, the mean annual air temperature at sea level decreases by about 0.4 °C (0.7 °F) per degree of latitude from the equator.[155] Earth's surface can be subdivided into specific latitudinal belts of approximately homogeneous climate. Ranging from the equator to the polar regions, these are the tropical (or equatorial), subtropical, temperate and polar climates.[156]
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This latitudinal rule has several anomalies:
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The commonly used Köppen climate classification system has five broad groups (humid tropics, arid, humid middle latitudes, continental and cold polar), which are further divided into more specific subtypes.[152] The Köppen system rates regions of terrain based on observed temperature and precipitation.
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The highest air temperature ever measured on Earth was 56.7 °C (134.1 °F) in Furnace Creek, California, in Death Valley, in 1913.[159] The lowest air temperature ever directly measured on Earth was −89.2 °C (−128.6 °F) at Vostok Station in 1983,[160] but satellites have used remote sensing to measure temperatures as low as −94.7 °C (−138.5 °F) in East Antarctica.[161] These temperature records are only measurements made with modern instruments from the 20th century onwards and likely do not reflect the full range of temperature on Earth.
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Above the troposphere, the atmosphere is usually divided into the stratosphere, mesosphere, and thermosphere.[146] Each layer has a different lapse rate, defining the rate of change in temperature with height. Beyond these, the exosphere thins out into the magnetosphere, where the geomagnetic fields interact with the solar wind.[162] Within the stratosphere is the ozone layer, a component that partially shields the surface from ultraviolet light and thus is important for life on Earth. The Kármán line, defined as 100 km above Earth's surface, is a working definition for the boundary between the atmosphere and outer space.[163]
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Thermal energy causes some of the molecules at the outer edge of the atmosphere to increase their velocity to the point where they can escape from Earth's gravity. This causes a slow but steady loss of the atmosphere into space. Because unfixed hydrogen has a low molecular mass, it can achieve escape velocity more readily, and it leaks into outer space at a greater rate than other gases.[164] The leakage of hydrogen into space contributes to the shifting of Earth's atmosphere and surface from an initially reducing state to its current oxidizing one. Photosynthesis provided a source of free oxygen, but the loss of reducing agents such as hydrogen is thought to have been a necessary precondition for the widespread accumulation of oxygen in the atmosphere.[165] Hence the ability of hydrogen to escape from the atmosphere may have influenced the nature of life that developed on Earth.[166] In the current, oxygen-rich atmosphere most hydrogen is converted into water before it has an opportunity to escape. Instead, most of the hydrogen loss comes from the destruction of methane in the upper atmosphere.[167]
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The gravity of Earth is the acceleration that is imparted to objects due to the distribution of mass within Earth. Near Earth's surface, gravitational acceleration is approximately 9.8 m/s2 (32 ft/s2). Local differences in topography, geology, and deeper tectonic structure cause local and broad, regional differences in Earth's gravitational field, known as gravity anomalies.[168]
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The main part of Earth's magnetic field is generated in the core, the site of a dynamo process that converts the kinetic energy of thermally and compositionally driven convection into electrical and magnetic field energy. The field extends outwards from the core, through the mantle, and up to Earth's surface, where it is, approximately, a dipole. The poles of the dipole are located close to Earth's geographic poles. At the equator of the magnetic field, the magnetic-field strength at the surface is 3.05×10−5 T, with a magnetic dipole moment of 7.79×1022 Am2 at epoch 2000, decreasing nearly 6% per century.[169] The convection movements in the core are chaotic; the magnetic poles drift and periodically change alignment. This causes secular variation of the main field and field reversals at irregular intervals averaging a few times every million years. The most recent reversal occurred approximately 700,000 years ago.[170][171]
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The extent of Earth's magnetic field in space defines the magnetosphere. Ions and electrons of the solar wind are deflected by the magnetosphere; solar wind pressure compresses the dayside of the magnetosphere, to about 10 Earth radii, and extends the nightside magnetosphere into a long tail.[172] Because the velocity of the solar wind is greater than the speed at which waves propagate through the solar wind, a supersonic bow shock precedes the dayside magnetosphere within the solar wind.[173] Charged particles are contained within the magnetosphere; the plasmasphere is defined by low-energy particles that essentially follow magnetic field lines as Earth rotates;[174][175] the ring current is defined by medium-energy particles that drift relative to the geomagnetic field, but with paths that are still dominated by the magnetic field,[176] and the Van Allen radiation belt are formed by high-energy particles whose motion is essentially random, but otherwise contained by the magnetosphere.[172][177]
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During magnetic storms and substorms, charged particles can be deflected from the outer magnetosphere and especially the magnetotail, directed along field lines into Earth's ionosphere, where atmospheric atoms can be excited and ionized, causing the aurora.[178]
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Earth's rotation period relative to the Sun—its mean solar day—is 86,400 seconds of mean solar time (86,400.0025 SI seconds).[179] Because Earth's solar day is now slightly longer than it was during the 19th century due to tidal deceleration, each day varies between 0 and 2 SI ms longer.[180][181]
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Earth's rotation period relative to the fixed stars, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is 86,164.0989 seconds of mean solar time (UT1), or 23h 56m 4.0989s.[2][n 19] Earth's rotation period relative to the precessing or moving mean March equinox, misnamed its sidereal day, is 86,164.0905 seconds of mean solar time (UT1) (23h 56m 4.0905s).[2] Thus the sidereal day is shorter than the stellar day by about 8.4 ms.[182] The length of the mean solar day in SI seconds is available from the IERS for the periods 1623–2005[183] and 1962–2005.[184]
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Apart from meteors within the atmosphere and low-orbiting satellites, the main apparent motion of celestial bodies in Earth's sky is to the west at a rate of 15°/h = 15'/min. For bodies near the celestial equator, this is equivalent to an apparent diameter of the Sun or the Moon every two minutes; from Earth's surface, the apparent sizes of the Sun and the Moon are approximately the same.[185][186]
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Earth orbits the Sun at an average distance of about 150 million km (93 million mi) every 365.2564 mean solar days, or one sidereal year. This gives an apparent movement of the Sun eastward with respect to the stars at a rate of about 1°/day, which is one apparent Sun or Moon diameter every 12 hours. Due to this motion, on average it takes 24 hours—a solar day—for Earth to complete a full rotation about its axis so that the Sun returns to the meridian. The orbital speed of Earth averages about 29.78 km/s (107,200 km/h; 66,600 mph), which is fast enough to travel a distance equal to Earth's diameter, about 12,742 km (7,918 mi), in seven minutes, and the distance to the Moon, 384,000 km (239,000 mi), in about 3.5 hours.[3]
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The Moon and Earth orbit a common barycenter every 27.32 days relative to the background stars. When combined with the Earth–Moon system's common orbit around the Sun, the period of the synodic month, from new moon to new moon, is 29.53 days. Viewed from the celestial north pole, the motion of Earth, the Moon, and their axial rotations are all counterclockwise. Viewed from a vantage point above the north poles of both the Sun and Earth, Earth orbits in a counterclockwise direction about the Sun. The orbital and axial planes are not precisely aligned: Earth's axis is tilted some 23.44 degrees from the perpendicular to the Earth–Sun plane (the ecliptic), and the Earth–Moon plane is tilted up to ±5.1 degrees against the Earth–Sun plane. Without this tilt, there would be an eclipse every two weeks, alternating between lunar eclipses and solar eclipses.[3][188]
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The Hill sphere, or the sphere of gravitational influence, of Earth is about 1.5 million km (930,000 mi) in radius.[189][n 20] This is the maximum distance at which Earth's gravitational influence is stronger than the more distant Sun and planets. Objects must orbit Earth within this radius, or they can become unbound by the gravitational perturbation of the Sun.
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Earth, along with the Solar System, is situated in the Milky Way and orbits about 28,000 light-years from its center. It is about 20 light-years above the galactic plane in the Orion Arm.[190]
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The axial tilt of Earth is approximately 23.439281°[2] with the axis of its orbit plane, always pointing towards the Celestial Poles. Due to Earth's axial tilt, the amount of sunlight reaching any given point on the surface varies over the course of the year. This causes the seasonal change in climate, with summer in the Northern Hemisphere occurring when the Tropic of Cancer is facing the Sun, and winter taking place when the Tropic of Capricorn in the Southern Hemisphere faces the Sun. During the summer, the day lasts longer, and the Sun climbs higher in the sky. In winter, the climate becomes cooler and the days shorter. In northern temperate latitudes, the Sun rises north of true east during the summer solstice, and sets north of true west, reversing in the winter. The Sun rises south of true east in the summer for the southern temperate zone and sets south of true west.
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Above the Arctic Circle, an extreme case is reached where there is no daylight at all for part of the year, up to six months at the North Pole itself, a polar night. In the Southern Hemisphere, the situation is exactly reversed, with the South Pole oriented opposite the direction of the North Pole. Six months later, this pole will experience a midnight sun, a day of 24 hours, again reversing with the South Pole.
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By astronomical convention, the four seasons can be determined by the solstices—the points in the orbit of maximum axial tilt toward or away from the Sun—and the equinoxes, when Earth's rotational axis is aligned with its orbital axis. In the Northern Hemisphere, winter solstice currently occurs around 21 December; summer solstice is near 21 June, spring equinox is around 20 March and autumnal equinox is about 22 or 23 September. In the Southern Hemisphere, the situation is reversed, with the summer and winter solstices exchanged and the spring and autumnal equinox dates swapped.[191]
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The angle of Earth's axial tilt is relatively stable over long periods of time. Its axial tilt does undergo nutation; a slight, irregular motion with a main period of 18.6 years.[192] The orientation (rather than the angle) of Earth's axis also changes over time, precessing around in a complete circle over each 25,800 year cycle; this precession is the reason for the difference between a sidereal year and a tropical year. Both of these motions are caused by the varying attraction of the Sun and the Moon on Earth's equatorial bulge. The poles also migrate a few meters across Earth's surface. This polar motion has multiple, cyclical components, which collectively are termed quasiperiodic motion. In addition to an annual component to this motion, there is a 14-month cycle called the Chandler wobble. Earth's rotational velocity also varies in a phenomenon known as length-of-day variation.[193]
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In modern times, Earth's perihelion occurs around 3 January, and its aphelion around 4 July. These dates change over time due to precession and other orbital factors, which follow cyclical patterns known as Milankovitch cycles. The changing Earth–Sun distance causes an increase of about 6.9%[n 21] in solar energy reaching Earth at perihelion relative to aphelion. Because the Southern Hemisphere is tilted toward the Sun at about the same time that Earth reaches the closest approach to the Sun, the Southern Hemisphere receives slightly more energy from the Sun than does the northern over the course of a year. This effect is much less significant than the total energy change due to the axial tilt, and most of the excess energy is absorbed by the higher proportion of water in the Southern Hemisphere.[194]
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A study from 2016 suggested that Planet Nine tilted all the planets of the Solar System, including Earth, by about six degrees.[195]
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A planet that can sustain life is termed habitable, even if life did not originate there. Earth provides liquid water—an environment where complex organic molecules can assemble and interact, and sufficient energy to sustain metabolism.[196] The distance of Earth from the Sun, as well as its orbital eccentricity, rate of rotation, axial tilt, geological history, sustaining atmosphere, and magnetic field all contribute to the current climatic conditions at the surface.[197]
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A planet's life forms inhabit ecosystems, whose total is sometimes said to form a "biosphere".[198] Earth's biosphere is thought to have begun evolving about 3.5 Gya.[65] The biosphere is divided into a number of biomes, inhabited by broadly similar plants and animals.[199] On land, biomes are separated primarily by differences in latitude, height above sea level and humidity. Terrestrial biomes lying within the Arctic or Antarctic Circles, at high altitudes or in extremely arid areas are relatively barren of plant and animal life; species diversity reaches a peak in humid lowlands at equatorial latitudes.[200]
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In July 2016, scientists reported identifying a set of 355 genes from the last universal common ancestor (LUCA) of all organisms living on Earth.[201]
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Earth has resources that have been exploited by humans.[203] Those termed non-renewable resources, such as fossil fuels, only renew over geological timescales.[204]
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Large deposits of fossil fuels are obtained from Earth's crust, consisting of coal, petroleum, and natural gas.[205] These deposits are used by humans both for energy production and as feedstock for chemical production.[206] Mineral ore bodies have also been formed within the crust through a process of ore genesis, resulting from actions of magmatism, erosion, and plate tectonics.[207] These bodies form concentrated sources for many metals and other useful elements.
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Earth's biosphere produces many useful biological products for humans, including food, wood, pharmaceuticals, oxygen, and the recycling of many organic wastes. The land-based ecosystem depends upon topsoil and fresh water, and the oceanic ecosystem depends upon dissolved nutrients washed down from the land.[208] In 1980, 50.53 million km2 (19.51 million sq mi) of Earth's land surface consisted of forest and woodlands, 67.88 million km2 (26.21 million sq mi) was grasslands and pasture, and 15.01 million km2 (5.80 million sq mi) was cultivated as croplands.[209] The estimated amount of irrigated land in 1993 was 2,481,250 km2 (958,020 sq mi).[13] Humans also live on the land by using building materials to construct shelters.
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Large areas of Earth's surface are subject to extreme weather such as tropical cyclones, hurricanes, or typhoons that dominate life in those areas. From 1980 to 2000, these events caused an average of 11,800 human deaths per year.[210] Many places are subject to earthquakes, landslides, tsunamis, volcanic eruptions, tornadoes, sinkholes, blizzards, floods, droughts, wildfires, and other calamities and disasters.
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Many localized areas are subject to human-made pollution of the air and water, acid rain and toxic substances, loss of vegetation (overgrazing, deforestation, desertification), loss of wildlife, species extinction, soil degradation, soil depletion and erosion.
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There is a scientific consensus linking human activities to global warming due to industrial carbon dioxide emissions. This is predicted to produce changes such as the melting of glaciers and ice sheets, more extreme temperature ranges, significant changes in weather and a global rise in average sea levels.[211]
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Cartography, the study and practice of map-making, and geography, the study of the lands, features, inhabitants and phenomena on Earth, have historically been the disciplines devoted to depicting Earth. Surveying, the determination of locations and distances, and to a lesser extent navigation, the determination of position and direction, have developed alongside cartography and geography, providing and suitably quantifying the requisite information.
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Earth's human population reached approximately seven billion on 31 October 2011.[213] Projections indicate that the world's human population will reach 9.2 billion in 2050.[214] Most of the growth is expected to take place in developing nations. Human population density varies widely around the world, but a majority live in Asia. By 2020, 60% of the world's population is expected to be living in urban, rather than rural, areas.[215]
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68% of the land mass of the world is in the northern hemisphere.[216] Partly due to the predominance of land mass, 90% of humans live in the northern hemisphere.[217]
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It is estimated that one-eighth of Earth's surface is suitable for humans to live on – three-quarters of Earth's surface is covered by oceans, leaving one-quarter as land. Half of that land area is desert (14%),[218] high mountains (27%),[219] or other unsuitable terrains. The northernmost permanent settlement in the world is Alert, on Ellesmere Island in Nunavut, Canada.[220] (82°28′N) The southernmost is the Amundsen–Scott South Pole Station, in Antarctica, almost exactly at the South Pole. (90°S)
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States claim the planet's entire land surface, except for parts of Antarctica and a few other unclaimed areas. Earth has never had a planetwide government, but the United Nations is the leading worldwide intergovernmental organization.
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The first human to orbit Earth was Yuri Gagarin on 12 April 1961.[221] In total, about 487 people have visited outer space and reached orbit as of 30 July 2010[update], and, of these, twelve have walked on the Moon.[222][223][224] Normally, the only humans in space are those on the International Space Station. The station's crew, made up of six people, is usually replaced every six months.[225] The farthest that humans have traveled from Earth is 400,171 km (248,655 mi), achieved during the Apollo 13 mission in 1970.[226]
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The Moon is a relatively large, terrestrial, planet-like natural satellite, with a diameter about one-quarter of Earth's. It is the largest moon in the Solar System relative to the size of its planet, although Charon is larger relative to the dwarf planet Pluto. The natural satellites of other planets are also referred to as "moons", after Earth's.
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The gravitational attraction between Earth and the Moon causes tides on Earth. The same effect on the Moon has led to its tidal locking: its rotation period is the same as the time it takes to orbit Earth. As a result, it always presents the same face to the planet. As the Moon orbits Earth, different parts of its face are illuminated by the Sun, leading to the lunar phases; the dark part of the face is separated from the light part by the solar terminator.
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Due to their tidal interaction, the Moon recedes from Earth at the rate of approximately 38 mm/a (1.5 in/year). Over millions of years, these tiny modifications—and the lengthening of Earth's day by about 23 µs/yr—add up to significant changes.[227] During the Devonian period, for example, (approximately 410 Mya) there were 400 days in a year, with each day lasting 21.8 hours.[228]
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The Moon may have dramatically affected the development of life by moderating the planet's climate. Paleontological evidence and computer simulations show that Earth's axial tilt is stabilized by tidal interactions with the Moon.[229] Some theorists think that without this stabilization against the torques applied by the Sun and planets to Earth's equatorial bulge, the rotational axis might be chaotically unstable, exhibiting chaotic changes over millions of years, as appears to be the case for Mars.[230]
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Viewed from Earth, the Moon is just far enough away to have almost the same apparent-sized disk as the Sun. The angular size (or solid angle) of these two bodies match because, although the Sun's diameter is about 400 times as large as the Moon's, it is also 400 times more distant.[186] This allows total and annular solar eclipses to occur on Earth.
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The most widely accepted theory of the Moon's origin, the giant-impact hypothesis, states that it formed from the collision of a Mars-size protoplanet called Theia with the early Earth. This hypothesis explains (among other things) the Moon's relative lack of iron and volatile elements and the fact that its composition is nearly identical to that of Earth's crust.[48]
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Earth has at least five co-orbital asteroids, including 3753 Cruithne and 2002 AA29.[231][232] A trojan asteroid companion, 2010 TK7, is librating around the leading Lagrange triangular point, L4, in Earth's orbit around the Sun.[233][234]
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The tiny near-Earth asteroid 2006 RH120 makes close approaches to the Earth–Moon system roughly every twenty years. During these approaches, it can orbit Earth for brief periods of time.[235]
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As of April 2018[update], there are 1,886 operational, human-made satellites orbiting Earth.[5] There are also inoperative satellites, including Vanguard 1, the oldest satellite currently in orbit, and over 16,000 pieces of tracked space debris.[n 3] Earth's largest artificial satellite is the International Space Station.
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The standard astronomical symbol of Earth consists of a cross circumscribed by a circle, ,[236] representing the four corners of the world.
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Human cultures have developed many views of the planet.[237] Earth is sometimes personified as a deity. In many cultures it is a mother goddess that is also the primary fertility deity,[238] and by the mid-20th century, the Gaia Principle compared Earth's environments and life as a single self-regulating organism leading to broad stabilization of the conditions of habitability.[239][240][241] Creation myths in many religions involve the creation of Earth by a supernatural deity or deities.[238]
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The Hindu Vedas (1500–900 BC) refer to the Earth as Bhūgola (भूगोल), which comes from Bhū (earth, ground) and Gola (ball, sphere, globe). It means the "globe of earth". There is no direct evidence that the Hindus of that time knew that the Earth was sphere-shaped, but this name has been used extensively since the inception of the Vedas.[citation needed]
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Scientific investigation has resulted in several culturally transformative shifts in people's view of the planet. Initial belief in a flat Earth was gradually displaced in the Greek colonies of southern Italy during the late 6th century BC by the idea of spherical Earth,[242][243][244] which was attributed to both the philosophers Pythagoras and Parmenides.[243][244] By the end of the 5th century BC, the sphericity of Earth was universally accepted among Greek intellectuals.[245] Earth was generally believed to be the center of the universe until the 16th century, when scientists first conclusively demonstrated that it was a moving object, comparable to the other planets in the Solar System.[246] Due to the efforts of influential Christian scholars and clerics such as James Ussher, who sought to determine the age of Earth through analysis of genealogies in Scripture, Westerners before the 19th century generally believed Earth to be a few thousand years old at most. It was only during the 19th century that geologists realized Earth's age was at least many millions of years.[247]
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Lord Kelvin used thermodynamics to estimate the age of Earth to be between 20 million and 400 million years in 1864, sparking a vigorous debate on the subject; it was only when radioactivity and radioactive dating were discovered in the late 19th and early 20th centuries that a reliable mechanism for determining Earth's age was established, proving the planet to be billions of years old.[248][249] The perception of Earth shifted again[further explanation needed] in the 20th century when humans first viewed it from orbit, and especially with photographs of Earth returned by the Apollo program.[250][251][252]
|
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Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
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Earth is the third planet from the Sun and the only astronomical object known to harbor life. According to radiometric dating estimation and other evidence, Earth formed over 4.5 billion years ago. Earth's gravity interacts with other objects in space, especially the Sun and the Moon, which is Earth's only natural satellite. Earth orbits around the Sun in 365.256 solar days, a period known as an Earth sidereal year. During this time, Earth rotates about its axis 366.256 times, that is, a sidereal year has 366.256 sidereal days.[n 6]
|
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|
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Earth's axis of rotation is tilted with respect to its orbital plane, producing seasons on Earth. The gravitational interaction between Earth and the Moon causes tides, stabilizes Earth's orientation on its axis, and gradually slows its rotation. Earth is the densest planet in the Solar System and the largest and most massive of the four rocky planets.
|
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|
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Earth's outer layer (lithosphere) is divided into several rigid tectonic plates that migrate across the surface over many millions of years. About 29% of Earth's surface is land consisting of continents and islands. The remaining 71% is covered with water, mostly by oceans but also lakes, rivers and other fresh water, which all together constitute the hydrosphere. The majority of Earth's polar regions are covered in ice, including the Antarctic ice sheet and the sea ice of the Arctic ice pack. Earth's interior remains active with a solid iron inner core, a liquid outer core that generates Earth's magnetic field, and a convecting mantle that drives plate tectonics.
|
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Within the first billion years of Earth's history, life appeared in the oceans and began to affect Earth's atmosphere and surface, leading to the proliferation of anaerobic and, later, aerobic organisms. Some geological evidence indicates that life may have arisen as early as 4.1 billion years ago. Since then, the combination of Earth's distance from the Sun, physical properties and geological history have allowed life to evolve and thrive. In the history of life on Earth, biodiversity has gone through long periods of expansion, occasionally punctuated by mass extinctions. Over 99% of all species that ever lived on Earth are extinct. Estimates of the number of species on Earth today vary widely; most species have not been described. Over 7.7 billion humans live on Earth and depend on its biosphere and natural resources for their survival.[23]
|
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The modern English word Earth developed, via Middle English,[n 7] from an Old English noun most often spelled eorðe.[24] It has cognates in every Germanic language, and their ancestral root has been reconstructed as *erþō. In its earliest attestation, the word eorðe was already being used to translate the many senses of Latin terra and Greek γῆ gē: the ground,[n 8] its soil,[n 9] dry land,[n 10] the human world,[n 11] the surface of the world (including the sea),[n 12] and the globe itself.[n 13] As with Roman Terra/Tellūs and Greek Gaia, Earth may have been a personified goddess in Germanic paganism: late Norse mythology included Jörð ('Earth'), a giantess often given as the mother of Thor.[33]
|
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|
15 |
+
Originally, earth was written in lowercase, and from early Middle English, its definite sense as "the globe" was expressed as the earth. By Early Modern English, many nouns were capitalized, and the earth became (and often remained) the Earth, particularly when referenced along with other heavenly bodies. More recently, the name is sometimes simply given as Earth, by analogy with the names of the other planets.[24] House styles now vary: Oxford spelling recognizes the lowercase form as the most common, with the capitalized form an acceptable variant. Another convention capitalizes "Earth" when appearing as a name (e.g. "Earth's atmosphere") but writes it in lowercase when preceded by the (e.g. "the atmosphere of the earth"). It almost always appears in lowercase in colloquial expressions such as "what on earth are you doing?"[34]
|
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|
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Occasionally, the name Terra /ˈtɛrə/ is used in scientific writing and especially in science fiction to distinguish our inhabited planet from others,[35] while in poetry Tellus /ˈtɛləs/ has been used to denote personification of the Earth.[36] The Greek poetic name Gaea (Gæa) /ˈdʒiːə/ is rare, though the alternative spelling Gaia has become common due to the Gaia hypothesis, in which case its pronunciation is /ˈɡaɪə/ rather than the more Classical /ˈɡeɪə/.[37]
|
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|
19 |
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There are a number of adjectives for the planet Earth. From Earth itself comes earthly. From Latin Terra come Terran /ˈtɛrən/,[38] Terrestrial /təˈrɛstriəl/,[39] and (via French) Terrene /təˈriːn/,[40] and from Latin Tellus come Tellurian /tɛˈlʊəriən/[41] and, more rarely, Telluric and Tellural. From Greek Gaia and Gaea comes Gaian and Gaean.
|
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|
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An inhabitant of the Earth is an Earthling, a Terran, a Terrestrial, a Tellurian or, rarely, an Earthian.
|
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+
|
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+
The oldest material found in the Solar System is dated to 4.5672±0.0006 billion years ago (BYA).[42] By 4.54±0.04 BYA[43] the primordial Earth had formed. The bodies in the Solar System formed and evolved with the Sun. In theory, a solar nebula partitions a volume out of a molecular cloud by gravitational collapse, which begins to spin and flatten into a circumstellar disk, and then the planets grow out of that disk with the Sun. A nebula contains gas, ice grains, and dust (including primordial nuclides). According to nebular theory, planetesimals formed by accretion, with the primordial Earth taking 10–20 million years (Mys) to form.[44]
|
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+
|
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+
A subject of research is the formation of the Moon, some 4.53 BYA.[45] A leading hypothesis is that it was formed by accretion from material loosed from Earth after a Mars-sized object, named Theia, hit Earth.[46] In this view, the mass of Theia was approximately 10 percent of Earth;[47] it hit Earth with a glancing blow and some of its mass merged with Earth.[48] Between approximately 4.1 and 3.8 BYA, numerous asteroid impacts during the Late Heavy Bombardment caused significant changes to the greater surface environment of the Moon and, by inference, to that of Earth.
|
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|
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Earth's atmosphere and oceans were formed by volcanic activity and outgassing.[49] Water vapor from these sources condensed into the oceans, augmented by water and ice from asteroids, protoplanets, and comets.[50] In this model, atmospheric "greenhouse gases" kept the oceans from freezing when the newly forming Sun had only 70% of its current luminosity.[51] By 3.5 BYA, Earth's magnetic field was established, which helped prevent the atmosphere from being stripped away by the solar wind.[52]
|
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+
A crust formed when the molten outer layer of Earth cooled to form a solid. The two models[53] that explain land mass propose either a steady growth to the present-day forms[54] or, more likely, a rapid growth[55] early in Earth history[56] followed by a long-term steady continental area.[57][58][59] Continents formed by plate tectonics, a process ultimately driven by the continuous loss of heat from Earth's interior. Over the period of hundreds of millions of years, the supercontinents have assembled and broken apart. Roughly 750 million years ago (MYA), one of the earliest known supercontinents, Rodinia, began to break apart. The continents later recombined to form Pannotia 600–540 MYA, then finally Pangaea, which also broke apart 180 MYA.[60]
|
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The present pattern of ice ages began about 40 MYA,[61] and then intensified during the Pleistocene about 3 MYA.[62] High-latitude regions have since undergone repeated cycles of glaciation and thaw, repeating about every 40,000–100,000 years. The last continental glaciation ended 10,000 years ago.[63]
|
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|
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Chemical reactions led to the first self-replicating molecules about four billion years ago. A half billion years later, the last common ancestor of all current life arose.[64] The evolution of photosynthesis allowed the Sun's energy to be harvested directly by life forms. The resultant molecular oxygen (O2) accumulated in the atmosphere and due to interaction with ultraviolet solar radiation, formed a protective ozone layer (O3) in the upper atmosphere.[65] The incorporation of smaller cells within larger ones resulted in the development of complex cells called eukaryotes.[66] True multicellular organisms formed as cells within colonies became increasingly specialized. Aided by the absorption of harmful ultraviolet radiation by the ozone layer, life colonized Earth's surface.[67] Among the earliest fossil evidence for life is microbial mat fossils found in 3.48 billion-year-old sandstone in Western Australia,[68] biogenic graphite found in 3.7 billion-year-old metasedimentary rocks in Western Greenland,[69] and remains of biotic material found in 4.1 billion-year-old rocks in Western Australia.[70][71] The earliest direct evidence of life on Earth is contained in 3.45 billion-year-old Australian rocks showing fossils of microorganisms.[72][73]
|
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During the Neoproterozoic, 750 to 580 MYA, much of Earth might have been covered in ice. This hypothesis has been termed "Snowball Earth", and it is of particular interest because it preceded the Cambrian explosion, when multicellular life forms significantly increased in complexity.[74] Following the Cambrian explosion, 535 MYA, there have been five mass extinctions.[75] The most recent such event was 66 MYA, when an asteroid impact triggered the extinction of the non-avian dinosaurs and other large reptiles, but spared some small animals such as mammals, which at the time resembled shrews. Mammalian life has diversified over the past 66 Mys, and several million years ago an African ape-like animal such as Orrorin tugenensis gained the ability to stand upright.[76] This facilitated tool use and encouraged communication that provided the nutrition and stimulation needed for a larger brain, which led to the evolution of humans. The development of agriculture, and then civilization, led to humans having an influence on Earth and the nature and quantity of other life forms that continues to this day.[77]
|
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Earth's expected long-term future is tied to that of the Sun. Over the next 1.1 billion years, solar luminosity will increase by 10%, and over the next 3.5 billion years by 40%.[78] Earth's increasing surface temperature will accelerate the inorganic carbon cycle, reducing CO2 concentration to levels lethally low for plants (10 ppm for C4 photosynthesis) in approximately 100–900 million years.[79][80] The lack of vegetation will result in the loss of oxygen in the atmosphere, making animal life impossible.[81] About a billion years from now, all surface water will have disappeared[82] and the mean global temperature will reach 70 °C (158 °F).[81] Earth is expected to be habitable until the end of photosynthesis about 500 million years from now,[79] but if nitrogen is removed from the atmosphere, life may continue until a runaway greenhouse effect occurs 2.3 billion years from now.[80] Anthropogenic emissions are "probably insufficient" to cause a runaway greenhouse at current solar luminosity.[83] Even if the Sun were eternal and stable, 27% of the water in the modern oceans will descend to the mantle in one billion years, due to reduced steam venting from mid-ocean ridges.[84]
|
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|
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+
The Sun will evolve to become a red giant in about 5 billion years. Models predict that the Sun will expand to roughly 1 AU (150 million km; 93 million mi), about 250 times its present radius.[78][85] Earth's fate is less clear. As a red giant, the Sun will lose roughly 30% of its mass, so, without tidal effects, Earth will move to an orbit 1.7 AU (250 million km; 160 million mi) from the Sun when the star reaches its maximum radius. Most, if not all, remaining life will be destroyed by the Sun's increased luminosity (peaking at about 5,000 times its present level).[78] A 2008 simulation indicates that Earth's orbit will eventually decay due to tidal effects and drag, causing it to enter the Sun's atmosphere and be vaporized.[85]
|
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|
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+
The shape of Earth is nearly spherical. There is a small flattening at the poles and bulging around the equator due to Earth's rotation.[89] To second order, Earth is approximately an oblate spheroid, whose equatorial diameter is 43 kilometres (27 mi) larger than the pole-to-pole diameter,[90] although the variation is less than 1% of the average radius of the Earth.
|
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The point on the surface farthest from Earth's center of mass is the summit of the equatorial Chimborazo volcano in Ecuador (6,384.4 km or 3,967.1 mi).[91][92][93][94] The average diameter of the reference spheroid is 12,742 kilometres (7,918 mi). Local topography deviates from this idealized spheroid, although on a global scale these deviations are small compared to Earth's radius: the maximum deviation of only 0.17% is at the Mariana Trench (10,911 metres or 35,797 feet below local sea level), whereas Mount Everest (8,848 metres or 29,029 feet above local sea level) represents a deviation of 0.14%.[n 14]
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In geodesy, the exact shape that Earth's oceans would adopt in the absence of land and perturbations such as tides and winds is called the geoid. More precisely, the geoid is the surface of gravitational equipotential at mean sea level.
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Earth's mass is approximately 5.97×1024 kg (5,970 Yg). It is composed mostly of iron (32.1%), oxygen (30.1%), silicon (15.1%), magnesium (13.9%), sulphur (2.9%), nickel (1.8%), calcium (1.5%), and aluminum (1.4%), with the remaining 1.2% consisting of trace amounts of other elements. Due to mass segregation, the core region is estimated to be primarily composed of iron (88.8%), with smaller amounts of nickel (5.8%), sulphur (4.5%), and less than 1% trace elements.[98]
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The most common rock constituents of the crust are nearly all oxides: chlorine, sulphur, and fluorine are the important exceptions to this and their total amount in any rock is usually much less than 1%. Over 99% of the crust is composed of 11 oxides, principally silica, alumina, iron oxides, lime, magnesia, potash and soda.[99][98][100]
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Earth's interior, like that of the other terrestrial planets, is divided into layers by their chemical or physical (rheological) properties. The outer layer is a chemically distinct silicate solid crust, which is underlain by a highly viscous solid mantle. The crust is separated from the mantle by the Mohorovičić discontinuity. The thickness of the crust varies from about 6 kilometres (3.7 mi) under the oceans to 30–50 km (19–31 mi) for the continents. The crust and the cold, rigid, top of the upper mantle are collectively known as the lithosphere, and it is of the lithosphere that the tectonic plates are composed. Beneath the lithosphere is the asthenosphere, a relatively low-viscosity layer on which the lithosphere rides. Important changes in crystal structure within the mantle occur at 410 and 660 km (250 and 410 mi) below the surface, spanning a transition zone that separates the upper and lower mantle. Beneath the mantle, an extremely low viscosity liquid outer core lies above a solid inner core.[101] Earth's inner core might rotate at a slightly higher angular velocity than the remainder of the planet, advancing by 0.1–0.5° per year.[102] The radius of the inner core is about one fifth of that of Earth.
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Earth's internal heat comes from a combination of residual heat from planetary accretion (about 20%) and heat produced through radioactive decay (80%).[105] The major heat-producing isotopes within Earth are potassium-40, uranium-238, and thorium-232.[106] At the center, the temperature may be up to 6,000 °C (10,830 °F),[107] and the pressure could reach 360 GPa (52 million psi).[108] Because much of the heat is provided by radioactive decay, scientists postulate that early in Earth's history, before isotopes with short half-lives were depleted, Earth's heat production was much higher. At approximately 3 Gyr, twice the present-day heat would have been produced, increasing the rates of mantle convection and plate tectonics, and allowing the production of uncommon igneous rocks such as komatiites that are rarely formed today.[105][109]
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The mean heat loss from Earth is 87 mW m−2, for a global heat loss of 4.42×1013 W.[111] A portion of the core's thermal energy is transported toward the crust by mantle plumes, a form of convection consisting of upwellings of higher-temperature rock. These plumes can produce hotspots and flood basalts.[112] More of the heat in Earth is lost through plate tectonics, by mantle upwelling associated with mid-ocean ridges. The final major mode of heat loss is through conduction through the lithosphere, the majority of which occurs under the oceans because the crust there is much thinner than that of the continents.[113]
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Earth's mechanically rigid outer layer, the lithosphere, is divided into tectonic plates. These plates are rigid segments that move relative to each other at one of three boundaries types: At convergent boundaries, two plates come together; at divergent boundaries, two plates are pulled apart; and at transform boundaries, two plates slide past one another laterally. Along these plate boundaries, earthquakes, volcanic activity, mountain-building, and oceanic trench formation can occur.[115] The tectonic plates ride on top of the asthenosphere, the solid but less-viscous part of the upper mantle that can flow and move along with the plates.[116]
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As the tectonic plates migrate, oceanic crust is subducted under the leading edges of the plates at convergent boundaries. At the same time, the upwelling of mantle material at divergent boundaries creates mid-ocean ridges. The combination of these processes recycles the oceanic crust back into the mantle. Due to this recycling, most of the ocean floor is less than 100 Ma old. The oldest oceanic crust is located in the Western Pacific and is estimated to be 200 Ma old.[117][118] By comparison, the oldest dated continental crust is 4,030 Ma.[119]
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The seven major plates are the Pacific, North American, Eurasian, African, Antarctic, Indo-Australian, and South American. Other notable plates include the Arabian Plate, the Caribbean Plate, the Nazca Plate off the west coast of South America and the Scotia Plate in the southern Atlantic Ocean. The Australian Plate fused with the Indian Plate between 50 and 55 Mya. The fastest-moving plates are the oceanic plates, with the Cocos Plate advancing at a rate of 75 mm/a (3.0 in/year)[120] and the Pacific Plate moving 52–69 mm/a (2.0–2.7 in/year). At the other extreme, the slowest-moving plate is the Eurasian Plate, progressing at a typical rate of 21 mm/a (0.83 in/year).[121]
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The total surface area of Earth is about 510 million km2 (197 million sq mi).[12] Of this, 70.8%,[12] or 361.13 million km2 (139.43 million sq mi), is below sea level and covered by ocean water.[122] Below the ocean's surface are much of the continental shelf, mountains, volcanoes,[90] oceanic trenches, submarine canyons, oceanic plateaus, abyssal plains, and a globe-spanning mid-ocean ridge system. The remaining 29.2%, or 148.94 million km2 (57.51 million sq mi), not covered by water has terrain that varies greatly from place to place and consists of mountains, deserts, plains, plateaus, and other landforms. Tectonics and erosion, volcanic eruptions, flooding, weathering, glaciation, the growth of coral reefs, and meteorite impacts are among the processes that constantly reshape Earth's surface over geological time.[123][124]
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The continental crust consists of lower density material such as the igneous rocks granite and andesite. Less common is basalt, a denser volcanic rock that is the primary constituent of the ocean floors.[125] Sedimentary rock is formed from the accumulation of sediment that becomes buried and compacted together. Nearly 75% of the continental surfaces are covered by sedimentary rocks, although they form about 5% of the crust.[126] The third form of rock material found on Earth is metamorphic rock, which is created from the transformation of pre-existing rock types through high pressures, high temperatures, or both. The most abundant silicate minerals on Earth's surface include quartz, feldspars, amphibole, mica, pyroxene and olivine.[127] Common carbonate minerals include calcite (found in limestone) and dolomite.[128]
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The elevation of the land surface varies from the low point of −418 m (−1,371 ft) at the Dead Sea, to a maximum altitude of 8,848 m (29,029 ft) at the top of Mount Everest. The mean height of land above sea level is about 797 m (2,615 ft).[129]
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The pedosphere is the outermost layer of Earth's continental surface and is composed of soil and subject to soil formation processes. The total arable land is 10.9% of the land surface, with 1.3% being permanent cropland.[130][131] Close to 40% of Earth's land surface is used for agriculture, or an estimated 16.7 million km2 (6.4 million sq mi) of cropland and 33.5 million km2 (12.9 million sq mi) of pastureland.[132]
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The abundance of water on Earth's surface is a unique feature that distinguishes the "Blue Planet" from other planets in the Solar System. Earth's hydrosphere consists chiefly of the oceans, but technically includes all water surfaces in the world, including inland seas, lakes, rivers, and underground waters down to a depth of 2,000 m (6,600 ft). The deepest underwater location is Challenger Deep of the Mariana Trench in the Pacific Ocean with a depth of 10,911.4 m (35,799 ft).[n 18][133]
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The mass of the oceans is approximately 1.35×1018 metric tons or about 1/4400 of Earth's total mass. The oceans cover an area of 361.8 million km2 (139.7 million sq mi) with a mean depth of 3,682 m (12,080 ft), resulting in an estimated volume of 1.332 billion km3 (320 million cu mi).[134] If all of Earth's crustal surface were at the same elevation as a smooth sphere, the depth of the resulting world ocean would be 2.7 to 2.8 km (1.68 to 1.74 mi).[135][136]
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About 97.5% of the water is saline; the remaining 2.5% is fresh water. Most fresh water, about 68.7%, is present as ice in ice caps and glaciers.[137]
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The average salinity of Earth's oceans is about 35 grams of salt per kilogram of sea water (3.5% salt).[138] Most of this salt was released from volcanic activity or extracted from cool igneous rocks.[139] The oceans are also a reservoir of dissolved atmospheric gases, which are essential for the survival of many aquatic life forms.[140] Sea water has an important influence on the world's climate, with the oceans acting as a large heat reservoir.[141] Shifts in the oceanic temperature distribution can cause significant weather shifts, such as the El Niño–Southern Oscillation.[142]
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The atmospheric pressure at Earth's sea level averages 101.325 kPa (14.696 psi),[143] with a scale height of about 8.5 km (5.3 mi).[3] A dry atmosphere is composed of 78.084% nitrogen, 20.946% oxygen, 0.934% argon, and trace amounts of carbon dioxide and other gaseous molecules.[143] Water vapor content varies between 0.01% and 4%[143] but averages about 1%.[3] The height of the troposphere varies with latitude, ranging between 8 km (5 mi) at the poles to 17 km (11 mi) at the equator, with some variation resulting from weather and seasonal factors.[144]
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Earth's biosphere has significantly altered its atmosphere. Oxygenic photosynthesis evolved 2.7 Gya, forming the primarily nitrogen–oxygen atmosphere of today.[65] This change enabled the proliferation of aerobic organisms and, indirectly, the formation of the ozone layer due to the subsequent conversion of atmospheric O2 into O3. The ozone layer blocks ultraviolet solar radiation, permitting life on land.[145] Other atmospheric functions important to life include transporting water vapor, providing useful gases, causing small meteors to burn up before they strike the surface, and moderating temperature.[146] This last phenomenon is known as the greenhouse effect: trace molecules within the atmosphere serve to capture thermal energy emitted from the ground, thereby raising the average temperature. Water vapor, carbon dioxide, methane, nitrous oxide, and ozone are the primary greenhouse gases in the atmosphere. Without this heat-retention effect, the average surface temperature would be −18 °C (0 °F), in contrast to the current +15 °C (59 °F),[147] and life on Earth probably would not exist in its current form.[148] In May 2017, glints of light, seen as twinkling from an orbiting satellite a million miles away, were found to be reflected light from ice crystals in the atmosphere.[149][150]
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Earth's atmosphere has no definite boundary, slowly becoming thinner and fading into outer space. Three-quarters of the atmosphere's mass is contained within the first 11 km (6.8 mi) of the surface. This lowest layer is called the troposphere. Energy from the Sun heats this layer, and the surface below, causing expansion of the air. This lower-density air then rises and is replaced by cooler, higher-density air. The result is atmospheric circulation that drives the weather and climate through redistribution of thermal energy.[151]
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The primary atmospheric circulation bands consist of the trade winds in the equatorial region below 30° latitude and the westerlies in the mid-latitudes between 30° and 60°.[152] Ocean currents are also important factors in determining climate, particularly the thermohaline circulation that distributes thermal energy from the equatorial oceans to the polar regions.[153]
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Water vapor generated through surface evaporation is transported by circulatory patterns in the atmosphere. When atmospheric conditions permit an uplift of warm, humid air, this water condenses and falls to the surface as precipitation.[151] Most of the water is then transported to lower elevations by river systems and usually returned to the oceans or deposited into lakes. This water cycle is a vital mechanism for supporting life on land and is a primary factor in the erosion of surface features over geological periods. Precipitation patterns vary widely, ranging from several meters of water per year to less than a millimeter. Atmospheric circulation, topographic features, and temperature differences determine the average precipitation that falls in each region.[154]
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The amount of solar energy reaching Earth's surface decreases with increasing latitude. At higher latitudes, the sunlight reaches the surface at lower angles, and it must pass through thicker columns of the atmosphere. As a result, the mean annual air temperature at sea level decreases by about 0.4 °C (0.7 °F) per degree of latitude from the equator.[155] Earth's surface can be subdivided into specific latitudinal belts of approximately homogeneous climate. Ranging from the equator to the polar regions, these are the tropical (or equatorial), subtropical, temperate and polar climates.[156]
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This latitudinal rule has several anomalies:
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The commonly used Köppen climate classification system has five broad groups (humid tropics, arid, humid middle latitudes, continental and cold polar), which are further divided into more specific subtypes.[152] The Köppen system rates regions of terrain based on observed temperature and precipitation.
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The highest air temperature ever measured on Earth was 56.7 °C (134.1 °F) in Furnace Creek, California, in Death Valley, in 1913.[159] The lowest air temperature ever directly measured on Earth was −89.2 °C (−128.6 °F) at Vostok Station in 1983,[160] but satellites have used remote sensing to measure temperatures as low as −94.7 °C (−138.5 °F) in East Antarctica.[161] These temperature records are only measurements made with modern instruments from the 20th century onwards and likely do not reflect the full range of temperature on Earth.
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Above the troposphere, the atmosphere is usually divided into the stratosphere, mesosphere, and thermosphere.[146] Each layer has a different lapse rate, defining the rate of change in temperature with height. Beyond these, the exosphere thins out into the magnetosphere, where the geomagnetic fields interact with the solar wind.[162] Within the stratosphere is the ozone layer, a component that partially shields the surface from ultraviolet light and thus is important for life on Earth. The Kármán line, defined as 100 km above Earth's surface, is a working definition for the boundary between the atmosphere and outer space.[163]
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Thermal energy causes some of the molecules at the outer edge of the atmosphere to increase their velocity to the point where they can escape from Earth's gravity. This causes a slow but steady loss of the atmosphere into space. Because unfixed hydrogen has a low molecular mass, it can achieve escape velocity more readily, and it leaks into outer space at a greater rate than other gases.[164] The leakage of hydrogen into space contributes to the shifting of Earth's atmosphere and surface from an initially reducing state to its current oxidizing one. Photosynthesis provided a source of free oxygen, but the loss of reducing agents such as hydrogen is thought to have been a necessary precondition for the widespread accumulation of oxygen in the atmosphere.[165] Hence the ability of hydrogen to escape from the atmosphere may have influenced the nature of life that developed on Earth.[166] In the current, oxygen-rich atmosphere most hydrogen is converted into water before it has an opportunity to escape. Instead, most of the hydrogen loss comes from the destruction of methane in the upper atmosphere.[167]
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The gravity of Earth is the acceleration that is imparted to objects due to the distribution of mass within Earth. Near Earth's surface, gravitational acceleration is approximately 9.8 m/s2 (32 ft/s2). Local differences in topography, geology, and deeper tectonic structure cause local and broad, regional differences in Earth's gravitational field, known as gravity anomalies.[168]
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The main part of Earth's magnetic field is generated in the core, the site of a dynamo process that converts the kinetic energy of thermally and compositionally driven convection into electrical and magnetic field energy. The field extends outwards from the core, through the mantle, and up to Earth's surface, where it is, approximately, a dipole. The poles of the dipole are located close to Earth's geographic poles. At the equator of the magnetic field, the magnetic-field strength at the surface is 3.05×10−5 T, with a magnetic dipole moment of 7.79×1022 Am2 at epoch 2000, decreasing nearly 6% per century.[169] The convection movements in the core are chaotic; the magnetic poles drift and periodically change alignment. This causes secular variation of the main field and field reversals at irregular intervals averaging a few times every million years. The most recent reversal occurred approximately 700,000 years ago.[170][171]
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The extent of Earth's magnetic field in space defines the magnetosphere. Ions and electrons of the solar wind are deflected by the magnetosphere; solar wind pressure compresses the dayside of the magnetosphere, to about 10 Earth radii, and extends the nightside magnetosphere into a long tail.[172] Because the velocity of the solar wind is greater than the speed at which waves propagate through the solar wind, a supersonic bow shock precedes the dayside magnetosphere within the solar wind.[173] Charged particles are contained within the magnetosphere; the plasmasphere is defined by low-energy particles that essentially follow magnetic field lines as Earth rotates;[174][175] the ring current is defined by medium-energy particles that drift relative to the geomagnetic field, but with paths that are still dominated by the magnetic field,[176] and the Van Allen radiation belt are formed by high-energy particles whose motion is essentially random, but otherwise contained by the magnetosphere.[172][177]
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During magnetic storms and substorms, charged particles can be deflected from the outer magnetosphere and especially the magnetotail, directed along field lines into Earth's ionosphere, where atmospheric atoms can be excited and ionized, causing the aurora.[178]
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Earth's rotation period relative to the Sun—its mean solar day—is 86,400 seconds of mean solar time (86,400.0025 SI seconds).[179] Because Earth's solar day is now slightly longer than it was during the 19th century due to tidal deceleration, each day varies between 0 and 2 SI ms longer.[180][181]
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Earth's rotation period relative to the fixed stars, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is 86,164.0989 seconds of mean solar time (UT1), or 23h 56m 4.0989s.[2][n 19] Earth's rotation period relative to the precessing or moving mean March equinox, misnamed its sidereal day, is 86,164.0905 seconds of mean solar time (UT1) (23h 56m 4.0905s).[2] Thus the sidereal day is shorter than the stellar day by about 8.4 ms.[182] The length of the mean solar day in SI seconds is available from the IERS for the periods 1623–2005[183] and 1962–2005.[184]
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Apart from meteors within the atmosphere and low-orbiting satellites, the main apparent motion of celestial bodies in Earth's sky is to the west at a rate of 15°/h = 15'/min. For bodies near the celestial equator, this is equivalent to an apparent diameter of the Sun or the Moon every two minutes; from Earth's surface, the apparent sizes of the Sun and the Moon are approximately the same.[185][186]
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Earth orbits the Sun at an average distance of about 150 million km (93 million mi) every 365.2564 mean solar days, or one sidereal year. This gives an apparent movement of the Sun eastward with respect to the stars at a rate of about 1°/day, which is one apparent Sun or Moon diameter every 12 hours. Due to this motion, on average it takes 24 hours—a solar day—for Earth to complete a full rotation about its axis so that the Sun returns to the meridian. The orbital speed of Earth averages about 29.78 km/s (107,200 km/h; 66,600 mph), which is fast enough to travel a distance equal to Earth's diameter, about 12,742 km (7,918 mi), in seven minutes, and the distance to the Moon, 384,000 km (239,000 mi), in about 3.5 hours.[3]
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The Moon and Earth orbit a common barycenter every 27.32 days relative to the background stars. When combined with the Earth–Moon system's common orbit around the Sun, the period of the synodic month, from new moon to new moon, is 29.53 days. Viewed from the celestial north pole, the motion of Earth, the Moon, and their axial rotations are all counterclockwise. Viewed from a vantage point above the north poles of both the Sun and Earth, Earth orbits in a counterclockwise direction about the Sun. The orbital and axial planes are not precisely aligned: Earth's axis is tilted some 23.44 degrees from the perpendicular to the Earth–Sun plane (the ecliptic), and the Earth–Moon plane is tilted up to ±5.1 degrees against the Earth–Sun plane. Without this tilt, there would be an eclipse every two weeks, alternating between lunar eclipses and solar eclipses.[3][188]
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The Hill sphere, or the sphere of gravitational influence, of Earth is about 1.5 million km (930,000 mi) in radius.[189][n 20] This is the maximum distance at which Earth's gravitational influence is stronger than the more distant Sun and planets. Objects must orbit Earth within this radius, or they can become unbound by the gravitational perturbation of the Sun.
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Earth, along with the Solar System, is situated in the Milky Way and orbits about 28,000 light-years from its center. It is about 20 light-years above the galactic plane in the Orion Arm.[190]
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The axial tilt of Earth is approximately 23.439281°[2] with the axis of its orbit plane, always pointing towards the Celestial Poles. Due to Earth's axial tilt, the amount of sunlight reaching any given point on the surface varies over the course of the year. This causes the seasonal change in climate, with summer in the Northern Hemisphere occurring when the Tropic of Cancer is facing the Sun, and winter taking place when the Tropic of Capricorn in the Southern Hemisphere faces the Sun. During the summer, the day lasts longer, and the Sun climbs higher in the sky. In winter, the climate becomes cooler and the days shorter. In northern temperate latitudes, the Sun rises north of true east during the summer solstice, and sets north of true west, reversing in the winter. The Sun rises south of true east in the summer for the southern temperate zone and sets south of true west.
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Above the Arctic Circle, an extreme case is reached where there is no daylight at all for part of the year, up to six months at the North Pole itself, a polar night. In the Southern Hemisphere, the situation is exactly reversed, with the South Pole oriented opposite the direction of the North Pole. Six months later, this pole will experience a midnight sun, a day of 24 hours, again reversing with the South Pole.
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By astronomical convention, the four seasons can be determined by the solstices—the points in the orbit of maximum axial tilt toward or away from the Sun—and the equinoxes, when Earth's rotational axis is aligned with its orbital axis. In the Northern Hemisphere, winter solstice currently occurs around 21 December; summer solstice is near 21 June, spring equinox is around 20 March and autumnal equinox is about 22 or 23 September. In the Southern Hemisphere, the situation is reversed, with the summer and winter solstices exchanged and the spring and autumnal equinox dates swapped.[191]
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The angle of Earth's axial tilt is relatively stable over long periods of time. Its axial tilt does undergo nutation; a slight, irregular motion with a main period of 18.6 years.[192] The orientation (rather than the angle) of Earth's axis also changes over time, precessing around in a complete circle over each 25,800 year cycle; this precession is the reason for the difference between a sidereal year and a tropical year. Both of these motions are caused by the varying attraction of the Sun and the Moon on Earth's equatorial bulge. The poles also migrate a few meters across Earth's surface. This polar motion has multiple, cyclical components, which collectively are termed quasiperiodic motion. In addition to an annual component to this motion, there is a 14-month cycle called the Chandler wobble. Earth's rotational velocity also varies in a phenomenon known as length-of-day variation.[193]
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In modern times, Earth's perihelion occurs around 3 January, and its aphelion around 4 July. These dates change over time due to precession and other orbital factors, which follow cyclical patterns known as Milankovitch cycles. The changing Earth–Sun distance causes an increase of about 6.9%[n 21] in solar energy reaching Earth at perihelion relative to aphelion. Because the Southern Hemisphere is tilted toward the Sun at about the same time that Earth reaches the closest approach to the Sun, the Southern Hemisphere receives slightly more energy from the Sun than does the northern over the course of a year. This effect is much less significant than the total energy change due to the axial tilt, and most of the excess energy is absorbed by the higher proportion of water in the Southern Hemisphere.[194]
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A study from 2016 suggested that Planet Nine tilted all the planets of the Solar System, including Earth, by about six degrees.[195]
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A planet that can sustain life is termed habitable, even if life did not originate there. Earth provides liquid water—an environment where complex organic molecules can assemble and interact, and sufficient energy to sustain metabolism.[196] The distance of Earth from the Sun, as well as its orbital eccentricity, rate of rotation, axial tilt, geological history, sustaining atmosphere, and magnetic field all contribute to the current climatic conditions at the surface.[197]
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A planet's life forms inhabit ecosystems, whose total is sometimes said to form a "biosphere".[198] Earth's biosphere is thought to have begun evolving about 3.5 Gya.[65] The biosphere is divided into a number of biomes, inhabited by broadly similar plants and animals.[199] On land, biomes are separated primarily by differences in latitude, height above sea level and humidity. Terrestrial biomes lying within the Arctic or Antarctic Circles, at high altitudes or in extremely arid areas are relatively barren of plant and animal life; species diversity reaches a peak in humid lowlands at equatorial latitudes.[200]
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In July 2016, scientists reported identifying a set of 355 genes from the last universal common ancestor (LUCA) of all organisms living on Earth.[201]
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Earth has resources that have been exploited by humans.[203] Those termed non-renewable resources, such as fossil fuels, only renew over geological timescales.[204]
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Large deposits of fossil fuels are obtained from Earth's crust, consisting of coal, petroleum, and natural gas.[205] These deposits are used by humans both for energy production and as feedstock for chemical production.[206] Mineral ore bodies have also been formed within the crust through a process of ore genesis, resulting from actions of magmatism, erosion, and plate tectonics.[207] These bodies form concentrated sources for many metals and other useful elements.
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Earth's biosphere produces many useful biological products for humans, including food, wood, pharmaceuticals, oxygen, and the recycling of many organic wastes. The land-based ecosystem depends upon topsoil and fresh water, and the oceanic ecosystem depends upon dissolved nutrients washed down from the land.[208] In 1980, 50.53 million km2 (19.51 million sq mi) of Earth's land surface consisted of forest and woodlands, 67.88 million km2 (26.21 million sq mi) was grasslands and pasture, and 15.01 million km2 (5.80 million sq mi) was cultivated as croplands.[209] The estimated amount of irrigated land in 1993 was 2,481,250 km2 (958,020 sq mi).[13] Humans also live on the land by using building materials to construct shelters.
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Large areas of Earth's surface are subject to extreme weather such as tropical cyclones, hurricanes, or typhoons that dominate life in those areas. From 1980 to 2000, these events caused an average of 11,800 human deaths per year.[210] Many places are subject to earthquakes, landslides, tsunamis, volcanic eruptions, tornadoes, sinkholes, blizzards, floods, droughts, wildfires, and other calamities and disasters.
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Many localized areas are subject to human-made pollution of the air and water, acid rain and toxic substances, loss of vegetation (overgrazing, deforestation, desertification), loss of wildlife, species extinction, soil degradation, soil depletion and erosion.
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There is a scientific consensus linking human activities to global warming due to industrial carbon dioxide emissions. This is predicted to produce changes such as the melting of glaciers and ice sheets, more extreme temperature ranges, significant changes in weather and a global rise in average sea levels.[211]
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Cartography, the study and practice of map-making, and geography, the study of the lands, features, inhabitants and phenomena on Earth, have historically been the disciplines devoted to depicting Earth. Surveying, the determination of locations and distances, and to a lesser extent navigation, the determination of position and direction, have developed alongside cartography and geography, providing and suitably quantifying the requisite information.
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Earth's human population reached approximately seven billion on 31 October 2011.[213] Projections indicate that the world's human population will reach 9.2 billion in 2050.[214] Most of the growth is expected to take place in developing nations. Human population density varies widely around the world, but a majority live in Asia. By 2020, 60% of the world's population is expected to be living in urban, rather than rural, areas.[215]
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68% of the land mass of the world is in the northern hemisphere.[216] Partly due to the predominance of land mass, 90% of humans live in the northern hemisphere.[217]
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It is estimated that one-eighth of Earth's surface is suitable for humans to live on – three-quarters of Earth's surface is covered by oceans, leaving one-quarter as land. Half of that land area is desert (14%),[218] high mountains (27%),[219] or other unsuitable terrains. The northernmost permanent settlement in the world is Alert, on Ellesmere Island in Nunavut, Canada.[220] (82°28′N) The southernmost is the Amundsen–Scott South Pole Station, in Antarctica, almost exactly at the South Pole. (90°S)
|
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States claim the planet's entire land surface, except for parts of Antarctica and a few other unclaimed areas. Earth has never had a planetwide government, but the United Nations is the leading worldwide intergovernmental organization.
|
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The first human to orbit Earth was Yuri Gagarin on 12 April 1961.[221] In total, about 487 people have visited outer space and reached orbit as of 30 July 2010[update], and, of these, twelve have walked on the Moon.[222][223][224] Normally, the only humans in space are those on the International Space Station. The station's crew, made up of six people, is usually replaced every six months.[225] The farthest that humans have traveled from Earth is 400,171 km (248,655 mi), achieved during the Apollo 13 mission in 1970.[226]
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The Moon is a relatively large, terrestrial, planet-like natural satellite, with a diameter about one-quarter of Earth's. It is the largest moon in the Solar System relative to the size of its planet, although Charon is larger relative to the dwarf planet Pluto. The natural satellites of other planets are also referred to as "moons", after Earth's.
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The gravitational attraction between Earth and the Moon causes tides on Earth. The same effect on the Moon has led to its tidal locking: its rotation period is the same as the time it takes to orbit Earth. As a result, it always presents the same face to the planet. As the Moon orbits Earth, different parts of its face are illuminated by the Sun, leading to the lunar phases; the dark part of the face is separated from the light part by the solar terminator.
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Due to their tidal interaction, the Moon recedes from Earth at the rate of approximately 38 mm/a (1.5 in/year). Over millions of years, these tiny modifications—and the lengthening of Earth's day by about 23 µs/yr—add up to significant changes.[227] During the Devonian period, for example, (approximately 410 Mya) there were 400 days in a year, with each day lasting 21.8 hours.[228]
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The Moon may have dramatically affected the development of life by moderating the planet's climate. Paleontological evidence and computer simulations show that Earth's axial tilt is stabilized by tidal interactions with the Moon.[229] Some theorists think that without this stabilization against the torques applied by the Sun and planets to Earth's equatorial bulge, the rotational axis might be chaotically unstable, exhibiting chaotic changes over millions of years, as appears to be the case for Mars.[230]
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Viewed from Earth, the Moon is just far enough away to have almost the same apparent-sized disk as the Sun. The angular size (or solid angle) of these two bodies match because, although the Sun's diameter is about 400 times as large as the Moon's, it is also 400 times more distant.[186] This allows total and annular solar eclipses to occur on Earth.
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The most widely accepted theory of the Moon's origin, the giant-impact hypothesis, states that it formed from the collision of a Mars-size protoplanet called Theia with the early Earth. This hypothesis explains (among other things) the Moon's relative lack of iron and volatile elements and the fact that its composition is nearly identical to that of Earth's crust.[48]
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Earth has at least five co-orbital asteroids, including 3753 Cruithne and 2002 AA29.[231][232] A trojan asteroid companion, 2010 TK7, is librating around the leading Lagrange triangular point, L4, in Earth's orbit around the Sun.[233][234]
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The tiny near-Earth asteroid 2006 RH120 makes close approaches to the Earth–Moon system roughly every twenty years. During these approaches, it can orbit Earth for brief periods of time.[235]
|
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As of April 2018[update], there are 1,886 operational, human-made satellites orbiting Earth.[5] There are also inoperative satellites, including Vanguard 1, the oldest satellite currently in orbit, and over 16,000 pieces of tracked space debris.[n 3] Earth's largest artificial satellite is the International Space Station.
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The standard astronomical symbol of Earth consists of a cross circumscribed by a circle, ,[236] representing the four corners of the world.
|
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|
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Human cultures have developed many views of the planet.[237] Earth is sometimes personified as a deity. In many cultures it is a mother goddess that is also the primary fertility deity,[238] and by the mid-20th century, the Gaia Principle compared Earth's environments and life as a single self-regulating organism leading to broad stabilization of the conditions of habitability.[239][240][241] Creation myths in many religions involve the creation of Earth by a supernatural deity or deities.[238]
|
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+
|
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+
The Hindu Vedas (1500–900 BC) refer to the Earth as Bhūgola (भूगोल), which comes from Bhū (earth, ground) and Gola (ball, sphere, globe). It means the "globe of earth". There is no direct evidence that the Hindus of that time knew that the Earth was sphere-shaped, but this name has been used extensively since the inception of the Vedas.[citation needed]
|
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+
|
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+
Scientific investigation has resulted in several culturally transformative shifts in people's view of the planet. Initial belief in a flat Earth was gradually displaced in the Greek colonies of southern Italy during the late 6th century BC by the idea of spherical Earth,[242][243][244] which was attributed to both the philosophers Pythagoras and Parmenides.[243][244] By the end of the 5th century BC, the sphericity of Earth was universally accepted among Greek intellectuals.[245] Earth was generally believed to be the center of the universe until the 16th century, when scientists first conclusively demonstrated that it was a moving object, comparable to the other planets in the Solar System.[246] Due to the efforts of influential Christian scholars and clerics such as James Ussher, who sought to determine the age of Earth through analysis of genealogies in Scripture, Westerners before the 19th century generally believed Earth to be a few thousand years old at most. It was only during the 19th century that geologists realized Earth's age was at least many millions of years.[247]
|
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|
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+
Lord Kelvin used thermodynamics to estimate the age of Earth to be between 20 million and 400 million years in 1864, sparking a vigorous debate on the subject; it was only when radioactivity and radioactive dating were discovered in the late 19th and early 20th centuries that a reliable mechanism for determining Earth's age was established, proving the planet to be billions of years old.[248][249] The perception of Earth shifted again[further explanation needed] in the 20th century when humans first viewed it from orbit, and especially with photographs of Earth returned by the Apollo program.[250][251][252]
|
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Solar System → Local Interstellar Cloud → Local Bubble → Gould Belt → Orion Arm → Milky Way → Milky Way subgroup → Local Group → Local Sheet → Virgo Supercluster → Laniakea Supercluster → Observable universe → UniverseEach arrow (→) may be read as "within" or "part of".
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1 |
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Terrorism is, in the broadest sense, the use of intentional violence for political or religious purposes.[1] It is used in this regard primarily to refer to violence during peacetime or in the context of war against non-combatants (mostly civilians and neutral military personnel).[2] The terms "terrorist" and "terrorism" originated during the French Revolution of the late 18th century[3] but gained mainstream popularity in the 1970s during the conflicts of Northern Ireland, the Basque Country and Palestine. The increased use of suicide attacks from the 1980s onwards was typified by the September 11 attacks in New York City and Washington, D.C. in 2001.
|
6 |
+
|
7 |
+
There are various different definitions of terrorism, with no universal agreement about it.[4][5] Terrorism is a charged term. It is often used with the connotation of something that is "morally wrong". Governments and non-state groups use the term to abuse or denounce opposing groups.[6][7][8][9][5] Varied political organizations have been accused of using terrorism to achieve their objectives. These include right-wing and left-wing political organizations, nationalist groups, religious groups, revolutionaries and ruling governments.[10] Legislation declaring terrorism a crime has been adopted in many states.[11] When terrorism is perpetrated by nation states, it is not considered terrorism by the state conducting it, making legality a largely grey-area issue.[12] There is no consensus as to whether or not terrorism should be regarded as a war crime.[11][13]
|
8 |
+
|
9 |
+
The Global Terrorism Database, maintained by the University of Maryland, College Park, has recorded more than 61,000 incidents of non-state terrorism, resulting in at least 140,000 deaths, between 2000 and 2014.[14]
|
10 |
+
|
11 |
+
Etymologically, the word terror is derived from the Latin verb Tersere, which later becomes Terrere. The latter form appears in European languages as early as the 12th century; its first known use in French is the word terrible in 1160. By 1356 the word terreur is in use. Terreur is the origin of the Middle English term terrour, which later becomes the modern word "terror".[15]
|
12 |
+
|
13 |
+
The term terroriste, meaning "terrorist", is first used in 1794 by the French philosopher François-Noël Babeuf, who denounces Maximilien Robespierre's Jacobin regime as a dictatorship.[16][17] In the years leading up to the Reign of Terror, the Brunswick Manifesto threatened Paris with an "exemplary, never to be forgotten vengeance: the city would be subjected to military punishment and total destruction" if the royal family was harmed, but this only increased the Revolution's will to abolish the monarchy.[18] Some writers attitudes about French Revolution grew less favorable after the French monarchy was abolished in 1792. During the Reign of Terror, which began in July 1793 and lasted thirteen months, Paris was governed by the Committee of Public safety who oversaw a regime of mass executions and public purges.[19]
|
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+
|
15 |
+
Prior to the French Revolution, ancient philosophers wrote about tyrannicide, as tyranny was seen as the greatest political threat to Greco-Roman civilization. Medieval philosophers were similarly occupied with the concept of tyranny, though the analysis of some theologians like Thomas Aquinas drew a distinction between usurpers, who could be killed by anyone, and legitimate rulers who abused their power – the latter, in Aquinas' view, could only be punished by a public authority. John of Salisbury was the first medieval Christian scholar to defend tyrannicide.[15]
|
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+
|
17 |
+
Most scholars today trace the origins of the modern tactic of terrorism to the Jewish Sicarii Zealots who attacked Romans and Jews in 1st century Palestine. They follow its development from the Persian Order of Assassins through to 19th-century anarchists. The "Reign of Terror" is usually regarded as an issue of etymology. The term terrorism has generally been used to describe violence by non-state actors rather than government violence since the 19th-century Anarchist Movement.[18][20][21]
|
18 |
+
|
19 |
+
In December 1795, Edmund Burke used the word "Terrorists" in a description of the new French government called 'Directory':[22]
|
20 |
+
|
21 |
+
At length, after a terrible struggle, the [Directory] Troops prevailed over the Citizens ... To secure them further, they have a strong corps of irregulars, ready armed. Thousands of those Hell-hounds called Terrorists, whom they had shut up in Prison on their last Revolution, as the Satellites of Tyranny, are let loose on the people.(emphasis added)
|
22 |
+
|
23 |
+
The terms "terrorism" and "terrorist" gained renewed currency in the 1970s as a result of the Israeli–Palestinian conflict,[23] the Northern Ireland conflict,[24] the Basque conflict,[25] and the operations of groups such as the Red Army Faction.[26] Leila Khaled was described as a terrorist in a 1970 issue of Life magazine.[27] A number of books on terrorism were published in the 1970s.[28] The topic came further to the fore after the 1983 Beirut barracks bombings[8] and again after the 2001 September 11 attacks[29][8][30] and the 2002 Bali bombings.[8]
|
24 |
+
|
25 |
+
There are over 109 different definitions of terrorism.[31] American political philosopher Michael Walzer in 2002 wrote: "Terrorism is the deliberate killing of innocent people, at random, to spread fear through a whole population and force the hand of its political leaders".[5] Bruce Hoffman, an American scholar, has noted that
|
26 |
+
|
27 |
+
It is not only individual agencies within the same governmental apparatus that cannot agree on a single definition of terrorism. Experts and other long-established scholars in the field are equally incapable of reaching a consensus.[32]
|
28 |
+
|
29 |
+
C. A. J. Coady has written that the question of how to define terrorism is "irresolvable" because "its natural home is in polemical, ideological and propagandist contexts".[12]
|
30 |
+
|
31 |
+
Experts disagree about "whether terrorism is wrong by definition or just wrong as a matter of fact; they disagree about whether terrorism should be defined in terms of its aims, or its methods, or both, or neither; they disagree about whether or not states can perpetrate terrorism; they even disagree about the importance or otherwise of terror for a definition of terrorism."[12]
|
32 |
+
|
33 |
+
State terrorism refers to acts of terrorism conducted by a state against its own citizens or against another state.[33]
|
34 |
+
|
35 |
+
In November 2004, a Secretary-General of the United Nations report described terrorism as any act "intended to cause death or serious bodily harm to civilians or non-combatants with the purpose of intimidating a population or compelling a government or an international organization to do or abstain from doing any act".[34] The international community has been slow to formulate a universally agreed, legally binding definition of this crime. These difficulties arise from the fact that the term "terrorism" is politically and emotionally charged.[35][36] In this regard, Angus Martyn, briefing the Australian parliament, stated,
|
36 |
+
|
37 |
+
The international community has never succeeded in developing an accepted comprehensive definition of terrorism. During the 1970s and 1980s, the United Nations attempts to define the term floundered mainly due to differences of opinion between various members about the use of violence in the context of conflicts over national liberation and self-determination.[37]
|
38 |
+
|
39 |
+
These divergences have made it impossible for the United Nations to conclude a Comprehensive Convention on International Terrorism that incorporates a single, all-encompassing, legally binding, criminal law definition of terrorism.[38] The international community has adopted a series of sectoral conventions that define and criminalize various types of terrorist activities.
|
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+
|
41 |
+
Since 1994, the United Nations General Assembly has repeatedly condemned terrorist acts using the following political description of terrorism:
|
42 |
+
|
43 |
+
Criminal acts intended or calculated to provoke a state of terror in the public, a group of persons or particular persons for political purposes are in any circumstance unjustifiable, whatever the considerations of a political, philosophical, ideological, racial, ethnic, religious or any other nature that may be invoked to justify them.[39]
|
44 |
+
|
45 |
+
Various legal systems and government agencies use different definitions of terrorism in their national legislation.
|
46 |
+
|
47 |
+
U.S. Code Title 22 Chapter 38, Section 2656f(d) defines terrorism as:
|
48 |
+
"Premeditated, politically motivated violence perpetrated against noncombatant targets by subnational groups or clandestine agents, usually intended to influence an audience".[40]
|
49 |
+
|
50 |
+
18 U.S.C. § 2331 defines "international terrorism" and "domestic terrorism" for purposes of Chapter 113B of the Code, entitled "Terrorism":
|
51 |
+
|
52 |
+
"International terrorism" means activities with the following three characteristics:
|
53 |
+
|
54 |
+
Involve violent acts or acts dangerous to human life that violate federal or state law;
|
55 |
+
|
56 |
+
Appear to be intended (i) to intimidate or coerce a civilian population; (ii) to influence the policy of a government by intimidation or coercion; or (iii) to affect the conduct of a government by mass destruction, assassination, or kidnapping; and
|
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+
|
58 |
+
occur primarily outside the territorial jurisdiction of the U.S., or transcend national boundaries in terms of the means by which they are accomplished, the persons they appear intended to intimidate or coerce, or the locale in which their perpetrators operate or seek asylum.
|
59 |
+
|
60 |
+
A definition proposed by Carsten Bockstette at the George C. Marshall European Center for Security Studies, underlines the psychological and tactical aspects of terrorism:
|
61 |
+
|
62 |
+
Terrorism is defined as political violence in an asymmetrical conflict that is designed to induce terror and psychic fear (sometimes indiscriminate) through the violent victimization and destruction of noncombatant targets (sometimes iconic symbols). Such acts are meant to send a message from an illicit clandestine organization. The purpose of terrorism is to exploit the media in order to achieve maximum attainable publicity as an amplifying force multiplier in order to influence the targeted audience(s) in order to reach short- and midterm political goals and/or desired long-term end states.[41]
|
63 |
+
|
64 |
+
Terrorists attack national symbols, which may negatively affect a government, while increasing the prestige of the given terrorist group or its ideology.[42]
|
65 |
+
|
66 |
+
Terrorist acts frequently have a political purpose.[44] Some official, governmental definitions of terrorism use the criterion of the illegitimacy or unlawfulness of the act.[45][better source needed] to distinguish between actions authorized by a government (and thus "lawful") and those of other actors, including individuals and small groups. For example, carrying out a strategic bombing on an enemy city, which is designed to affect civilian support for a cause, would not be considered terrorism if it were authorized by a government. This criterion is inherently problematic and is not universally accepted,[attribution needed] because: it denies the existence of state terrorism.[46] An associated term is violent non-state actor.[47]
|
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+
|
68 |
+
According to Ali Khan, the distinction lies ultimately in a political judgment.[48]
|
69 |
+
|
70 |
+
Having the moral charge in our vocabulary of 'something morally wrong', the term 'terrorism' is often used to abuse or denounce opposite parties, either governments or non-state-groups.[6][7][8][9][5]
|
71 |
+
|
72 |
+
Those labeled "terrorists" by their opponents rarely identify themselves as such, and typically use other terms or terms specific to their situation, such as separatist, freedom fighter, liberator, revolutionary, vigilante, militant, paramilitary, guerrilla, rebel, patriot, or any similar-meaning word in other languages and cultures. Jihadi, mujaheddin, and fedayeen are similar Arabic words that have entered the English lexicon. It is common for both parties in a conflict to describe each other as terrorists.[49]
|
73 |
+
|
74 |
+
On whether particular terrorist acts, such as killing non-combatants, can be justified as the lesser evil in a particular circumstance, philosophers have expressed different views: while, according to David Rodin, utilitarian philosophers can (in theory) conceive of cases in which the evil of terrorism is outweighed by the good that could not be achieved in a less morally costly way, in practice the "harmful effects of undermining the convention of non-combatant immunity is thought to outweigh the goods that may be achieved by particular acts of terrorism".[50] Among the non-utilitarian philosophers, Michael Walzer argued that terrorism can be morally justified in only one specific case: when "a nation or community faces the extreme threat of complete destruction and the only way it can preserve itself is by intentionally targeting non-combatants, then it is morally entitled to do so".[50][51]
|
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+
|
76 |
+
In his book Inside Terrorism Bruce Hoffman offered an explanation of why the term terrorism becomes distorted:
|
77 |
+
|
78 |
+
On one point, at least, everyone agrees: terrorism is a pejorative term. It is a word with intrinsically negative connotations that is generally applied to one's enemies and opponents, or to those with whom one disagrees and would otherwise prefer to ignore. 'What is called terrorism,' Brian Jenkins has written, 'thus seems to depend on one's point of view. Use of the term implies a moral judgment; and if one party can successfully attach the label terrorist to its opponent, then it has indirectly persuaded others to adopt its moral viewpoint.' Hence the decision to call someone or label some organization terrorist becomes almost unavoidably subjective, depending largely on whether one sympathizes with or opposes the person/group/cause concerned. If one identifies with the victim of the violence, for example, then the act is terrorism. If, however, one identifies with the perpetrator, the violent act is regarded in a more sympathetic, if not positive (or, at the worst, an ambivalent) light; and it is not terrorism.[52][53][54]
|
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+
|
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+
The pejorative connotations of the word can be summed up in the aphorism, "One man's terrorist is another man's freedom fighter".[49] This is exemplified when a group using irregular military methods is an ally of a state against a mutual enemy, but later falls out with the state and starts to use those methods against its former ally. During World War II, the Malayan People's Anti-Japanese Army was allied with the British, but during the Malayan Emergency, members of its successor (the Malayan Races Liberation Army), were branded "terrorists" by the British.[55][56] More recently, Ronald Reagan and others in the American administration frequently called the mujaheddin "freedom fighters" during the Soviet–Afghan War[57] yet twenty years later, when a new generation of Afghan men were fighting against what they perceive to be a regime installed by foreign powers, their attacks were labelled "terrorism" by George W. Bush.[58][59][60] Groups accused of terrorism understandably prefer terms reflecting legitimate military or ideological action.[61][62][63] Leading terrorism researcher Professor Martin Rudner, director of the Canadian Centre of Intelligence and Security Studies at Ottawa's Carleton University, defines "terrorist acts" as unlawful attacks for political or other ideological goals, and said:
|
81 |
+
|
82 |
+
There is the famous statement: 'One man's terrorist is another man's freedom fighter.' But that is grossly misleading. It assesses the validity of the cause when terrorism is an act. One can have a perfectly beautiful cause and yet if one commits terrorist acts, it is terrorism regardless.[64]
|
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+
|
84 |
+
Some groups, when involved in a "liberation" struggle, have been called "terrorists" by the Western governments or media. Later, these same persons, as leaders of the liberated nations, are called "statesmen" by similar organizations. Two examples of this phenomenon are the Nobel Peace Prize laureates Menachem Begin and Nelson Mandela.[65][66][67][68][69][70] WikiLeaks editor Julian Assange has been called a "terrorist" by Sarah Palin and Joe Biden.[71][72]
|
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+
|
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+
Sometimes, states that are close allies, for reasons of history, culture and politics, can disagree over whether or not members of a certain organization are terrorists. For instance, for many years, some branches of the United States government refused to label members of the Provisional Irish Republican Army (IRA) as terrorists while the IRA was using methods against one of the United States' closest allies (the United Kingdom) that the UK branded as terrorism. This was highlighted by the Quinn v. Robinson case.[73][74]
|
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+
|
88 |
+
Media outlets who wish to convey impartiality may limit their usage of "terrorist" and "terrorism" because they are loosely defined, potentially controversial in nature, and subjective terms.[75][76]
|
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+
|
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Depending on how broadly the term is defined, the roots and practice of terrorism can be traced at least to the 1st-century AD.[77] Sicarii Zealots, though some dispute whether the group, a radical offshoot of the Zealots which was active in Judaea Province at the beginning of the 1st century AD, was in fact terrorist. According to the contemporary Jewish-Roman historian Josephus, after the Zealotry rebellion against Roman rule in Judea, when some prominent Jewish collaborators with Roman rule were killed,[78][79] Judas of Galilee formed a small and more extreme offshoot of the Zealots, the Sicarii, in 6 AD.[80] Their terror was directed against Jewish "collaborators", including temple priests, Sadducees, Herodians, and other wealthy elites.[81]
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The term "terrorism" itself was originally used to describe the actions of the Jacobin Club during the "Reign of Terror" in the French Revolution. "Terror is nothing other than justice, prompt, severe, inflexible", said Jacobin leader Maximilien Robespierre. In 1795, Edmund Burke denounced the Jacobins for letting "thousands of those hell-hounds called Terrorists ... loose on the people" of France.
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In January 1858, Italian patriot Felice Orsini threw three bombs in an attempt to assassinate French Emperor Napoleon III.[82] Eight bystanders were killed and 142 injured.[82] The incident played a crucial role as an inspiration for the development of the early terrorist groups.[82]
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Arguably the first organization to utilize modern terrorist techniques was the Irish Republican Brotherhood,[83] founded in 1858 as a revolutionary Irish nationalist group[84] that carried out attacks in England.[85] The group initiated the Fenian dynamite campaign in 1881, one of the first modern terror campaigns.[86] Instead of earlier forms of terrorism based on political assassination, this campaign used modern, timed explosives with the express aim of sowing fear in the very heart of metropolitan Britain, in order to achieve political gains.[87]
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Another early terrorist group was Narodnaya Volya, founded in Russia in 1878 as a revolutionary anarchist group inspired by Sergei Nechayev and "propaganda by the deed" theorist Carlo Pisacane.[77][88][89] The group developed ideas – such as targeted killing of the 'leaders of oppression' – that were to become the hallmark of subsequent violence by small non-state groups, and they were convinced that the developing technologies of the age – such as the invention of dynamite, which they were the first anarchist group to make widespread use of[90] – enabled them to strike directly and with discrimination.[91]
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Scholars of terrorism refer to four major waves of global terrorism: "the Anarchist, the Anti-Colonial, the New Left, and the Religious. The first three have been completed and lasted around 40 years; the fourth is now in its third decade."[92]
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Terrorist incidents, 1970–2015. A total of 157,520 incidents are plotted. Orange: 1970–1999, Red: 2000–2015
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Top 10 Countries (2000–2014)
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Worldwide non-state terrorist incidents 1970–2017
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Depending on the country, the political system, and the time in history, the types of terrorism are varying.
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In early 1975, the Law Enforcement Assistant Administration in the United States formed the National Advisory Committee on Criminal Justice Standards and Goals. One of the five volumes that the committee wrote was titled Disorders and Terrorism, produced by the Task Force on Disorders and Terrorism under the direction of H. H. A. Cooper, Director of the Task Force staff.
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The Task Force defines terrorism as "a tactic or technique by means of which a violent act or the threat thereof is used for the prime purpose of creating overwhelming fear for coercive purposes". It classified disorders and terrorism into six categories:[96]
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Other sources have defined the typology of terrorism in different ways, for example, broadly classifying it into domestic terrorism and international terrorism, or using categories such as vigilante terrorism or insurgent terrorism.[100] One way the typology of terrorism may be defined:[101][102]
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Individuals and groups choose terrorism as a tactic because it can:
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Attacks on "collaborators" are used to intimidate people from cooperating with the state in order to undermine state control. This strategy was used in Ireland, in Kenya, in Algeria and in Cyprus during their independence struggles.[104]
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Stated motives for the September 11 attacks included inspiring more fighters to join the cause of repelling the United States from Muslim countries with a successful high-profile attack. The attacks prompted some criticism from domestic and international observers regarding perceived injustices in U.S. foreign policy that provoked the attacks, but the larger practical effect was that the United States government declared a War on Terror that resulted in substantial military engagements in several Muslim-majority countries. Various commentators have inferred that al-Qaeda expected a military response, and welcomed it as a provocation that would result in more Muslims fight the United States. Some commentators believe that the resulting anger and suspicion directed toward innocent Muslims living in Western countries and the indignities inflicted upon them by security forces and the general public also contributes to radicalization of new recruits.[103] Despite criticism that the Iraqi government had no involvement with the September 11 attacks, Bush declared the 2003 invasion of Iraq to be part of the War on Terror. The resulting backlash and instability enabled the rise of Islamic State of Iraq and the Levant and the temporary creation of an Islamic caliphate holding territory in Iraq and Syria, until ISIL lost its territory through military defeats.
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Attacks used to draw international attention to struggles that are otherwise unreported have included the Palestinian airplane hijackings in 1970 and the 1975 Dutch train hostage crisis.
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Specific political or social causes have included:
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Causes for right-wing terrorism have included white nationalism, ethnonationalism, fascism, anti-socialism, the anti-abortion movement, and tax resistance.
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Sometimes terrorists on the same side fight for different reasons. For example, in the Chechen–Russian conflict secular Chechens using terrorist tactics fighting for national independence are allied with radical Islamist terrorists who have arrived from other countries.[105]
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Various personal and social factors may influence the personal choice of whether or not to join a terrorist group or attempt an act of terror, including:
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A report conducted by Paul Gill, John Horgan and Paige Deckert[dubious – discuss] found that for "lone wolf" terrorists:[106]
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Ariel Merari, a psychologist who has studied the psychological profiles of suicide terrorists since 1983 through media reports that contained biographical details, interviews with the suicides’ families, and interviews with jailed would-be suicide attackers, concluded that they were unlikely to be psychologically abnormal.[107] In comparison to economic theories of criminal behaviour, Scott Atran found that suicide terrorists exhibit none of the socially dysfunctional attributes – such as fatherless, friendless, jobless situations – or suicidal symptoms. By which he means, they do not kill themselves simply out of hopelessness or a sense of 'having nothing to lose'.[108]
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Abrahm suggests that terrorist organizations do not select terrorism for its political effectiveness.[109] Individual terrorists tend to be motivated more by a desire for social solidarity with other members of their organization than by political platforms or strategic objectives, which are often murky and undefined.[109]
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Michael Mousseau shows possible relationships between the type of economy within a country and ideology associated with terrorism.[example needed][110] Many terrorists have a history of domestic violence.[111]
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Terrorism is most common in nations with intermediate political freedom, and it is least common in the most democratic nations.[112][113][114][115]
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Some examples of "terrorism" in non-democratic nations include ETA in Spain under Francisco Franco (although the group's terrorist activities increased sharply after Franco's death),[116] the Organization of Ukrainian Nationalists in pre-war Poland,[117] the Shining Path in Peru under Alberto Fujimori,[118] the Kurdistan Workers Party when Turkey was ruled by military leaders and the ANC in South Africa.[119] Democracies, such as Japan, the United Kingdom, the United States, Israel, Indonesia, India, Spain, Germany, Italy and the Philippines, have experienced domestic terrorism.
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While a democratic nation espousing civil liberties may claim a sense of higher moral ground than other regimes, an act of terrorism within such a state may cause a dilemma: whether to maintain its civil liberties and thus risk being perceived as ineffective in dealing with the problem; or alternatively to restrict its civil liberties and thus risk delegitimizing its claim of supporting civil liberties.[120] For this reason, homegrown terrorism has started to be seen as a greater threat, as stated by former CIA Director Michael Hayden.[121] This dilemma, some social theorists would conclude, may very well play into the initial plans of the acting terrorist(s); namely, to delegitimize the state and cause a systematic shift towards anarchy via the accumulation of negative sentiments towards the state system.[122]
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Terrorist acts throughout history have been performed on religious grounds with the goal to either spread or enforce a system of belief, viewpoint or opinion.[124][dubious – discuss][irrelevant citation] The validity and scope of religious terrorism is limited to an individual's view or a group's view or interpretation of that belief system's teachings.[citation needed][needs context]
|
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According to the Global Terrorism Index by the University of Maryland, College Park, religious extremism has overtaken national separatism and become the main driver of terrorist attacks around the world. Since 9/11 there has been a five-fold increase in deaths from terrorist attacks. The majority of incidents over the past several years can be tied to groups with a religious agenda. Before 2000, it was nationalist separatist terrorist organizations such as the IRA and Chechen rebels who were behind the most attacks. The number of incidents from nationalist separatist groups has remained relatively stable in the years since while religious extremism has grown. The prevalence of Islamist groups in Iraq, Afghanistan, Pakistan, Nigeria and Syria is the main driver behind these trends.[125]
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Four of the terrorist groups that have been most active since 2001 are Boko Haram, Al Qaeda, the Taliban and ISIL. These groups have been most active in Iraq, Afghanistan, Pakistan, Nigeria and Syria. 80 percent of all deaths from terrorism occurred in one of these five countries.[125]
|
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Terrorism in Pakistan has become a great problem. From the summer of 2007 until late 2009, more than 1,500 people were killed in suicide and other attacks on civilians[126] for reasons attributed to a number of causes – sectarian violence between Sunni and Shia Muslims; easy availability of guns and explosives; the existence of a "Kalashnikov culture"; an influx of ideologically driven Muslims based in or near Pakistan, who originated from various nations around the world and the subsequent war against the pro-Soviet Afghans in the 1980s which blew back into Pakistan; the presence of Islamist insurgent groups and forces such as the Taliban and Lashkar-e-Taiba. On July 2, 2013 in Lahore, 50 Muslim scholars of the Sunni Ittehad Council (SIC) issued a collective fatwa against suicide bombings, the killing of innocent people, bomb attacks, and targeted killings declaring them as Haraam or forbidden.[127]
|
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In 2015, the Southern Poverty Law Center released a report on terrorism in the United States. The report (titled The Age of the Wolf) found that during that period, "more people have been killed in America by non-Islamic domestic terrorists than jihadists."[128] The "virulent racist and anti-semitic" ideology of the ultra-right wing Christian Identity movement is usually accompanied by anti-government sentiments.[129] Adherents of Christian Identity believe that whites of European descent can be traced back to the "Lost Tribes of Israel" and many consider Jews to be the Satanic offspring of Eve and the Serpent.[129] This group has committed hate crimes, bombings and other acts of terrorism. Its influence ranges from the Ku Klux Klan and neo-Nazi groups to the anti-government militia and sovereign citizen movements.[129] Christian Identity's origins can be traced back to Anglo-Israelism, which held the view that British people were descendants of ancient Israelites. By the 1930s, the movement had been infected with anti-Semitism, and eventually Christian Identity theology diverged from traditional Anglo-Israelism, and developed what is known as the "two seed" theory.[129] According to the two-seed theory, the Jewish people are descended from Cain and the serpent (not from Shem).[129] The white European seedline is descended from the "lost tribes" of Israel. They hold themselves to "God's laws", not to "man's laws", and they do not feel bound to a government that they consider run by Jews and the New World Order.[129]
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Israel has had problems with Jewish religious terrorism. During British Mandate occupation, Irgun is among Zionist groups labelled as terrorist organization by British government and United Nations,[130] for violent terror attacks against Britons and Arabs.[131][132] Another extremist group Lehi openly declared its members as 'terrorists.'[133][134] Historian William Cleveland stated many Jews justified any action, even terrorism, taken in the cause of the creation of a Jewish state.[135] In 1995, Yigal Amir assassinated Israeli Prime Minister Yitzhak Rabin. For Amir, killing Rabin was an exemplary act that symbolized the fight against an illegitimate government that was prepared to cede Jewish Holy Land to the Palestinians. [136]
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The perpetrators of acts of terrorism can be individuals, groups, or states. According to some definitions, clandestine or semi-clandestine state actors may carry out terrorist acts outside the framework of a state of war. The most common image of terrorism is that it is carried out by small and secretive cells, highly motivated to serve a particular cause and many of the most deadly operations in recent times, such as the September 11 attacks, the London underground bombing, 2008 Mumbai attacks and the 2002 Bali bombing were planned and carried out by a close clique, composed of close friends, family members and other strong social networks. These groups benefited from the free flow of information and efficient telecommunications to succeed where others had failed.[137]
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Over the years, much research has been conducted to distill a terrorist profile to explain these individuals' actions through their psychology and socio-economic circumstances.[138] Others, like Roderick Hindery, have sought to discern profiles in the propaganda tactics used by terrorists. Some security organizations designate these groups as violent non-state actors.[citation needed] A 2007 study by economist Alan B. Krueger found that terrorists were less likely to come from an impoverished background (28 percent vs. 33 percent) and more likely to have at least a high-school education (47 percent vs. 38 percent). Another analysis found only 16 percent of terrorists came from impoverished families, vs. 30 percent of male Palestinians, and over 60 percent had gone beyond high school, vs. 15 percent of the populace.A study into the poverty-stricken conditions and whether or not,terrorists are more likely to come from here,show that people who grew up in these situations tend to show aggression and frustration towards others. This theory is largely debated for the simple fact that just because one is frustrated,does not make them a potential terrorist.[31][139]
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To avoid detection, a terrorist will look, dress, and behave normally until executing the assigned mission. Some claim that attempts to profile terrorists based on personality, physical, or sociological traits are not useful.[140] The physical and behavioral description of the terrorist could describe almost any normal person.[141] the majority of terrorist attacks are carried out by military age men, aged 16–40.[141]
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|
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Groups not part of the state apparatus of in opposition to the state are most commonly referred to as a "terrorist" in the media.
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According to the Global Terrorism Database, the most active terrorist group in the period 1970 to 2010 was Shining Path (with 4,517 attacks), followed by Farabundo Marti National Liberation Front (FMLN), Irish Republican Army (IRA), Basque Fatherland and Freedom (ETA), Revolutionary Armed Forces of Colombia (FARC), Taliban, Liberation Tigers of Tamil Eelam, New People's Army, National Liberation Army of Colombia (ELN), and Kurdistan Workers Party (PKK).[142]
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A state can sponsor terrorism by funding or harboring a terrorist group. Opinions as to which acts of violence by states consist of state-sponsored terrorism vary widely. When states provide funding for groups considered by some to be terrorist, they rarely acknowledge them as such.[143][citation needed]
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Civilization is based on a clearly defined and widely accepted yet often unarticulated hierarchy. Violence done by those higher on the hierarchy to those lower is nearly always invisible, that is, unnoticed. When it is noticed, it is fully rationalized. Violence done by those lower on the hierarchy to those higher is unthinkable, and when it does occur it is regarded with shock, horror, and the fetishization of the victims.
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As with "terrorism" the concept of "state terrorism" is controversial.[145] The Chairman of the United Nations Counter-Terrorism Committee has stated that the Committee was conscious of 12 international Conventions on the subject, and none of them referred to State terrorism, which was not an international legal concept. If States abused their power, they should be judged against international conventions dealing with war crimes, international human rights law, and international humanitarian law.[146] Former United Nations Secretary-General Kofi Annan has said that it is "time to set aside debates on so-called 'state terrorism'. The use of force by states is already thoroughly regulated under international law".[147] he made clear that, "regardless of the differences between governments on the question of the definition of terrorism, what is clear and what we can all agree on is that any deliberate attack on innocent civilians [or non-combatants], regardless of one's cause, is unacceptable and fits into the definition of terrorism."[148]
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State terrorism has been used to refer to terrorist acts committed by governmental agents or forces. This involves the use of state resources employed by a state's foreign policies, such as using its military to directly perform acts of terrorism. Professor of Political Science Michael Stohl cites the examples that include the German bombing of London, the Japanese bombing of Pearl Harbor, the British firebombing of Dresden, and the U.S. atomic bombings of Hiroshima and Nagasaki during World War II. He argues that "the use of terror tactics is common in international relations and the state has been and remains a more likely employer of terrorism within the international system than insurgents." He cites the first strike option as an example of the "terror of coercive diplomacy" as a form of this, which holds the world hostage with the implied threat of using nuclear weapons in "crisis management" and he argues that the institutionalized form of terrorism has occurred as a result of changes that took place following World War II. In this analysis, state terrorism exhibited as a form of foreign policy was shaped by the presence and use of weapons of mass destruction, and the legitimizing of such violent behavior led to an increasingly accepted form of this behavior by the state.[149][150][151]
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Charles Stewart Parnell described William Ewart Gladstone's Irish Coercion Act as terrorism in his "no-Rent manifesto" in 1881, during the Irish Land War.[152] The concept is used to describe political repressions by governments against their own civilian populations with the purpose of inciting fear. For example, taking and executing civilian hostages or extrajudicial elimination campaigns are commonly considered "terror" or terrorism, for example during the Red Terror or the Great Terror.[153] Such actions are often described as democide or genocide, which have been argued to be equivalent to state terrorism.[154] Empirical studies on this have found that democracies have little democide.[155][156] Western democracies, including the United States, have supported state terrorism[157] and mass killings,[158] with some examples being the Indonesian mass killings of 1965–66 and Operation Condor.[159][160][161]
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The connection between terrorism and tourism has been widely studied since the Luxor massacre in Egypt.[162][163] In the 1970s, the targets of terrorists were politicians and chiefs of police while now, international tourists and visitors are selected as the main targets of attacks. The attacks on the World Trade Center and the Pentagon on September 11, 2001, were the symbolic center, which marked a new epoch in the use of civil transport against the main power of the planet.[164] From this event onwards, the spaces of leisure that characterized the pride of West, were conceived as dangerous and frightful.[165][166]
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State sponsors have constituted a major form of funding; for example, Palestine Liberation Organization, Democratic Front for the Liberation of Palestine and other groups considered to be terrorist organizations, were funded by the Soviet Union.[167][168] The Stern Gang received funding from Italian Fascist officers in Beirut to undermine the British Mandate for Palestine.[169]
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"Revolutionary tax" is another major form of funding, and essentially a euphemism for "protection money".[167] Revolutionary taxes "play a secondary role as one other means of intimidating the target population".[167]
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Other major sources of funding include kidnapping for ransoms, smuggling (including wildlife smuggling),[170] fraud, and robbery.[167] The Islamic State in Iraq and the Levant has reportedly received funding "via private donations from the Gulf states".[171]
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The Financial Action Task Force is an inter-governmental body whose mandate, since October 2001, has included combating terrorist financing.[172]
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Terrorist attacks are often targeted to maximize fear and publicity, usually using explosives or poison.[174]
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Terrorist groups usually methodically plan attacks in advance, and may train participants, plant undercover agents, and raise money from supporters or through organized crime. Communications occur through modern telecommunications, or through old-fashioned methods such as couriers. There is concern about terrorist attacks employing weapons of mass destruction.
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Terrorism is a form of asymmetric warfare, and is more common when direct conventional warfare will not be effective because opposing forces vary greatly in power.[175]
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The context in which terrorist tactics are used is often a large-scale, unresolved political conflict. The type of conflict varies widely; historical examples include:
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Responses to terrorism are broad in scope. They can include re-alignments of the political spectrum and reassessments of fundamental values.
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Specific types of responses include:
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The term "counter-terrorism" has a narrower connotation, implying that it is directed at terrorist actors.
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According to a report by Dana Priest and William M. Arkin in The Washington Post, "Some 1,271 government organizations and 1,931 private companies work on programs related to counterterrorism, homeland security and intelligence in about 10,000 locations across the United States."[176]
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America's thinking on how to defeat radical Islamists is split along two very different schools of thought. Republicans, typically follow what is known as the Bush Doctrine, advocate the military model of taking the fight to the enemy and seeking to democratize the Middle East. Democrats, by contrast, generally propose the law enforcement model of better cooperation with nations and more security at home.[177] In the introduction of the U.S. Army / Marine Corps Counterinsurgency Field Manual, Sarah Sewall states the need for "U.S. forces to make securing the civilian, rather than destroying the enemy, their top priority. The civilian population is the center of gravity – the deciding factor in the struggle.... Civilian deaths create an extended family of enemies – new insurgent recruits or informants – and erode support of the host nation." Sewall sums up the book's key points on how to win this battle: "Sometimes, the more you protect your force, the less secure you may be.... Sometimes, the more force is used, the less effective it is.... The more successful the counterinsurgency is, the less force can be used and the more risk must be accepted.... Sometimes, doing nothing is the best reaction."[178] This strategy, often termed "courageous restraint", has certainly led to some success on the Middle East battlefield. However, it does not address the fact that terrorists are mostly homegrown.[177]
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Terrorism research, called terrorism and counter-terrorism research, is an interdisciplinary academic field which seeks to understand the causes of terrorism, how to prevent it as well as its impact in the broadest sense. Terrorism research can be carried out in both military and civilian contexts, for example by research centres such as the British Centre for the Study of Terrorism and Political Violence, the Norwegian Centre for Violence and Traumatic Stress Studies, and the International Centre for Counter-Terrorism (ICCT). There are several academic journals devoted to the field.[179]
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One of the agreements that promote the international legal anti-terror framework is the Code of Conduct Towards Achieving a World Free of Terrorism that was adopted at the 73rd session of the United Nations General Assembly in 2018. The Code of Conduct was initiated by Kazakhstan President Nursultan Nazarbayev. Its main goal is to implement a wide range of international commitments to counter terrorism and establish a broad global coalition towards achieving a world free of terrorism by 2045. The Code was signed by more than 70 countries.[180]
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Mass media exposure may be a primary goal of those carrying out terrorism, to expose issues that would otherwise be ignored by the media. Some consider this to be manipulation and exploitation of the media.[181]
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The Internet has created a new channel for groups to spread their messages.[182] This has created a cycle of measures and counter measures by groups in support of and in opposition to terrorist movements. The United Nations has created its own online counter-terrorism resource.[183]
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The mass media will, on occasion, censor organizations involved in terrorism (through self-restraint or regulation) to discourage further terrorism. This may encourage organizations to perform more extreme acts of terrorism to be shown in the mass media. Conversely James F. Pastor explains the significant relationship between terrorism and the media, and the underlying benefit each receives from the other.[184]
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There is always a point at which the terrorist ceases to manipulate the media gestalt. A point at which the violence may well escalate, but beyond which the terrorist has become symptomatic of the media gestalt itself. Terrorism as we ordinarily understand it is innately media-related.
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Former British Prime Minister Margaret Thatcher famously spoke of the close connection between terrorism and the media, calling publicity 'the oxygen of terrorism'.[186]
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Jones and Libicki (2008) created a list of all the terrorist groups they could find that were active between 1968 and 2006. They found 648. Of those, 136 splintered and 244 were still active in 2006.[188] Of the ones that ended, 43 percent converted to nonviolent political actions, like the Irish Republican Army in Northern Ireland. Law enforcement took out 40 percent. Ten percent won. Only 20 groups, 7 percent, were destroyed by military force.
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Forty-two groups became large enough to be labeled an insurgency; 38 of those had ended by 2006. Of those, 47 percent converted to nonviolent political actors. Only 5 percent were taken out by law enforcement. 26 percent won. 21 percent succumbed to military force.[189] Jones and Libicki concluded that military force may be necessary to deal with large insurgencies but are only occasionally decisive, because the military is too often seen as a bigger threat to civilians than the terrorists. To avoid that, the rules of engagement must be conscious of collateral damage and work to minimize it.
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Another researcher, Audrey Cronin, lists six primary ways that terrorist groups end:[190]
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The following terrorism databases are or were made publicly available for research purposes, and track specific acts of terrorism:
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The following public report and index provides a summary of key global trends and patterns in terrorism around the world
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The following publicly available resources index electronic and bibliographic resources on the subject of terrorism
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The following terrorism databases are maintained in secrecy by the United States Government for intelligence and counter-terrorism purposes:
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Jones and Libicki (2008) includes a table of 268 terrorist groups active between 1968 and 2006 with their status as of 2006: still active, splintered, converted to nonviolence, removed by law enforcement or military, or won. (These data are not in a convenient machine-readable format but are available.)
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1 |
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Terrorism is, in the broadest sense, the use of intentional violence for political or religious purposes.[1] It is used in this regard primarily to refer to violence during peacetime or in the context of war against non-combatants (mostly civilians and neutral military personnel).[2] The terms "terrorist" and "terrorism" originated during the French Revolution of the late 18th century[3] but gained mainstream popularity in the 1970s during the conflicts of Northern Ireland, the Basque Country and Palestine. The increased use of suicide attacks from the 1980s onwards was typified by the September 11 attacks in New York City and Washington, D.C. in 2001.
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There are various different definitions of terrorism, with no universal agreement about it.[4][5] Terrorism is a charged term. It is often used with the connotation of something that is "morally wrong". Governments and non-state groups use the term to abuse or denounce opposing groups.[6][7][8][9][5] Varied political organizations have been accused of using terrorism to achieve their objectives. These include right-wing and left-wing political organizations, nationalist groups, religious groups, revolutionaries and ruling governments.[10] Legislation declaring terrorism a crime has been adopted in many states.[11] When terrorism is perpetrated by nation states, it is not considered terrorism by the state conducting it, making legality a largely grey-area issue.[12] There is no consensus as to whether or not terrorism should be regarded as a war crime.[11][13]
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The Global Terrorism Database, maintained by the University of Maryland, College Park, has recorded more than 61,000 incidents of non-state terrorism, resulting in at least 140,000 deaths, between 2000 and 2014.[14]
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Etymologically, the word terror is derived from the Latin verb Tersere, which later becomes Terrere. The latter form appears in European languages as early as the 12th century; its first known use in French is the word terrible in 1160. By 1356 the word terreur is in use. Terreur is the origin of the Middle English term terrour, which later becomes the modern word "terror".[15]
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The term terroriste, meaning "terrorist", is first used in 1794 by the French philosopher François-Noël Babeuf, who denounces Maximilien Robespierre's Jacobin regime as a dictatorship.[16][17] In the years leading up to the Reign of Terror, the Brunswick Manifesto threatened Paris with an "exemplary, never to be forgotten vengeance: the city would be subjected to military punishment and total destruction" if the royal family was harmed, but this only increased the Revolution's will to abolish the monarchy.[18] Some writers attitudes about French Revolution grew less favorable after the French monarchy was abolished in 1792. During the Reign of Terror, which began in July 1793 and lasted thirteen months, Paris was governed by the Committee of Public safety who oversaw a regime of mass executions and public purges.[19]
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Prior to the French Revolution, ancient philosophers wrote about tyrannicide, as tyranny was seen as the greatest political threat to Greco-Roman civilization. Medieval philosophers were similarly occupied with the concept of tyranny, though the analysis of some theologians like Thomas Aquinas drew a distinction between usurpers, who could be killed by anyone, and legitimate rulers who abused their power – the latter, in Aquinas' view, could only be punished by a public authority. John of Salisbury was the first medieval Christian scholar to defend tyrannicide.[15]
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Most scholars today trace the origins of the modern tactic of terrorism to the Jewish Sicarii Zealots who attacked Romans and Jews in 1st century Palestine. They follow its development from the Persian Order of Assassins through to 19th-century anarchists. The "Reign of Terror" is usually regarded as an issue of etymology. The term terrorism has generally been used to describe violence by non-state actors rather than government violence since the 19th-century Anarchist Movement.[18][20][21]
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In December 1795, Edmund Burke used the word "Terrorists" in a description of the new French government called 'Directory':[22]
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At length, after a terrible struggle, the [Directory] Troops prevailed over the Citizens ... To secure them further, they have a strong corps of irregulars, ready armed. Thousands of those Hell-hounds called Terrorists, whom they had shut up in Prison on their last Revolution, as the Satellites of Tyranny, are let loose on the people.(emphasis added)
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The terms "terrorism" and "terrorist" gained renewed currency in the 1970s as a result of the Israeli–Palestinian conflict,[23] the Northern Ireland conflict,[24] the Basque conflict,[25] and the operations of groups such as the Red Army Faction.[26] Leila Khaled was described as a terrorist in a 1970 issue of Life magazine.[27] A number of books on terrorism were published in the 1970s.[28] The topic came further to the fore after the 1983 Beirut barracks bombings[8] and again after the 2001 September 11 attacks[29][8][30] and the 2002 Bali bombings.[8]
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There are over 109 different definitions of terrorism.[31] American political philosopher Michael Walzer in 2002 wrote: "Terrorism is the deliberate killing of innocent people, at random, to spread fear through a whole population and force the hand of its political leaders".[5] Bruce Hoffman, an American scholar, has noted that
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26 |
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It is not only individual agencies within the same governmental apparatus that cannot agree on a single definition of terrorism. Experts and other long-established scholars in the field are equally incapable of reaching a consensus.[32]
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C. A. J. Coady has written that the question of how to define terrorism is "irresolvable" because "its natural home is in polemical, ideological and propagandist contexts".[12]
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30 |
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Experts disagree about "whether terrorism is wrong by definition or just wrong as a matter of fact; they disagree about whether terrorism should be defined in terms of its aims, or its methods, or both, or neither; they disagree about whether or not states can perpetrate terrorism; they even disagree about the importance or otherwise of terror for a definition of terrorism."[12]
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32 |
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33 |
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State terrorism refers to acts of terrorism conducted by a state against its own citizens or against another state.[33]
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34 |
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In November 2004, a Secretary-General of the United Nations report described terrorism as any act "intended to cause death or serious bodily harm to civilians or non-combatants with the purpose of intimidating a population or compelling a government or an international organization to do or abstain from doing any act".[34] The international community has been slow to formulate a universally agreed, legally binding definition of this crime. These difficulties arise from the fact that the term "terrorism" is politically and emotionally charged.[35][36] In this regard, Angus Martyn, briefing the Australian parliament, stated,
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37 |
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The international community has never succeeded in developing an accepted comprehensive definition of terrorism. During the 1970s and 1980s, the United Nations attempts to define the term floundered mainly due to differences of opinion between various members about the use of violence in the context of conflicts over national liberation and self-determination.[37]
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These divergences have made it impossible for the United Nations to conclude a Comprehensive Convention on International Terrorism that incorporates a single, all-encompassing, legally binding, criminal law definition of terrorism.[38] The international community has adopted a series of sectoral conventions that define and criminalize various types of terrorist activities.
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Since 1994, the United Nations General Assembly has repeatedly condemned terrorist acts using the following political description of terrorism:
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Criminal acts intended or calculated to provoke a state of terror in the public, a group of persons or particular persons for political purposes are in any circumstance unjustifiable, whatever the considerations of a political, philosophical, ideological, racial, ethnic, religious or any other nature that may be invoked to justify them.[39]
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Various legal systems and government agencies use different definitions of terrorism in their national legislation.
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U.S. Code Title 22 Chapter 38, Section 2656f(d) defines terrorism as:
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"Premeditated, politically motivated violence perpetrated against noncombatant targets by subnational groups or clandestine agents, usually intended to influence an audience".[40]
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18 U.S.C. § 2331 defines "international terrorism" and "domestic terrorism" for purposes of Chapter 113B of the Code, entitled "Terrorism":
|
51 |
+
|
52 |
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"International terrorism" means activities with the following three characteristics:
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53 |
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54 |
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Involve violent acts or acts dangerous to human life that violate federal or state law;
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55 |
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Appear to be intended (i) to intimidate or coerce a civilian population; (ii) to influence the policy of a government by intimidation or coercion; or (iii) to affect the conduct of a government by mass destruction, assassination, or kidnapping; and
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occur primarily outside the territorial jurisdiction of the U.S., or transcend national boundaries in terms of the means by which they are accomplished, the persons they appear intended to intimidate or coerce, or the locale in which their perpetrators operate or seek asylum.
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A definition proposed by Carsten Bockstette at the George C. Marshall European Center for Security Studies, underlines the psychological and tactical aspects of terrorism:
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Terrorism is defined as political violence in an asymmetrical conflict that is designed to induce terror and psychic fear (sometimes indiscriminate) through the violent victimization and destruction of noncombatant targets (sometimes iconic symbols). Such acts are meant to send a message from an illicit clandestine organization. The purpose of terrorism is to exploit the media in order to achieve maximum attainable publicity as an amplifying force multiplier in order to influence the targeted audience(s) in order to reach short- and midterm political goals and/or desired long-term end states.[41]
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Terrorists attack national symbols, which may negatively affect a government, while increasing the prestige of the given terrorist group or its ideology.[42]
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Terrorist acts frequently have a political purpose.[44] Some official, governmental definitions of terrorism use the criterion of the illegitimacy or unlawfulness of the act.[45][better source needed] to distinguish between actions authorized by a government (and thus "lawful") and those of other actors, including individuals and small groups. For example, carrying out a strategic bombing on an enemy city, which is designed to affect civilian support for a cause, would not be considered terrorism if it were authorized by a government. This criterion is inherently problematic and is not universally accepted,[attribution needed] because: it denies the existence of state terrorism.[46] An associated term is violent non-state actor.[47]
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According to Ali Khan, the distinction lies ultimately in a political judgment.[48]
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Having the moral charge in our vocabulary of 'something morally wrong', the term 'terrorism' is often used to abuse or denounce opposite parties, either governments or non-state-groups.[6][7][8][9][5]
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72 |
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Those labeled "terrorists" by their opponents rarely identify themselves as such, and typically use other terms or terms specific to their situation, such as separatist, freedom fighter, liberator, revolutionary, vigilante, militant, paramilitary, guerrilla, rebel, patriot, or any similar-meaning word in other languages and cultures. Jihadi, mujaheddin, and fedayeen are similar Arabic words that have entered the English lexicon. It is common for both parties in a conflict to describe each other as terrorists.[49]
|
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On whether particular terrorist acts, such as killing non-combatants, can be justified as the lesser evil in a particular circumstance, philosophers have expressed different views: while, according to David Rodin, utilitarian philosophers can (in theory) conceive of cases in which the evil of terrorism is outweighed by the good that could not be achieved in a less morally costly way, in practice the "harmful effects of undermining the convention of non-combatant immunity is thought to outweigh the goods that may be achieved by particular acts of terrorism".[50] Among the non-utilitarian philosophers, Michael Walzer argued that terrorism can be morally justified in only one specific case: when "a nation or community faces the extreme threat of complete destruction and the only way it can preserve itself is by intentionally targeting non-combatants, then it is morally entitled to do so".[50][51]
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In his book Inside Terrorism Bruce Hoffman offered an explanation of why the term terrorism becomes distorted:
|
77 |
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|
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On one point, at least, everyone agrees: terrorism is a pejorative term. It is a word with intrinsically negative connotations that is generally applied to one's enemies and opponents, or to those with whom one disagrees and would otherwise prefer to ignore. 'What is called terrorism,' Brian Jenkins has written, 'thus seems to depend on one's point of view. Use of the term implies a moral judgment; and if one party can successfully attach the label terrorist to its opponent, then it has indirectly persuaded others to adopt its moral viewpoint.' Hence the decision to call someone or label some organization terrorist becomes almost unavoidably subjective, depending largely on whether one sympathizes with or opposes the person/group/cause concerned. If one identifies with the victim of the violence, for example, then the act is terrorism. If, however, one identifies with the perpetrator, the violent act is regarded in a more sympathetic, if not positive (or, at the worst, an ambivalent) light; and it is not terrorism.[52][53][54]
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The pejorative connotations of the word can be summed up in the aphorism, "One man's terrorist is another man's freedom fighter".[49] This is exemplified when a group using irregular military methods is an ally of a state against a mutual enemy, but later falls out with the state and starts to use those methods against its former ally. During World War II, the Malayan People's Anti-Japanese Army was allied with the British, but during the Malayan Emergency, members of its successor (the Malayan Races Liberation Army), were branded "terrorists" by the British.[55][56] More recently, Ronald Reagan and others in the American administration frequently called the mujaheddin "freedom fighters" during the Soviet–Afghan War[57] yet twenty years later, when a new generation of Afghan men were fighting against what they perceive to be a regime installed by foreign powers, their attacks were labelled "terrorism" by George W. Bush.[58][59][60] Groups accused of terrorism understandably prefer terms reflecting legitimate military or ideological action.[61][62][63] Leading terrorism researcher Professor Martin Rudner, director of the Canadian Centre of Intelligence and Security Studies at Ottawa's Carleton University, defines "terrorist acts" as unlawful attacks for political or other ideological goals, and said:
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There is the famous statement: 'One man's terrorist is another man's freedom fighter.' But that is grossly misleading. It assesses the validity of the cause when terrorism is an act. One can have a perfectly beautiful cause and yet if one commits terrorist acts, it is terrorism regardless.[64]
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Some groups, when involved in a "liberation" struggle, have been called "terrorists" by the Western governments or media. Later, these same persons, as leaders of the liberated nations, are called "statesmen" by similar organizations. Two examples of this phenomenon are the Nobel Peace Prize laureates Menachem Begin and Nelson Mandela.[65][66][67][68][69][70] WikiLeaks editor Julian Assange has been called a "terrorist" by Sarah Palin and Joe Biden.[71][72]
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Sometimes, states that are close allies, for reasons of history, culture and politics, can disagree over whether or not members of a certain organization are terrorists. For instance, for many years, some branches of the United States government refused to label members of the Provisional Irish Republican Army (IRA) as terrorists while the IRA was using methods against one of the United States' closest allies (the United Kingdom) that the UK branded as terrorism. This was highlighted by the Quinn v. Robinson case.[73][74]
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Media outlets who wish to convey impartiality may limit their usage of "terrorist" and "terrorism" because they are loosely defined, potentially controversial in nature, and subjective terms.[75][76]
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Depending on how broadly the term is defined, the roots and practice of terrorism can be traced at least to the 1st-century AD.[77] Sicarii Zealots, though some dispute whether the group, a radical offshoot of the Zealots which was active in Judaea Province at the beginning of the 1st century AD, was in fact terrorist. According to the contemporary Jewish-Roman historian Josephus, after the Zealotry rebellion against Roman rule in Judea, when some prominent Jewish collaborators with Roman rule were killed,[78][79] Judas of Galilee formed a small and more extreme offshoot of the Zealots, the Sicarii, in 6 AD.[80] Their terror was directed against Jewish "collaborators", including temple priests, Sadducees, Herodians, and other wealthy elites.[81]
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The term "terrorism" itself was originally used to describe the actions of the Jacobin Club during the "Reign of Terror" in the French Revolution. "Terror is nothing other than justice, prompt, severe, inflexible", said Jacobin leader Maximilien Robespierre. In 1795, Edmund Burke denounced the Jacobins for letting "thousands of those hell-hounds called Terrorists ... loose on the people" of France.
|
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In January 1858, Italian patriot Felice Orsini threw three bombs in an attempt to assassinate French Emperor Napoleon III.[82] Eight bystanders were killed and 142 injured.[82] The incident played a crucial role as an inspiration for the development of the early terrorist groups.[82]
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|
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Arguably the first organization to utilize modern terrorist techniques was the Irish Republican Brotherhood,[83] founded in 1858 as a revolutionary Irish nationalist group[84] that carried out attacks in England.[85] The group initiated the Fenian dynamite campaign in 1881, one of the first modern terror campaigns.[86] Instead of earlier forms of terrorism based on political assassination, this campaign used modern, timed explosives with the express aim of sowing fear in the very heart of metropolitan Britain, in order to achieve political gains.[87]
|
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|
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Another early terrorist group was Narodnaya Volya, founded in Russia in 1878 as a revolutionary anarchist group inspired by Sergei Nechayev and "propaganda by the deed" theorist Carlo Pisacane.[77][88][89] The group developed ideas – such as targeted killing of the 'leaders of oppression' – that were to become the hallmark of subsequent violence by small non-state groups, and they were convinced that the developing technologies of the age – such as the invention of dynamite, which they were the first anarchist group to make widespread use of[90] – enabled them to strike directly and with discrimination.[91]
|
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|
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Scholars of terrorism refer to four major waves of global terrorism: "the Anarchist, the Anti-Colonial, the New Left, and the Religious. The first three have been completed and lasted around 40 years; the fourth is now in its third decade."[92]
|
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|
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Terrorist incidents, 1970–2015. A total of 157,520 incidents are plotted. Orange: 1970–1999, Red: 2000–2015
|
103 |
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|
104 |
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Top 10 Countries (2000–2014)
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105 |
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|
106 |
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Worldwide non-state terrorist incidents 1970–2017
|
107 |
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|
108 |
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Depending on the country, the political system, and the time in history, the types of terrorism are varying.
|
109 |
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|
110 |
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In early 1975, the Law Enforcement Assistant Administration in the United States formed the National Advisory Committee on Criminal Justice Standards and Goals. One of the five volumes that the committee wrote was titled Disorders and Terrorism, produced by the Task Force on Disorders and Terrorism under the direction of H. H. A. Cooper, Director of the Task Force staff.
|
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|
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The Task Force defines terrorism as "a tactic or technique by means of which a violent act or the threat thereof is used for the prime purpose of creating overwhelming fear for coercive purposes". It classified disorders and terrorism into six categories:[96]
|
113 |
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|
114 |
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Other sources have defined the typology of terrorism in different ways, for example, broadly classifying it into domestic terrorism and international terrorism, or using categories such as vigilante terrorism or insurgent terrorism.[100] One way the typology of terrorism may be defined:[101][102]
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|
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Individuals and groups choose terrorism as a tactic because it can:
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|
118 |
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Attacks on "collaborators" are used to intimidate people from cooperating with the state in order to undermine state control. This strategy was used in Ireland, in Kenya, in Algeria and in Cyprus during their independence struggles.[104]
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119 |
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|
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Stated motives for the September 11 attacks included inspiring more fighters to join the cause of repelling the United States from Muslim countries with a successful high-profile attack. The attacks prompted some criticism from domestic and international observers regarding perceived injustices in U.S. foreign policy that provoked the attacks, but the larger practical effect was that the United States government declared a War on Terror that resulted in substantial military engagements in several Muslim-majority countries. Various commentators have inferred that al-Qaeda expected a military response, and welcomed it as a provocation that would result in more Muslims fight the United States. Some commentators believe that the resulting anger and suspicion directed toward innocent Muslims living in Western countries and the indignities inflicted upon them by security forces and the general public also contributes to radicalization of new recruits.[103] Despite criticism that the Iraqi government had no involvement with the September 11 attacks, Bush declared the 2003 invasion of Iraq to be part of the War on Terror. The resulting backlash and instability enabled the rise of Islamic State of Iraq and the Levant and the temporary creation of an Islamic caliphate holding territory in Iraq and Syria, until ISIL lost its territory through military defeats.
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|
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Attacks used to draw international attention to struggles that are otherwise unreported have included the Palestinian airplane hijackings in 1970 and the 1975 Dutch train hostage crisis.
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+
|
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Specific political or social causes have included:
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125 |
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|
126 |
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Causes for right-wing terrorism have included white nationalism, ethnonationalism, fascism, anti-socialism, the anti-abortion movement, and tax resistance.
|
127 |
+
|
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Sometimes terrorists on the same side fight for different reasons. For example, in the Chechen–Russian conflict secular Chechens using terrorist tactics fighting for national independence are allied with radical Islamist terrorists who have arrived from other countries.[105]
|
129 |
+
|
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Various personal and social factors may influence the personal choice of whether or not to join a terrorist group or attempt an act of terror, including:
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131 |
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|
132 |
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A report conducted by Paul Gill, John Horgan and Paige Deckert[dubious – discuss] found that for "lone wolf" terrorists:[106]
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Ariel Merari, a psychologist who has studied the psychological profiles of suicide terrorists since 1983 through media reports that contained biographical details, interviews with the suicides’ families, and interviews with jailed would-be suicide attackers, concluded that they were unlikely to be psychologically abnormal.[107] In comparison to economic theories of criminal behaviour, Scott Atran found that suicide terrorists exhibit none of the socially dysfunctional attributes – such as fatherless, friendless, jobless situations – or suicidal symptoms. By which he means, they do not kill themselves simply out of hopelessness or a sense of 'having nothing to lose'.[108]
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135 |
+
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136 |
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Abrahm suggests that terrorist organizations do not select terrorism for its political effectiveness.[109] Individual terrorists tend to be motivated more by a desire for social solidarity with other members of their organization than by political platforms or strategic objectives, which are often murky and undefined.[109]
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|
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Michael Mousseau shows possible relationships between the type of economy within a country and ideology associated with terrorism.[example needed][110] Many terrorists have a history of domestic violence.[111]
|
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+
|
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Terrorism is most common in nations with intermediate political freedom, and it is least common in the most democratic nations.[112][113][114][115]
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141 |
+
|
142 |
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Some examples of "terrorism" in non-democratic nations include ETA in Spain under Francisco Franco (although the group's terrorist activities increased sharply after Franco's death),[116] the Organization of Ukrainian Nationalists in pre-war Poland,[117] the Shining Path in Peru under Alberto Fujimori,[118] the Kurdistan Workers Party when Turkey was ruled by military leaders and the ANC in South Africa.[119] Democracies, such as Japan, the United Kingdom, the United States, Israel, Indonesia, India, Spain, Germany, Italy and the Philippines, have experienced domestic terrorism.
|
143 |
+
|
144 |
+
While a democratic nation espousing civil liberties may claim a sense of higher moral ground than other regimes, an act of terrorism within such a state may cause a dilemma: whether to maintain its civil liberties and thus risk being perceived as ineffective in dealing with the problem; or alternatively to restrict its civil liberties and thus risk delegitimizing its claim of supporting civil liberties.[120] For this reason, homegrown terrorism has started to be seen as a greater threat, as stated by former CIA Director Michael Hayden.[121] This dilemma, some social theorists would conclude, may very well play into the initial plans of the acting terrorist(s); namely, to delegitimize the state and cause a systematic shift towards anarchy via the accumulation of negative sentiments towards the state system.[122]
|
145 |
+
|
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Terrorist acts throughout history have been performed on religious grounds with the goal to either spread or enforce a system of belief, viewpoint or opinion.[124][dubious – discuss][irrelevant citation] The validity and scope of religious terrorism is limited to an individual's view or a group's view or interpretation of that belief system's teachings.[citation needed][needs context]
|
147 |
+
|
148 |
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According to the Global Terrorism Index by the University of Maryland, College Park, religious extremism has overtaken national separatism and become the main driver of terrorist attacks around the world. Since 9/11 there has been a five-fold increase in deaths from terrorist attacks. The majority of incidents over the past several years can be tied to groups with a religious agenda. Before 2000, it was nationalist separatist terrorist organizations such as the IRA and Chechen rebels who were behind the most attacks. The number of incidents from nationalist separatist groups has remained relatively stable in the years since while religious extremism has grown. The prevalence of Islamist groups in Iraq, Afghanistan, Pakistan, Nigeria and Syria is the main driver behind these trends.[125]
|
149 |
+
|
150 |
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Four of the terrorist groups that have been most active since 2001 are Boko Haram, Al Qaeda, the Taliban and ISIL. These groups have been most active in Iraq, Afghanistan, Pakistan, Nigeria and Syria. 80 percent of all deaths from terrorism occurred in one of these five countries.[125]
|
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|
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Terrorism in Pakistan has become a great problem. From the summer of 2007 until late 2009, more than 1,500 people were killed in suicide and other attacks on civilians[126] for reasons attributed to a number of causes – sectarian violence between Sunni and Shia Muslims; easy availability of guns and explosives; the existence of a "Kalashnikov culture"; an influx of ideologically driven Muslims based in or near Pakistan, who originated from various nations around the world and the subsequent war against the pro-Soviet Afghans in the 1980s which blew back into Pakistan; the presence of Islamist insurgent groups and forces such as the Taliban and Lashkar-e-Taiba. On July 2, 2013 in Lahore, 50 Muslim scholars of the Sunni Ittehad Council (SIC) issued a collective fatwa against suicide bombings, the killing of innocent people, bomb attacks, and targeted killings declaring them as Haraam or forbidden.[127]
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In 2015, the Southern Poverty Law Center released a report on terrorism in the United States. The report (titled The Age of the Wolf) found that during that period, "more people have been killed in America by non-Islamic domestic terrorists than jihadists."[128] The "virulent racist and anti-semitic" ideology of the ultra-right wing Christian Identity movement is usually accompanied by anti-government sentiments.[129] Adherents of Christian Identity believe that whites of European descent can be traced back to the "Lost Tribes of Israel" and many consider Jews to be the Satanic offspring of Eve and the Serpent.[129] This group has committed hate crimes, bombings and other acts of terrorism. Its influence ranges from the Ku Klux Klan and neo-Nazi groups to the anti-government militia and sovereign citizen movements.[129] Christian Identity's origins can be traced back to Anglo-Israelism, which held the view that British people were descendants of ancient Israelites. By the 1930s, the movement had been infected with anti-Semitism, and eventually Christian Identity theology diverged from traditional Anglo-Israelism, and developed what is known as the "two seed" theory.[129] According to the two-seed theory, the Jewish people are descended from Cain and the serpent (not from Shem).[129] The white European seedline is descended from the "lost tribes" of Israel. They hold themselves to "God's laws", not to "man's laws", and they do not feel bound to a government that they consider run by Jews and the New World Order.[129]
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Israel has had problems with Jewish religious terrorism. During British Mandate occupation, Irgun is among Zionist groups labelled as terrorist organization by British government and United Nations,[130] for violent terror attacks against Britons and Arabs.[131][132] Another extremist group Lehi openly declared its members as 'terrorists.'[133][134] Historian William Cleveland stated many Jews justified any action, even terrorism, taken in the cause of the creation of a Jewish state.[135] In 1995, Yigal Amir assassinated Israeli Prime Minister Yitzhak Rabin. For Amir, killing Rabin was an exemplary act that symbolized the fight against an illegitimate government that was prepared to cede Jewish Holy Land to the Palestinians. [136]
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The perpetrators of acts of terrorism can be individuals, groups, or states. According to some definitions, clandestine or semi-clandestine state actors may carry out terrorist acts outside the framework of a state of war. The most common image of terrorism is that it is carried out by small and secretive cells, highly motivated to serve a particular cause and many of the most deadly operations in recent times, such as the September 11 attacks, the London underground bombing, 2008 Mumbai attacks and the 2002 Bali bombing were planned and carried out by a close clique, composed of close friends, family members and other strong social networks. These groups benefited from the free flow of information and efficient telecommunications to succeed where others had failed.[137]
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Over the years, much research has been conducted to distill a terrorist profile to explain these individuals' actions through their psychology and socio-economic circumstances.[138] Others, like Roderick Hindery, have sought to discern profiles in the propaganda tactics used by terrorists. Some security organizations designate these groups as violent non-state actors.[citation needed] A 2007 study by economist Alan B. Krueger found that terrorists were less likely to come from an impoverished background (28 percent vs. 33 percent) and more likely to have at least a high-school education (47 percent vs. 38 percent). Another analysis found only 16 percent of terrorists came from impoverished families, vs. 30 percent of male Palestinians, and over 60 percent had gone beyond high school, vs. 15 percent of the populace.A study into the poverty-stricken conditions and whether or not,terrorists are more likely to come from here,show that people who grew up in these situations tend to show aggression and frustration towards others. This theory is largely debated for the simple fact that just because one is frustrated,does not make them a potential terrorist.[31][139]
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To avoid detection, a terrorist will look, dress, and behave normally until executing the assigned mission. Some claim that attempts to profile terrorists based on personality, physical, or sociological traits are not useful.[140] The physical and behavioral description of the terrorist could describe almost any normal person.[141] the majority of terrorist attacks are carried out by military age men, aged 16–40.[141]
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Groups not part of the state apparatus of in opposition to the state are most commonly referred to as a "terrorist" in the media.
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According to the Global Terrorism Database, the most active terrorist group in the period 1970 to 2010 was Shining Path (with 4,517 attacks), followed by Farabundo Marti National Liberation Front (FMLN), Irish Republican Army (IRA), Basque Fatherland and Freedom (ETA), Revolutionary Armed Forces of Colombia (FARC), Taliban, Liberation Tigers of Tamil Eelam, New People's Army, National Liberation Army of Colombia (ELN), and Kurdistan Workers Party (PKK).[142]
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A state can sponsor terrorism by funding or harboring a terrorist group. Opinions as to which acts of violence by states consist of state-sponsored terrorism vary widely. When states provide funding for groups considered by some to be terrorist, they rarely acknowledge them as such.[143][citation needed]
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Civilization is based on a clearly defined and widely accepted yet often unarticulated hierarchy. Violence done by those higher on the hierarchy to those lower is nearly always invisible, that is, unnoticed. When it is noticed, it is fully rationalized. Violence done by those lower on the hierarchy to those higher is unthinkable, and when it does occur it is regarded with shock, horror, and the fetishization of the victims.
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As with "terrorism" the concept of "state terrorism" is controversial.[145] The Chairman of the United Nations Counter-Terrorism Committee has stated that the Committee was conscious of 12 international Conventions on the subject, and none of them referred to State terrorism, which was not an international legal concept. If States abused their power, they should be judged against international conventions dealing with war crimes, international human rights law, and international humanitarian law.[146] Former United Nations Secretary-General Kofi Annan has said that it is "time to set aside debates on so-called 'state terrorism'. The use of force by states is already thoroughly regulated under international law".[147] he made clear that, "regardless of the differences between governments on the question of the definition of terrorism, what is clear and what we can all agree on is that any deliberate attack on innocent civilians [or non-combatants], regardless of one's cause, is unacceptable and fits into the definition of terrorism."[148]
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State terrorism has been used to refer to terrorist acts committed by governmental agents or forces. This involves the use of state resources employed by a state's foreign policies, such as using its military to directly perform acts of terrorism. Professor of Political Science Michael Stohl cites the examples that include the German bombing of London, the Japanese bombing of Pearl Harbor, the British firebombing of Dresden, and the U.S. atomic bombings of Hiroshima and Nagasaki during World War II. He argues that "the use of terror tactics is common in international relations and the state has been and remains a more likely employer of terrorism within the international system than insurgents." He cites the first strike option as an example of the "terror of coercive diplomacy" as a form of this, which holds the world hostage with the implied threat of using nuclear weapons in "crisis management" and he argues that the institutionalized form of terrorism has occurred as a result of changes that took place following World War II. In this analysis, state terrorism exhibited as a form of foreign policy was shaped by the presence and use of weapons of mass destruction, and the legitimizing of such violent behavior led to an increasingly accepted form of this behavior by the state.[149][150][151]
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Charles Stewart Parnell described William Ewart Gladstone's Irish Coercion Act as terrorism in his "no-Rent manifesto" in 1881, during the Irish Land War.[152] The concept is used to describe political repressions by governments against their own civilian populations with the purpose of inciting fear. For example, taking and executing civilian hostages or extrajudicial elimination campaigns are commonly considered "terror" or terrorism, for example during the Red Terror or the Great Terror.[153] Such actions are often described as democide or genocide, which have been argued to be equivalent to state terrorism.[154] Empirical studies on this have found that democracies have little democide.[155][156] Western democracies, including the United States, have supported state terrorism[157] and mass killings,[158] with some examples being the Indonesian mass killings of 1965–66 and Operation Condor.[159][160][161]
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The connection between terrorism and tourism has been widely studied since the Luxor massacre in Egypt.[162][163] In the 1970s, the targets of terrorists were politicians and chiefs of police while now, international tourists and visitors are selected as the main targets of attacks. The attacks on the World Trade Center and the Pentagon on September 11, 2001, were the symbolic center, which marked a new epoch in the use of civil transport against the main power of the planet.[164] From this event onwards, the spaces of leisure that characterized the pride of West, were conceived as dangerous and frightful.[165][166]
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State sponsors have constituted a major form of funding; for example, Palestine Liberation Organization, Democratic Front for the Liberation of Palestine and other groups considered to be terrorist organizations, were funded by the Soviet Union.[167][168] The Stern Gang received funding from Italian Fascist officers in Beirut to undermine the British Mandate for Palestine.[169]
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"Revolutionary tax" is another major form of funding, and essentially a euphemism for "protection money".[167] Revolutionary taxes "play a secondary role as one other means of intimidating the target population".[167]
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Other major sources of funding include kidnapping for ransoms, smuggling (including wildlife smuggling),[170] fraud, and robbery.[167] The Islamic State in Iraq and the Levant has reportedly received funding "via private donations from the Gulf states".[171]
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The Financial Action Task Force is an inter-governmental body whose mandate, since October 2001, has included combating terrorist financing.[172]
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Terrorist attacks are often targeted to maximize fear and publicity, usually using explosives or poison.[174]
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Terrorist groups usually methodically plan attacks in advance, and may train participants, plant undercover agents, and raise money from supporters or through organized crime. Communications occur through modern telecommunications, or through old-fashioned methods such as couriers. There is concern about terrorist attacks employing weapons of mass destruction.
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Terrorism is a form of asymmetric warfare, and is more common when direct conventional warfare will not be effective because opposing forces vary greatly in power.[175]
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The context in which terrorist tactics are used is often a large-scale, unresolved political conflict. The type of conflict varies widely; historical examples include:
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Responses to terrorism are broad in scope. They can include re-alignments of the political spectrum and reassessments of fundamental values.
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Specific types of responses include:
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The term "counter-terrorism" has a narrower connotation, implying that it is directed at terrorist actors.
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According to a report by Dana Priest and William M. Arkin in The Washington Post, "Some 1,271 government organizations and 1,931 private companies work on programs related to counterterrorism, homeland security and intelligence in about 10,000 locations across the United States."[176]
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America's thinking on how to defeat radical Islamists is split along two very different schools of thought. Republicans, typically follow what is known as the Bush Doctrine, advocate the military model of taking the fight to the enemy and seeking to democratize the Middle East. Democrats, by contrast, generally propose the law enforcement model of better cooperation with nations and more security at home.[177] In the introduction of the U.S. Army / Marine Corps Counterinsurgency Field Manual, Sarah Sewall states the need for "U.S. forces to make securing the civilian, rather than destroying the enemy, their top priority. The civilian population is the center of gravity – the deciding factor in the struggle.... Civilian deaths create an extended family of enemies – new insurgent recruits or informants – and erode support of the host nation." Sewall sums up the book's key points on how to win this battle: "Sometimes, the more you protect your force, the less secure you may be.... Sometimes, the more force is used, the less effective it is.... The more successful the counterinsurgency is, the less force can be used and the more risk must be accepted.... Sometimes, doing nothing is the best reaction."[178] This strategy, often termed "courageous restraint", has certainly led to some success on the Middle East battlefield. However, it does not address the fact that terrorists are mostly homegrown.[177]
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Terrorism research, called terrorism and counter-terrorism research, is an interdisciplinary academic field which seeks to understand the causes of terrorism, how to prevent it as well as its impact in the broadest sense. Terrorism research can be carried out in both military and civilian contexts, for example by research centres such as the British Centre for the Study of Terrorism and Political Violence, the Norwegian Centre for Violence and Traumatic Stress Studies, and the International Centre for Counter-Terrorism (ICCT). There are several academic journals devoted to the field.[179]
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One of the agreements that promote the international legal anti-terror framework is the Code of Conduct Towards Achieving a World Free of Terrorism that was adopted at the 73rd session of the United Nations General Assembly in 2018. The Code of Conduct was initiated by Kazakhstan President Nursultan Nazarbayev. Its main goal is to implement a wide range of international commitments to counter terrorism and establish a broad global coalition towards achieving a world free of terrorism by 2045. The Code was signed by more than 70 countries.[180]
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Mass media exposure may be a primary goal of those carrying out terrorism, to expose issues that would otherwise be ignored by the media. Some consider this to be manipulation and exploitation of the media.[181]
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The Internet has created a new channel for groups to spread their messages.[182] This has created a cycle of measures and counter measures by groups in support of and in opposition to terrorist movements. The United Nations has created its own online counter-terrorism resource.[183]
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The mass media will, on occasion, censor organizations involved in terrorism (through self-restraint or regulation) to discourage further terrorism. This may encourage organizations to perform more extreme acts of terrorism to be shown in the mass media. Conversely James F. Pastor explains the significant relationship between terrorism and the media, and the underlying benefit each receives from the other.[184]
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There is always a point at which the terrorist ceases to manipulate the media gestalt. A point at which the violence may well escalate, but beyond which the terrorist has become symptomatic of the media gestalt itself. Terrorism as we ordinarily understand it is innately media-related.
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Former British Prime Minister Margaret Thatcher famously spoke of the close connection between terrorism and the media, calling publicity 'the oxygen of terrorism'.[186]
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Jones and Libicki (2008) created a list of all the terrorist groups they could find that were active between 1968 and 2006. They found 648. Of those, 136 splintered and 244 were still active in 2006.[188] Of the ones that ended, 43 percent converted to nonviolent political actions, like the Irish Republican Army in Northern Ireland. Law enforcement took out 40 percent. Ten percent won. Only 20 groups, 7 percent, were destroyed by military force.
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Forty-two groups became large enough to be labeled an insurgency; 38 of those had ended by 2006. Of those, 47 percent converted to nonviolent political actors. Only 5 percent were taken out by law enforcement. 26 percent won. 21 percent succumbed to military force.[189] Jones and Libicki concluded that military force may be necessary to deal with large insurgencies but are only occasionally decisive, because the military is too often seen as a bigger threat to civilians than the terrorists. To avoid that, the rules of engagement must be conscious of collateral damage and work to minimize it.
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Another researcher, Audrey Cronin, lists six primary ways that terrorist groups end:[190]
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The following terrorism databases are or were made publicly available for research purposes, and track specific acts of terrorism:
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The following public report and index provides a summary of key global trends and patterns in terrorism around the world
|
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The following publicly available resources index electronic and bibliographic resources on the subject of terrorism
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The following terrorism databases are maintained in secrecy by the United States Government for intelligence and counter-terrorism purposes:
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Jones and Libicki (2008) includes a table of 268 terrorist groups active between 1968 and 2006 with their status as of 2006: still active, splintered, converted to nonviolence, removed by law enforcement or military, or won. (These data are not in a convenient machine-readable format but are available.)
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1 |
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Testicle or testis (plural testes) is the male reproductive gland or gonad in all animals, including humans. It is homologous to the female ovary. The functions of the testes are to produce both sperm and androgens, primarily testosterone. Testosterone release is controlled by the anterior pituitary luteinizing hormone; whereas sperm production is controlled both by the anterior pituitary follicle-stimulating hormone and gonadal testosterone.
|
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Males have two testicles of similar size contained within the scrotum, which is an extension of the abdominal wall. Scrotal asymmetry is not unusual: one testicle extends farther down into the scrotum than the other due to differences in the anatomy of the vasculature.
|
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|
7 |
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The volume of the testicle can be estimated by palpating it and comparing it to ellipsoids of known sizes. Another method is to use calipers (an orchidometer) or a ruler either on the person or on an ultrasound image to obtain the three measurements of the x, y, and z axes (length, depth and width). These measurements can then be used to calculate the volume, using the formula for the volume of an ellipsoid:
|
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|
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The dimensions of the average adult testicle are up to 2 inches long, 0.8 inches in breadth, and 1.2 inches in height (5 × 2 × 3 cm). The Tanner scale for the maturity of male genitals assigns a maturity stage to the calculated volume ranging from stage I, a volume of less than 1.5 ml; to stage V, a volume greater than 20 ml. Normal volume is 15 to 25 ml; the average is 18 cm³ per testis (range 12 cm³ to 30 cm³.[1]
|
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|
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The testes are covered by a tough membranous shell called the tunica albuginea. Within the testes are very fine coiled tubes called seminiferous tubules. The tubules are lined with a layer of cells (germ cells) that develop from puberty through old age into sperm cells (also known as spermatozoa or male gametes). The developing sperm travel through the seminiferous tubules to the rete testis located in the mediastinum testis, to the efferent ducts, and then to the epididymis where newly created sperm cells mature (see spermatogenesis). The sperm move into the vas deferens, and are eventually expelled through the urethra and out of the urethral orifice through muscular contractions.
|
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|
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Blood supply and lymphatic drainage of the testes and scrotum are distinct:
|
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Many anatomical features of the adult testis reflect its developmental origin in the abdomen. The layers of tissue enclosing each testicle are derived from the layers of the anterior abdominal wall. Notably, the cremasteric muscle arises from the internal oblique muscle.
|
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|
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Large molecules cannot pass from the blood into the lumen of a seminiferous tubule due to the presence of tight junctions between adjacent Sertoli cells. The spermatogonia are in the basal compartment (deep to the level of the tight junctions) and the more mature forms such as primary and secondary spermatocytes and spermatids are in the adluminal compartment.
|
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The function of the blood–testis barrier may be to prevent an auto-immune reaction. Mature sperm (and their antigens) arise long after immune tolerance is established in infancy. Therefore, since sperm are antigenically different from self tissue, a male animal can react immunologically to his own sperm. In fact, he is capable of making antibodies against them.
|
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Injection of sperm antigens causes inflammation of the testis (auto-immune orchitis) and reduced fertility. Thus, the blood–testis barrier may reduce the likelihood that sperm proteins will induce an immune response, reducing fertility and so progeny.
|
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+
Spermatogenesis is enhanced at temperatures slightly less than core body temperature.[5] The spermatogenesis is less efficient at lower and higher temperatures than 33 °C.[5] Because the testes are located outside the body, the smooth tissue of the scrotum can move them closer or further away from the body.[5] The temperature of the testes is maintained at 35 degrees Celsius (95 degrees Fahrenheit), i.e. two degrees below the body temperature of 37 degrees Celsius (98.6 degrees Fahrenheit). Higher temperatures affect spermatogenesis.[6] There are a number of mechanisms to maintain the testes at the optimum temperature.[5]
|
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|
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+
The cremasteric muscle is part of the spermatic cord. When this muscle contracts, the cord is shortened and the testicle is moved closer up toward the body, which provides slightly more warmth to maintain optimal testicular temperature. When cooling is required, the cremasteric muscle relaxes and the testicle is lowered away from the warm body and is able to cool. Contraction also occurs in response to stress (the testicles rise up toward the body in an effort to protect them in a fight).
|
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The cremaster muscle can reflexively raise each testicle individually if properly triggered. This phenomenon is known as the cremasteric reflex. The testicles can also be lifted voluntarily using the pubococcygeus muscle, which partially activates related muscles.
|
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|
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+
The human genome includes approximately 20,000 protein coding genes: 80% of these genes are expressed in adult testes.[7] The testes have the highest fraction of tissue type-specific genes compared to other organs and tissues:[8] about 1000 of them are highly specific for the testes,[7] and about 2,200 show an elevated pattern of expression here. A majority of these genes encode for proteins that are expressed in the seminiferous tubules and have functions related to spermatogenesis.[9][8] Sperm cells express proteins that result in the development of flagella; these same proteins are expressed in the female in cells lining the fallopian tube, and cause the development of cilia. In other words, sperm cell flagella and Fallopian tube cilia are homologous structures. The testis-specific proteins that show the highest level of expression are protamines.
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|
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+
There are two phases in which the testes grow substantially; namely in embryonic and pubertal age.
|
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|
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During mammalian development, the gonads are at first capable of becoming either ovaries or testes.[10] In humans, starting at about week 4 the gonadal rudiments are present within the intermediate mesoderm adjacent to the developing kidneys. At about week 6, sex cords develop within the forming testes. These are made up of early Sertoli cells that surround and nurture the germ cells that migrate into the gonads shortly before sex determination begins. In males, the sex-specific gene SRY that is found on the Y-chromosome initiates sex determination by downstream regulation of sex-determining factors, (such as GATA4, SOX9 and AMH), which leads to development of the male phenotype, including directing development of the early bipotential gonad down the male path of development.
|
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|
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+
Testes follow the "path of descent" from high in the posterior fetal abdomen to the inguinal ring and beyond to the inguinal canal and into the scrotum. In most cases (97% full-term, 70% preterm), both testes have descended by birth. In most other cases, only one testis fails to descend (cryptorchidism) and that will probably express itself within a year.
|
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|
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+
The testes grow in response to the start of spermatogenesis. Size depends on lytic function, sperm production (amount of spermatogenesis present in testis), interstitial fluid, and Sertoli cell fluid production. After puberty, the volume of the testes can be increased by over 500% as compared to the pre-pubertal size.[citation needed] Testicles are fully descended before one reaches puberty.
|
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|
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+
|
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|
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+
Testicular enlargement is an unspecific sign of various testicular diseases, and can be defined as a testicular size of more than 5 cm (long axis) x 3 cm (short axis).[14]
|
42 |
+
|
43 |
+
Blue balls is a slang term for a temporary fluid congestion in the testicles and prostate region caused by prolonged sexual arousal.
|
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+
|
45 |
+
Testicular prostheses are available to mimic the appearance and feel of one or both testicles, when absent as from injury or as treatment in association to gender dysphoria. There have also been some instances of their implantation in dogs.[15]
|
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+
|
47 |
+
To some extent, it is possible to change testicular size. Short of direct injury or subjecting them to adverse conditions, e.g., higher temperature than they are normally accustomed to, they can be shrunk by competing against their intrinsic hormonal function through the use of externally administered steroidal hormones. Steroids taken for muscle enhancement (especially anabolic steroids) often have the undesired side effect of testicular shrinkage.
|
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|
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+
Similarly, stimulation of testicular functions via gonadotropic-like hormones may enlarge their size. Testes may shrink or atrophy during hormone replacement therapy or through chemical castration.
|
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+
|
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+
In all cases, the loss in testes volume corresponds with a loss of spermatogenesis.
|
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+
|
53 |
+
Testicles of a male calf or other livestock are cooked and eaten in a dish sometimes called Rocky Mountain oysters.[16]
|
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+
|
55 |
+
As early as 330 BC, Aristotle prescribed the ligation (tying off) of the left testicle in men wishing to have boys.[17] In the Middle Ages, men who wanted a boy sometimes had their left testicle removed. This was because people believed that the right testicle made "boy" sperm and the left made "girl" sperm.[18]
|
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|
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+
One theory about the etymology of the word testis is based on Roman law. The original Latin word testis, "witness", was used in the firmly established legal principle "Testis unus, testis nullus" (one witness [equals] no witness), meaning that testimony by any one person in court was to be disregarded unless corroborated by the testimony of at least another. This led to the common practice of producing two witnesses, bribed to testify the same way in cases of lawsuits with ulterior motives. Since such "witnesses" always came in pairs, the meaning was accordingly extended, often in the diminutive (testiculus, testiculi).[citation needed]
|
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+
|
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+
Another theory says that testis is influenced by a loan translation, from Greek parastatēs "defender (in law), supporter" that is "two glands side by side".[19]
|
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+
|
61 |
+
In sharks, the testicle on the right side is usually larger, and in many bird and mammal species, the left may be the larger. The primitive jawless fish have only a single testis, located in the midline of the body, although even this forms from the fusion of paired structures in the embryo.[20]
|
62 |
+
|
63 |
+
In seasonal breeders, the weight of the testes often increases during the breeding season.[21] The testicles of a dromedary camel are 7–10 cm (2.8–3.9 in) long, 4.5 cm (1.8 in) deep and 5 cm (2.0 in) in width. The right testicle is often smaller than the left.[22] The testicles of a male red fox attain their greatest weight in December–February.[23] Spermatogenesis in male golden jackals occurs 10–12 days before the females enter estrus and, during this time, males' testicles triple in weight.[24]
|
64 |
+
|
65 |
+
The basal condition for mammals is to have internal testes.[25] The testes of monotremes,[26][27] xenarthrans,[27] and elephants[28] remain within the abdomen. There are also some marsupials with external testes[29][30][31] and Boreoeutherian mammals with internal testes, such as the rhinoceros.[32] Cetaceans such as whales and dolphins also have internal testes.[33][34] As external testes would increase drag in the water they have internal testes which are kept cool by special circulatory systems that cool the arterial blood going to the testes by placing the arteries near veins bringing cooled venous blood from the skin.[35][36] In odobenids and phocids, the location of the testes is para-abdominal, though otariids have scrotal testes.[37]
|
66 |
+
|
67 |
+
Boreoeutherian land mammals, the large group of mammals that includes humans, have externalized testes.[38] Their testes function best at temperatures lower than their core body temperature. Their testes are located outside of the body, suspended by the spermatic cord within the scrotum.
|
68 |
+
|
69 |
+
There are several hypotheses why most boreotherian mammals have external testes which operate best at a temperature that is slightly less than the core body temperature, e.g. that it is stuck with enzymes evolved in a colder temperature due to external testes evolving for different reasons, that the lower temperature of the testes simply is more efficient for sperm production.
|
70 |
+
|
71 |
+
1) More efficient. The classic hypothesis is that cooler temperature of the testes allows for more efficient fertile spermatogenesis. In other words, there are no possible enzymes operating at normal core body temperature that are as efficient as the ones evolved, at least none appearing in our evolution so far.
|
72 |
+
|
73 |
+
The early mammals had lower body temperatures and thus their testes worked efficiently within their body. However it is argued that boreotherian mammals have higher body temperatures than the other mammals and had to develop external testes to keep them cool. It is argued that those mammals with internal testes, such as the monotremes, armadillos, sloths, elephants, and rhinoceroses, have a lower core body temperatures than those mammals with external testes.[citation needed]
|
74 |
+
|
75 |
+
However, the question remains why birds despite having very high core body temperatures have internal testes and did not evolve external testes.[39] It was once theorized that birds used their air sacs to cool the testes internally, but later studies revealed that birds' testes are able to function at core body temperature.[39]
|
76 |
+
|
77 |
+
Some mammals which have seasonal breeding cycles keep their testes internal until the breeding season at which point their testes descend and increase in size and become external.[40]
|
78 |
+
|
79 |
+
2) Irreversible adaptation to sperm competition. It has been suggested that the ancestor of the boreoeutherian mammals was a small mammal that required very large testes (perhaps rather like those of a hamster) for sperm competition and thus had to place its testes outside the body.[41] This led to enzymes involved in spermatogenesis, spermatogenic DNA polymerase beta and recombinase activities evolving a unique temperature optimum, slightly less than core body temperature. When the boreoeutherian mammals then diversified into forms that were larger and/or did not require intense sperm competition they still produced enzymes that operated best at cooler temperatures and had to keep their testes outside the body. This position is made less parsimonious by the fact that the kangaroo, a non-boreoeutherian mammal, has external testicles. The ancestors of kangaroos might, separately from boreotherian mammals, have also been subject to heavy sperm competition and thus developed external testes, however, kangaroo external testes are suggestive of a possible adaptive function for external testes in large animals.
|
80 |
+
|
81 |
+
3) Protection from abdominal cavity pressure changes. One argument for the evolution of external testes is that it protects the testes from abdominal cavity pressure changes caused by jumping and galloping.[42]
|
82 |
+
|
83 |
+
4) Protection against DNA damage. Mild, transient scrotal heat stress causes DNA damage, reduced fertility and abnormal embryonic development in mice.[43] DNA strand breaks were found in spermatocytes recovered from testicles subjected to 40 °C or 42 °C for 30 minutes.[43] These findings suggest that the external location of the testicles provides the adaptive benefit of protecting spermatogenic cells from heat-induced DNA damage that could otherwise lead to infertility and germline mutation.
|
84 |
+
|
85 |
+
The relative size of testes is often influenced by mating systems.[44] Testicular size as a proportion of body weight varies widely. In the mammalian kingdom, there is a tendency for testicular size to correspond with multiple mates (e.g., harems, polygamy). Production of testicular output sperm and spermatic fluid is also larger in polygamous animals, possibly a spermatogenic competition for survival. The testes of the right whale are likely to be the largest of any animal, each weighing around 500 kg (1,100 lb).[45]
|
86 |
+
|
87 |
+
Among the Hominidae, gorillas have little female promiscuity and sperm competition and the testes are small compared to body weight (0.03%). Chimpanzees have high promiscuity and large testes compared to body weight (0.3%). Human testicular size falls between these extremes (0.08%).[46]
|
88 |
+
|
89 |
+
Testis weight also varies in seasonal breeders like coyotes.[21]
|
90 |
+
|
91 |
+
Under a tough membranous shell called the tunica albuginea, the testis of amniotes, as well as some teleost fish, contains very fine coiled tubes called seminiferous tubules.
|
92 |
+
|
93 |
+
Amphibians and most fish do not possess seminiferous tubules. Instead, the sperm are produced in spherical structures called sperm ampullae. These are seasonal structures, releasing their contents during the breeding season, and then being reabsorbed by the body. Before the next breeding season, new sperm ampullae begin to form and ripen. The ampullae are otherwise essentially identical to the seminiferous tubules in higher vertebrates, including the same range of cell types.[20]
|
94 |
+
|
95 |
+
Testicle
|
96 |
+
|
97 |
+
Testicle
|
98 |
+
|
99 |
+
Testicle hanging on cremaster muscle. These are two healthy testicles. Heat causes them to descend, allowing cooling.
|
100 |
+
|
101 |
+
A healthy scrotum containing normal size testes. The scrotum is in tight condition. The image also shows the texture.
|
102 |
+
|
103 |
+
Testicle of a cat: 1: Extremitas capitata, 2: Extremitas caudata, 3: Margo epididymalis, 4: Margo liber, 5: Mesorchium, 6: Epididymis, 7: testicular artery and vene, 8: Ductus deferens
|
104 |
+
|
105 |
+
Testis surface
|
106 |
+
|
107 |
+
Testis cross section
|
108 |
+
|
109 |
+
The right testis, exposed by laying open the tunica vaginalis.
|
110 |
+
|
111 |
+
Microscopic view of Rabbit testis 100×
|
112 |
+
|
113 |
+
Testicle
|
en/5665.html.txt
ADDED
@@ -0,0 +1,113 @@
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|
1 |
+
|
2 |
+
|
3 |
+
Testicle or testis (plural testes) is the male reproductive gland or gonad in all animals, including humans. It is homologous to the female ovary. The functions of the testes are to produce both sperm and androgens, primarily testosterone. Testosterone release is controlled by the anterior pituitary luteinizing hormone; whereas sperm production is controlled both by the anterior pituitary follicle-stimulating hormone and gonadal testosterone.
|
4 |
+
|
5 |
+
Males have two testicles of similar size contained within the scrotum, which is an extension of the abdominal wall. Scrotal asymmetry is not unusual: one testicle extends farther down into the scrotum than the other due to differences in the anatomy of the vasculature.
|
6 |
+
|
7 |
+
The volume of the testicle can be estimated by palpating it and comparing it to ellipsoids of known sizes. Another method is to use calipers (an orchidometer) or a ruler either on the person or on an ultrasound image to obtain the three measurements of the x, y, and z axes (length, depth and width). These measurements can then be used to calculate the volume, using the formula for the volume of an ellipsoid:
|
8 |
+
|
9 |
+
The dimensions of the average adult testicle are up to 2 inches long, 0.8 inches in breadth, and 1.2 inches in height (5 × 2 × 3 cm). The Tanner scale for the maturity of male genitals assigns a maturity stage to the calculated volume ranging from stage I, a volume of less than 1.5 ml; to stage V, a volume greater than 20 ml. Normal volume is 15 to 25 ml; the average is 18 cm³ per testis (range 12 cm³ to 30 cm³.[1]
|
10 |
+
|
11 |
+
The testes are covered by a tough membranous shell called the tunica albuginea. Within the testes are very fine coiled tubes called seminiferous tubules. The tubules are lined with a layer of cells (germ cells) that develop from puberty through old age into sperm cells (also known as spermatozoa or male gametes). The developing sperm travel through the seminiferous tubules to the rete testis located in the mediastinum testis, to the efferent ducts, and then to the epididymis where newly created sperm cells mature (see spermatogenesis). The sperm move into the vas deferens, and are eventually expelled through the urethra and out of the urethral orifice through muscular contractions.
|
12 |
+
|
13 |
+
Blood supply and lymphatic drainage of the testes and scrotum are distinct:
|
14 |
+
|
15 |
+
Many anatomical features of the adult testis reflect its developmental origin in the abdomen. The layers of tissue enclosing each testicle are derived from the layers of the anterior abdominal wall. Notably, the cremasteric muscle arises from the internal oblique muscle.
|
16 |
+
|
17 |
+
Large molecules cannot pass from the blood into the lumen of a seminiferous tubule due to the presence of tight junctions between adjacent Sertoli cells. The spermatogonia are in the basal compartment (deep to the level of the tight junctions) and the more mature forms such as primary and secondary spermatocytes and spermatids are in the adluminal compartment.
|
18 |
+
|
19 |
+
The function of the blood–testis barrier may be to prevent an auto-immune reaction. Mature sperm (and their antigens) arise long after immune tolerance is established in infancy. Therefore, since sperm are antigenically different from self tissue, a male animal can react immunologically to his own sperm. In fact, he is capable of making antibodies against them.
|
20 |
+
|
21 |
+
Injection of sperm antigens causes inflammation of the testis (auto-immune orchitis) and reduced fertility. Thus, the blood–testis barrier may reduce the likelihood that sperm proteins will induce an immune response, reducing fertility and so progeny.
|
22 |
+
|
23 |
+
Spermatogenesis is enhanced at temperatures slightly less than core body temperature.[5] The spermatogenesis is less efficient at lower and higher temperatures than 33 °C.[5] Because the testes are located outside the body, the smooth tissue of the scrotum can move them closer or further away from the body.[5] The temperature of the testes is maintained at 35 degrees Celsius (95 degrees Fahrenheit), i.e. two degrees below the body temperature of 37 degrees Celsius (98.6 degrees Fahrenheit). Higher temperatures affect spermatogenesis.[6] There are a number of mechanisms to maintain the testes at the optimum temperature.[5]
|
24 |
+
|
25 |
+
The cremasteric muscle is part of the spermatic cord. When this muscle contracts, the cord is shortened and the testicle is moved closer up toward the body, which provides slightly more warmth to maintain optimal testicular temperature. When cooling is required, the cremasteric muscle relaxes and the testicle is lowered away from the warm body and is able to cool. Contraction also occurs in response to stress (the testicles rise up toward the body in an effort to protect them in a fight).
|
26 |
+
|
27 |
+
The cremaster muscle can reflexively raise each testicle individually if properly triggered. This phenomenon is known as the cremasteric reflex. The testicles can also be lifted voluntarily using the pubococcygeus muscle, which partially activates related muscles.
|
28 |
+
|
29 |
+
The human genome includes approximately 20,000 protein coding genes: 80% of these genes are expressed in adult testes.[7] The testes have the highest fraction of tissue type-specific genes compared to other organs and tissues:[8] about 1000 of them are highly specific for the testes,[7] and about 2,200 show an elevated pattern of expression here. A majority of these genes encode for proteins that are expressed in the seminiferous tubules and have functions related to spermatogenesis.[9][8] Sperm cells express proteins that result in the development of flagella; these same proteins are expressed in the female in cells lining the fallopian tube, and cause the development of cilia. In other words, sperm cell flagella and Fallopian tube cilia are homologous structures. The testis-specific proteins that show the highest level of expression are protamines.
|
30 |
+
|
31 |
+
There are two phases in which the testes grow substantially; namely in embryonic and pubertal age.
|
32 |
+
|
33 |
+
During mammalian development, the gonads are at first capable of becoming either ovaries or testes.[10] In humans, starting at about week 4 the gonadal rudiments are present within the intermediate mesoderm adjacent to the developing kidneys. At about week 6, sex cords develop within the forming testes. These are made up of early Sertoli cells that surround and nurture the germ cells that migrate into the gonads shortly before sex determination begins. In males, the sex-specific gene SRY that is found on the Y-chromosome initiates sex determination by downstream regulation of sex-determining factors, (such as GATA4, SOX9 and AMH), which leads to development of the male phenotype, including directing development of the early bipotential gonad down the male path of development.
|
34 |
+
|
35 |
+
Testes follow the "path of descent" from high in the posterior fetal abdomen to the inguinal ring and beyond to the inguinal canal and into the scrotum. In most cases (97% full-term, 70% preterm), both testes have descended by birth. In most other cases, only one testis fails to descend (cryptorchidism) and that will probably express itself within a year.
|
36 |
+
|
37 |
+
The testes grow in response to the start of spermatogenesis. Size depends on lytic function, sperm production (amount of spermatogenesis present in testis), interstitial fluid, and Sertoli cell fluid production. After puberty, the volume of the testes can be increased by over 500% as compared to the pre-pubertal size.[citation needed] Testicles are fully descended before one reaches puberty.
|
38 |
+
|
39 |
+
|
40 |
+
|
41 |
+
Testicular enlargement is an unspecific sign of various testicular diseases, and can be defined as a testicular size of more than 5 cm (long axis) x 3 cm (short axis).[14]
|
42 |
+
|
43 |
+
Blue balls is a slang term for a temporary fluid congestion in the testicles and prostate region caused by prolonged sexual arousal.
|
44 |
+
|
45 |
+
Testicular prostheses are available to mimic the appearance and feel of one or both testicles, when absent as from injury or as treatment in association to gender dysphoria. There have also been some instances of their implantation in dogs.[15]
|
46 |
+
|
47 |
+
To some extent, it is possible to change testicular size. Short of direct injury or subjecting them to adverse conditions, e.g., higher temperature than they are normally accustomed to, they can be shrunk by competing against their intrinsic hormonal function through the use of externally administered steroidal hormones. Steroids taken for muscle enhancement (especially anabolic steroids) often have the undesired side effect of testicular shrinkage.
|
48 |
+
|
49 |
+
Similarly, stimulation of testicular functions via gonadotropic-like hormones may enlarge their size. Testes may shrink or atrophy during hormone replacement therapy or through chemical castration.
|
50 |
+
|
51 |
+
In all cases, the loss in testes volume corresponds with a loss of spermatogenesis.
|
52 |
+
|
53 |
+
Testicles of a male calf or other livestock are cooked and eaten in a dish sometimes called Rocky Mountain oysters.[16]
|
54 |
+
|
55 |
+
As early as 330 BC, Aristotle prescribed the ligation (tying off) of the left testicle in men wishing to have boys.[17] In the Middle Ages, men who wanted a boy sometimes had their left testicle removed. This was because people believed that the right testicle made "boy" sperm and the left made "girl" sperm.[18]
|
56 |
+
|
57 |
+
One theory about the etymology of the word testis is based on Roman law. The original Latin word testis, "witness", was used in the firmly established legal principle "Testis unus, testis nullus" (one witness [equals] no witness), meaning that testimony by any one person in court was to be disregarded unless corroborated by the testimony of at least another. This led to the common practice of producing two witnesses, bribed to testify the same way in cases of lawsuits with ulterior motives. Since such "witnesses" always came in pairs, the meaning was accordingly extended, often in the diminutive (testiculus, testiculi).[citation needed]
|
58 |
+
|
59 |
+
Another theory says that testis is influenced by a loan translation, from Greek parastatēs "defender (in law), supporter" that is "two glands side by side".[19]
|
60 |
+
|
61 |
+
In sharks, the testicle on the right side is usually larger, and in many bird and mammal species, the left may be the larger. The primitive jawless fish have only a single testis, located in the midline of the body, although even this forms from the fusion of paired structures in the embryo.[20]
|
62 |
+
|
63 |
+
In seasonal breeders, the weight of the testes often increases during the breeding season.[21] The testicles of a dromedary camel are 7–10 cm (2.8–3.9 in) long, 4.5 cm (1.8 in) deep and 5 cm (2.0 in) in width. The right testicle is often smaller than the left.[22] The testicles of a male red fox attain their greatest weight in December–February.[23] Spermatogenesis in male golden jackals occurs 10–12 days before the females enter estrus and, during this time, males' testicles triple in weight.[24]
|
64 |
+
|
65 |
+
The basal condition for mammals is to have internal testes.[25] The testes of monotremes,[26][27] xenarthrans,[27] and elephants[28] remain within the abdomen. There are also some marsupials with external testes[29][30][31] and Boreoeutherian mammals with internal testes, such as the rhinoceros.[32] Cetaceans such as whales and dolphins also have internal testes.[33][34] As external testes would increase drag in the water they have internal testes which are kept cool by special circulatory systems that cool the arterial blood going to the testes by placing the arteries near veins bringing cooled venous blood from the skin.[35][36] In odobenids and phocids, the location of the testes is para-abdominal, though otariids have scrotal testes.[37]
|
66 |
+
|
67 |
+
Boreoeutherian land mammals, the large group of mammals that includes humans, have externalized testes.[38] Their testes function best at temperatures lower than their core body temperature. Their testes are located outside of the body, suspended by the spermatic cord within the scrotum.
|
68 |
+
|
69 |
+
There are several hypotheses why most boreotherian mammals have external testes which operate best at a temperature that is slightly less than the core body temperature, e.g. that it is stuck with enzymes evolved in a colder temperature due to external testes evolving for different reasons, that the lower temperature of the testes simply is more efficient for sperm production.
|
70 |
+
|
71 |
+
1) More efficient. The classic hypothesis is that cooler temperature of the testes allows for more efficient fertile spermatogenesis. In other words, there are no possible enzymes operating at normal core body temperature that are as efficient as the ones evolved, at least none appearing in our evolution so far.
|
72 |
+
|
73 |
+
The early mammals had lower body temperatures and thus their testes worked efficiently within their body. However it is argued that boreotherian mammals have higher body temperatures than the other mammals and had to develop external testes to keep them cool. It is argued that those mammals with internal testes, such as the monotremes, armadillos, sloths, elephants, and rhinoceroses, have a lower core body temperatures than those mammals with external testes.[citation needed]
|
74 |
+
|
75 |
+
However, the question remains why birds despite having very high core body temperatures have internal testes and did not evolve external testes.[39] It was once theorized that birds used their air sacs to cool the testes internally, but later studies revealed that birds' testes are able to function at core body temperature.[39]
|
76 |
+
|
77 |
+
Some mammals which have seasonal breeding cycles keep their testes internal until the breeding season at which point their testes descend and increase in size and become external.[40]
|
78 |
+
|
79 |
+
2) Irreversible adaptation to sperm competition. It has been suggested that the ancestor of the boreoeutherian mammals was a small mammal that required very large testes (perhaps rather like those of a hamster) for sperm competition and thus had to place its testes outside the body.[41] This led to enzymes involved in spermatogenesis, spermatogenic DNA polymerase beta and recombinase activities evolving a unique temperature optimum, slightly less than core body temperature. When the boreoeutherian mammals then diversified into forms that were larger and/or did not require intense sperm competition they still produced enzymes that operated best at cooler temperatures and had to keep their testes outside the body. This position is made less parsimonious by the fact that the kangaroo, a non-boreoeutherian mammal, has external testicles. The ancestors of kangaroos might, separately from boreotherian mammals, have also been subject to heavy sperm competition and thus developed external testes, however, kangaroo external testes are suggestive of a possible adaptive function for external testes in large animals.
|
80 |
+
|
81 |
+
3) Protection from abdominal cavity pressure changes. One argument for the evolution of external testes is that it protects the testes from abdominal cavity pressure changes caused by jumping and galloping.[42]
|
82 |
+
|
83 |
+
4) Protection against DNA damage. Mild, transient scrotal heat stress causes DNA damage, reduced fertility and abnormal embryonic development in mice.[43] DNA strand breaks were found in spermatocytes recovered from testicles subjected to 40 °C or 42 °C for 30 minutes.[43] These findings suggest that the external location of the testicles provides the adaptive benefit of protecting spermatogenic cells from heat-induced DNA damage that could otherwise lead to infertility and germline mutation.
|
84 |
+
|
85 |
+
The relative size of testes is often influenced by mating systems.[44] Testicular size as a proportion of body weight varies widely. In the mammalian kingdom, there is a tendency for testicular size to correspond with multiple mates (e.g., harems, polygamy). Production of testicular output sperm and spermatic fluid is also larger in polygamous animals, possibly a spermatogenic competition for survival. The testes of the right whale are likely to be the largest of any animal, each weighing around 500 kg (1,100 lb).[45]
|
86 |
+
|
87 |
+
Among the Hominidae, gorillas have little female promiscuity and sperm competition and the testes are small compared to body weight (0.03%). Chimpanzees have high promiscuity and large testes compared to body weight (0.3%). Human testicular size falls between these extremes (0.08%).[46]
|
88 |
+
|
89 |
+
Testis weight also varies in seasonal breeders like coyotes.[21]
|
90 |
+
|
91 |
+
Under a tough membranous shell called the tunica albuginea, the testis of amniotes, as well as some teleost fish, contains very fine coiled tubes called seminiferous tubules.
|
92 |
+
|
93 |
+
Amphibians and most fish do not possess seminiferous tubules. Instead, the sperm are produced in spherical structures called sperm ampullae. These are seasonal structures, releasing their contents during the breeding season, and then being reabsorbed by the body. Before the next breeding season, new sperm ampullae begin to form and ripen. The ampullae are otherwise essentially identical to the seminiferous tubules in higher vertebrates, including the same range of cell types.[20]
|
94 |
+
|
95 |
+
Testicle
|
96 |
+
|
97 |
+
Testicle
|
98 |
+
|
99 |
+
Testicle hanging on cremaster muscle. These are two healthy testicles. Heat causes them to descend, allowing cooling.
|
100 |
+
|
101 |
+
A healthy scrotum containing normal size testes. The scrotum is in tight condition. The image also shows the texture.
|
102 |
+
|
103 |
+
Testicle of a cat: 1: Extremitas capitata, 2: Extremitas caudata, 3: Margo epididymalis, 4: Margo liber, 5: Mesorchium, 6: Epididymis, 7: testicular artery and vene, 8: Ductus deferens
|
104 |
+
|
105 |
+
Testis surface
|
106 |
+
|
107 |
+
Testis cross section
|
108 |
+
|
109 |
+
The right testis, exposed by laying open the tunica vaginalis.
|
110 |
+
|
111 |
+
Microscopic view of Rabbit testis 100×
|
112 |
+
|
113 |
+
Testicle
|
en/5666.html.txt
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1 |
+
|
2 |
+
|
3 |
+
Cryptodira
|
4 |
+
Pleurodira
|
5 |
+
†Paracryptodira
|
6 |
+
|
7 |
+
Turtles are reptiles of the order Testudines characterized by a special bony or cartilaginous shell developed from their ribs and acting as a shield.[3] "Turtle" may refer to the order as a whole (American English) or to fresh-water and sea-dwelling testudines (British English).[4] The order Testudines includes both extant (living) and extinct species. The earliest known members of this group date from the Middle Jurassic,[1] making turtles one of the oldest reptile groups and a more ancient group than snakes or crocodilians. Of the 356 known species[2] alive today, some are highly endangered.[2]
|
8 |
+
|
9 |
+
Turtles are ectotherms—animals commonly called cold-blooded—meaning that their internal temperature varies according to the ambient environment. However, because of their high metabolic rate, leatherback sea turtles have a body temperature that is noticeably higher than that of the surrounding water. Turtles are classified as amniotes, along with other reptiles, birds, and mammals. Like other amniotes, turtles breathe air and do not lay eggs underwater, although many species live in or around water.
|
10 |
+
|
11 |
+
Differences exist in usage of the common terms turtle, tortoise, and terrapin, depending on the variety of English being used.[5] These terms are common names and do not reflect precise biological or taxonomic distinctions.[6]
|
12 |
+
|
13 |
+
Turtle may either refer to the order as a whole, or to particular turtles that make up a form taxon that is not monophyletic, or may be limited to only aquatic species. Tortoise usually refers to any land-dwelling, non-swimming chelonian.[7] Terrapin is used to describe several species of small, edible, hard-shell turtles, typically those found in brackish waters.
|
14 |
+
|
15 |
+
In North America, all chelonians are commonly called turtles. Tortoise is used only in reference to fully terrestrial turtles or, more narrowly, only those members of Testudinidae, the family of modern land tortoises.[8][7] Terrapin may refer to small semi-aquatic turtles that live in fresh and brackish water, in particular the diamondback terrapin (Malaclemys terrapin).[9][10][11][12] Although the members of the genus Terrapene dwell mostly on land, they are referred to as box turtles rather than tortoises.[6] The American Society of Ichthyologists and Herpetologists uses "turtle" to describe all species of the order Testudines, regardless of whether they are land-dwelling or sea-dwelling, and uses "tortoise" as a more specific term for slow-moving terrestrial species.[5]
|
16 |
+
|
17 |
+
In the United Kingdom, the word turtle is used for water-dwelling species, including ones known in the US as terrapins, but not for terrestrial species, which are known only as tortoises.
|
18 |
+
|
19 |
+
The word chelonian is popular among veterinarians, scientists, and conservationists working with these animals as a catch-all name for any member of the superorder Chelonia, which includes all turtles living and extinct, as well as their immediate ancestors. Chelonia is based on the Greek word for turtles, χελώνη chelone; Greek χέλυς chelys "tortoise" is also used in the formation of scientific names of chelonians.[13] Testudines, on the other hand, is based on the Latin word for tortoise, testudo.[14] Terrapin comes from an Algonquian word for turtle.[8][15]
|
20 |
+
|
21 |
+
Some languages do not have this distinction, as all of these are referred to by the same name. For example, in Spanish, the word tortuga is used for turtles, tortoises, and terrapins. A sea-dwelling turtle is tortuga marina, a freshwater species tortuga de río, and a tortoise tortuga terrestre.[16]
|
22 |
+
|
23 |
+
The largest living chelonian is the leatherback sea turtle (Dermochelys coriacea), which reaches a shell length of 200 cm (6.6 ft) and can reach a weight of over 900 kg (2,000 lb). Freshwater turtles are generally smaller, but with the largest species, the Asian softshell turtle Pelochelys cantorii, a few individuals have been reported up to 200 cm (6.6 ft). This dwarfs even the better-known alligator snapping turtle, the largest chelonian in North America, which attains a shell length of up to 80 cm (2.6 ft) and weighs as much as 113.4 kg (250 lb).[17]
|
24 |
+
|
25 |
+
Giant tortoises of the genera Geochelone, Meiolania, and others were relatively widely distributed around the world into prehistoric times, and are known to have existed in North and South America, Australia, and Africa. They became extinct at the same time as the appearance of man, and it is assumed humans hunted them for food. The only surviving giant tortoises are on the Seychelles and Galápagos Islands and can grow to over 130 cm (51 in) in length, and weigh about 300 kg (660 lb).[18]
|
26 |
+
|
27 |
+
The largest ever chelonian was Archelon ischyros, a Late Cretaceous sea turtle known to have been up to 4.6 m (15 ft) long.[19]
|
28 |
+
|
29 |
+
The smallest turtle is the speckled padloper tortoise of South Africa. It measures no more than 8 cm (3.1 in) in length and weighs about 140 g (4.9 oz). Two other species of small turtles are the American mud turtles and musk turtles that live in an area that ranges from Canada to South America. The shell length of many species in this group is less than 13 cm (5.1 in) in length.
|
30 |
+
|
31 |
+
Turtles are divided into two groups, according to how they retract their necks into their shells (something the ancestral Proganochelys could not do). The mechanism of neck retraction differs phylogenetically: the suborder Pleurodira retracts laterally to the side, anterior to shoulder girdles, while the suborder Cryptodira retracts straight back, between shoulder girdles.[20] These motions are largely due to the morphology and arrangement of cervical vertebrae. Of all recent turtles, the cervical column consists of nine joints and eight vertebrae, which are individually independent.[21] Since these vertebrae are not fused and are rounded, the neck is more flexible, being able to bend in the backwards and sideways directions.[20] The primary function and evolutionary implication of neck retraction is thought to be for feeding rather than protection.[22] Neck retraction and reciprocal extension allows the turtle to reach out further to capture prey while swimming. Neck expansion creates suction when the head is thrust forward and the oropharynx is expanded, and this morphology suggests the retraction function is for feeding purposes as the suction helps catch prey.[22] The protection the shell provides the head when it is retracted is therefore not the main function of retraction, thus is an exaptation.[23] As for the difference between the two methods of retraction, both Pleurodirans and Cryptodirans use the quick extension of the neck as a method of predation, so the difference in retraction mechanism is not due to a difference in ecological niche.[24]
|
32 |
+
|
33 |
+
Most turtles that spend most of their lives on land have their eyes looking down at objects in front of them. Some aquatic turtles, such as snapping turtles and soft-shelled turtles, have eyes closer to the top of the head. These species of turtle can hide from predators in shallow water, where they lie entirely submerged except for their eyes and nostrils. Near their eyes, sea turtles possess glands that produce salty tears that rid their body of excess salt taken in from the water they drink.
|
34 |
+
|
35 |
+
Turtles have rigid beaks and use their jaws to cut and chew food. Instead of having teeth, which they appear to have lost about 150–200 million years ago,[25] the upper and lower jaws of the turtle are covered by horny ridges. Carnivorous turtles usually have knife-sharp ridges for slicing through their prey. Herbivorous turtles have serrated-edged ridges that help them cut through tough plants. They use their tongues to swallow food, but unlike most reptiles, they cannot stick out their tongues to catch food.
|
36 |
+
|
37 |
+
The upper shell of the turtle is called the carapace. The lower shell that encases the belly is called the plastron. The carapace and plastron are joined together on the turtle's sides by bony structures called bridges. The inner layer of a turtle's shell is made up of about 60 bones that include portions of the backbone and the ribs, meaning the turtle cannot crawl out of its shell. In most turtles, the outer layer of the shell is covered by horny scales called scutes that are part of its outer skin, or epidermis. Scutes are made up of the fibrous protein keratin that also makes up the scales of other reptiles. These scutes overlap the seams between the shell bones and add strength to the shell. Some turtles do not have horny scutes; for example, the leatherback sea turtle and the soft-shelled turtles have shells covered with leathery skin instead.
|
38 |
+
|
39 |
+
The shape of the shell gives helpful clues about how a turtle lives. Most tortoises have a large, dome-shaped shell that makes it difficult for predators to crush the shell between their jaws. One of the few exceptions is the African pancake tortoise, which has a flat, flexible shell that allows it to hide in rock crevices. Most aquatic turtles have flat, streamlined shells, which aid in swimming and diving. American snapping turtles and musk turtles have small, cross-shaped plastrons that give them more efficient leg movement for walking along the bottom of ponds and streams. Another exception is the Belawan Turtle (Cirebon, West Java), which has sunken-back soft-shell.
|
40 |
+
|
41 |
+
The color of a turtle's shell may vary. Shells are commonly colored brown, black, or olive green. In some species, shells may have red, orange, yellow, or grey markings, often spots, lines, or irregular blotches. One of the most colorful turtles is the eastern painted turtle, which includes a yellow plastron and a black or olive shell with red markings around the rim.
|
42 |
+
|
43 |
+
Tortoises, being land-based, have rather heavy shells. In contrast, aquatic and soft-shelled turtles have lighter shells that help them avoid sinking in water and swim faster with more agility. These lighter shells have large spaces called fontanelles between the shell bones. The shells of leatherback sea turtles are extremely light because they lack scutes and contain many fontanelles.
|
44 |
+
|
45 |
+
It has been suggested by Jackson (2002) that the turtle shell can function as pH buffer. To endure through anoxic conditions, such as winter periods trapped beneath ice or within anoxic mud at the bottom of ponds, turtles utilize two general physiological mechanisms. In the case of prolonged periods of anoxia, it has been shown that the turtle shell both releases carbonate buffers and uptakes lactic acid.[26]
|
46 |
+
|
47 |
+
Respiration, for many amniotes, is achieved by the contraction and relaxation of specific muscle groups (i.e. intercostals, abdominal muscles, and/or a diaphragm) attached to an internal rib-cage that can expand or contract the body wall thus assisting airflow in and out of the lungs.[27] The ribs of Chelonians, however, are fused with their carapace and external to their pelvic and pectoral girdles, a feature unique among turtles. This rigid shell is not capable of expansion, and by rendering their rib-cage immobile, Testudines have had to evolve special adaptations for respiration.[28][29]
|
48 |
+
|
49 |
+
Turtle pulmonary ventilation occurs by using specific groups of abdominal muscles attached to their viscera and shell that pull the lungs ventrally during inspiration, where air is drawn in via a negative pressure gradient (Boyle's Law).[27] In expiration, the contraction of the transversus abdominis is the driving force by propelling the viscera into the lungs and expelling air under positive pressure.[28] Conversely, the relaxing and flattening of the oblique abdominis muscle pulls the transversus back down which, once again, draws air back into the lungs.[28] Important auxiliary muscles used for ventilatory processes are the pectoralis, which is used in conjunction with the transverse abdominis during inspiration, and the serratus, which moves with the abdominal oblique accompanying expiration.
|
50 |
+
|
51 |
+
The lungs of Testudines are multi-chambered and attached their entire length down the carapace. The number of chambers can vary between taxa, though most commonly they have three lateral chambers, three medial chambers, and one terminal chamber.[30] As previously mentioned, the act of specific abdominal muscles pulling down the viscera (or pushing back up) is what allows for respiration in turtles. Specifically, it is the turtles large liver that pulls or pushes on the lungs.[28] Ventral to the lungs, in the coelomic cavity, the liver of turtles is attached directly to the right lung, and their stomach is directly attached to the left lung by the ventral mesopneumonium, which is attached to their liver by the ventral mesentery.[28] When the liver is pulled down, inspiration begins. Supporting the lungs is the post-pulmonary septum, which is found in all Testudines, and is thought to prevent the lungs from collapsing.[31]
|
52 |
+
|
53 |
+
As mentioned above, the outer layer of the shell is part of the skin; each scute (or plate) on the shell corresponds to a single modified scale. The remainder of the skin has much smaller scales, similar to the skin of other reptiles. Turtles do not molt their skins all at once as snakes do, but continuously in small pieces. When turtles are kept in aquaria, small sheets of dead skin can be seen in the water (often appearing to be a thin piece of plastic) having been sloughed off when the animals deliberately rub themselves against a piece of wood or stone. Tortoises also shed skin, but dead skin is allowed to accumulate into thick knobs and plates that provide protection to parts of the body outside the shell.
|
54 |
+
|
55 |
+
By counting the rings formed by the stack of smaller, older scutes on top of the larger, newer ones, it is possible to estimate the age of a turtle, if one knows how many scutes are produced in a year.[32] This method is not very accurate, partly because growth rate is not constant, but also because some of the scutes eventually fall away from the shell.
|
56 |
+
|
57 |
+
Terrestrial tortoises have short, sturdy feet. Tortoises are famous for moving slowly, in part because of their heavy, cumbersome shells, which restrict stride length.
|
58 |
+
|
59 |
+
Amphibious turtles normally have limbs similar to those of tortoises, except that the feet are webbed and often have long claws. These turtles swim using all four feet in a way similar to the dog paddle, with the feet on the left and right side of the body alternately providing thrust. Large turtles tend to swim less than smaller ones, and the very big species, such as alligator snapping turtles, hardly swim at all, preferring to walk along the bottom of the river or lake. As well as webbed feet, turtles have very long claws, used to help them clamber onto riverbanks and floating logs upon which they bask. Male turtles tend to have particularly long claws, and these appear to be used to stimulate the female while mating. While most turtles have webbed feet, some, such as the pig-nosed turtle, have true flippers, with the digits being fused into paddles and the claws being relatively small. These species swim in the same way as sea turtles do (see below).
|
60 |
+
|
61 |
+
Sea turtles are almost entirely aquatic and have flippers instead of feet. Sea turtles fly through the water, using the up-and-down motion of the front flippers to generate thrust; the back feet are not used for propulsion but may be used as rudders for steering. Compared with freshwater turtles, sea turtles have very limited mobility on land, and apart from the dash from the nest to the sea as hatchlings, male sea turtles normally never leave the sea. Females must come back onto land to lay eggs. They move very slowly and laboriously, dragging themselves forwards with their flippers.
|
62 |
+
|
63 |
+
Turtles are thought to have exceptional night vision due to the unusually large number of rod cells in their retinas. Turtles have color vision with a wealth of cone subtypes with sensitivities ranging from the near ultraviolet (UVA) to red. Some land turtles have very poor pursuit movement abilities, which are normally found only in predators that hunt quick-moving prey, but carnivorous turtles are able to move their heads quickly to snap.
|
64 |
+
|
65 |
+
While typically thought of as mute, turtles make various sounds when communicating. Tortoises may be vocal when courting and mating. Various species of both freshwater and sea turtles emit numerous types of calls, often short and low frequency, from the time they are in the egg to when they are adults. These vocalizations may serve to create group cohesion when migrating.[33]
|
66 |
+
|
67 |
+
It has been reported that wood turtles are better than white rats at learning to navigate mazes.[34] Case studies exist of turtles playing.[34] They do, however, have a very low encephalization quotient (relative brain to body mass), and their hard shells enable them to live without fast reflexes or elaborate predator avoidance strategies.[35]
|
68 |
+
In the laboratory, turtles (Pseudemys nelsoni) can learn novel operant tasks and have demonstrated a long-term memory of at least 7.5 months.[36] Similarly, giant tortoises can learn and remember tasks, and master lessons much faster when trained in groups.[37] Remarkably, tortoises that were tested 9 years after the initial training still retained the operant conditioning.[38]
|
69 |
+
|
70 |
+
Turtles are known for displaying a wide variety of mating behaviors, however, they are not known for forming pair-bonds or for being part of a social group.[39] Once fertilization has occurred and an offspring has been produced, neither parent will provide care for the offspring once it's hatched.[39] Females generally outnumber males in various turtle species (such as Green turtles), and as a result, most males will engage in multiple copulation with multiple partners throughout their lifespan.[40] However, due to the sexual dimorphism present in most turtle species, males must develop different courting strategies or use alternate methods to gain access to a potential mate.[41] Most terrestrial species have males that are larger than females, and fighting between males often determines a hierarchical order in which the higher up the order an individual is, the better the chance is of the individual getting access to a potential mate.[41] For most semi-aquatic species and bottom-walking aquatic species, combat occurs less often. Males belonging to semi-aquatic and bottom-walking species instead often use their larger size advantage to forcibly mate with a female.[41] In fully aquatic species, males are often smaller than females and therefore they cannot use the same strategy as their semi-aquatic relatives, which relies on overpowering the females with strength. Males in this category resort to using courtship displays in an attempt to gain mating access to a female.[41]
|
71 |
+
|
72 |
+
Wood turtles are an example of a terrestrial species where the males have a hierarchical ranking system based on dominance through fighting, and it's shown that the males with the highest rank and thus the most wins in fights have the most offspring.[39]
|
73 |
+
|
74 |
+
Galapagos tortoises are another example of a species which has a hierarchical rank that is determined by dominance displays, and access to food and mates is regulated by this dominance hierarchy.[42] Two male saddle backs most often compete for access to cactus trees, which is their source of food.[43] The winner is the individual who stretches their neck the highest, and that individual gets access to the cactus tree, which can attract potential mates.[43]
|
75 |
+
|
76 |
+
The male scorpion mud turtle is an example of a bottom-walking aquatic species that relies on overpowering females with its larger size as a mating strategy.[41] The male approaches the female from the rear, and often resorts to aggressive methods such as biting the female's tail or hind limbs, followed by a mounting behavior in which the male clasps the edges of her carapace with his forelimbs and hind limbs to hold her in position.[44] The male follows this action by laterally waving his head and sometimes biting the female's head in an attempt to get her to withdraw her head into her shell. This exposes her cloaca, and with it exposed, the male can attempt copulation by trying to insert his grasping tail.[44]
|
77 |
+
|
78 |
+
Male radiated tortoises are also known to use the force mating strategy wherein they use surrounding vegetation to trap or prevent females from escaping, then pin them down for copulation.[45]
|
79 |
+
|
80 |
+
Red-eared sliders are an example of a fully aquatic species in which the male performs a courtship behavior. In this case the male extends his forelegs with the palms facing out and flutters his forelegs in the female's face.[46] Female choice is important in this method, and the females of some species, such as green sea turtles, aren't always receptive. As such, they've evolved certain behaviors to avoid the male's attempts at copulation, such as swimming away, confronting the male followed by biting, or a refusal position in which the female assumes a vertical position with her limbs widely outspread and her plastron facing the male.[40] If the water is too shallow to perform the refusal position, the females will resort to beaching themselves, which is a proven deterrent method, as the males will not follow them ashore.[40]
|
81 |
+
|
82 |
+
Although many turtles spend large amounts of their lives underwater, all turtles and tortoises breathe air and must surface at regular intervals to refill their lungs. Immersion periods vary between 60 seconds and 1 hour depending on the species.[47] They can also spend much or all of their lives on dry land. Aquatic respiration in Australian freshwater turtles is currently being studied. Some species have large cloacal cavities that are lined with many finger-like projections. These projections, called papillae, have a rich blood supply and increase the surface area of the cloaca. The turtles can take up dissolved oxygen from the water using these papillae, in much the same way that fish use gills to respire.[48]
|
83 |
+
|
84 |
+
Like other reptiles, turtles lay eggs that are slightly soft and leathery. The eggs of the largest species are spherical while the eggs of the rest are elongated. Their albumen is white and contains a different protein from bird eggs, such that it will not coagulate when cooked. Turtle eggs prepared to eat consist mainly of yolk. In some species, temperature determines whether an egg develops into a male or a female: a higher temperature causes a female, a lower temperature causes a male. Large numbers of eggs are deposited in holes dug into mud or sand. They are then covered and left to incubate by themselves. Depending on the species, the eggs will typically take 70–120 days to hatch.[citation needed] When the turtles hatch, they squirm their way to the surface and head toward the water. There are no known species in which the mother cares for her young.
|
85 |
+
|
86 |
+
Sea turtles lay their eggs on dry, sandy beaches. Immature sea turtles are not cared for by the adults. Turtles can take many years to reach breeding age, and in many cases, breed every few years rather than annually.
|
87 |
+
|
88 |
+
Researchers have recently discovered a turtle's organs do not gradually break down or become less efficient over time, unlike most other animals. It was found that the liver, lungs, and kidneys of a centenarian turtle are virtually indistinguishable from those of its immature counterpart. This has inspired genetic researchers to begin examining the turtle genome for longevity genes.[49]
|
89 |
+
|
90 |
+
A group of turtles is known as a bale.[50]
|
91 |
+
|
92 |
+
A turtle's diet varies greatly depending on the environment in which it lives. Adult turtles typically eat aquatic plants;[citation needed] invertebrates such as insects, snails, and worms; and have been reported to occasionally eat dead marine animals. Several small freshwater species are carnivorous, eating small fish and a wide range of aquatic life. However, protein is essential to turtle growth and juvenile turtles are purely carnivorous.
|
93 |
+
|
94 |
+
Sea turtles typically feed on jellyfish, sponges, and other soft-bodied organisms. Some species with stronger jaws have been observed to eat shellfish, while others, such as the green sea turtle, do not eat meat at all and, instead, have a diet largely made up of algae.[51]
|
95 |
+
|
96 |
+
Based on body fossils, the first proto-turtles are believed to have existed in the late Triassic Period of the Mesozoic era, about 220 million years ago, and their shell, which has remained a remarkably stable body plan, is thought to have evolved from bony extensions of their backbones and broad ribs that expanded and grew together to form a complete shell that offered protection at every stage of its evolution, even when the bony component of the shell was not complete. This is supported by fossils of the freshwater Odontochelys semitestacea or "half-shelled turtle with teeth", from the late Triassic, which have been found near Guangling in southwest China. Odontochelys displays a complete bony plastron and an incomplete carapace, similar to an early stage of turtle embryonic development.[52] Prior to this discovery, the earliest-known fossil turtle ancestors, like Proganochelys, were terrestrial and had a complete shell, offering no clue to the evolution of this remarkable anatomical feature. By the late Jurassic, turtles had radiated widely, and their fossil history becomes easier to read.
|
97 |
+
|
98 |
+
Their exact ancestry has been disputed. It was believed they are the only surviving branch of the ancient evolutionary grade Anapsida, which includes groups such as procolophonids, millerettids, protorothyrids, and pareiasaurs. All anapsid skulls lack a temporal opening while all other extant amniotes have temporal openings (although in mammals, the hole has become the zygomatic arch). The millerettids, protorothyrids, and pareiasaurs became extinct in the late Permian period and the procolophonoids during the Triassic.[53]
|
99 |
+
|
100 |
+
However, it was later suggested that the anapsid-like turtle skull may be due to reversion rather than to anapsid descent. More recent morphological phylogenetic studies with this in mind placed turtles firmly within diapsids, slightly closer to Squamata than to Archosauria.[54][55] All molecular studies have strongly upheld the placement of turtles within diapsids; some place turtles within Archosauria,[56] or, more commonly, as a sister group to extant archosaurs,[57][58][59][60] though an analysis conducted by Lyson et al. (2012) recovered turtles as the sister group of lepidosaurs instead.[61] Reanalysis of prior phylogenies suggests that they classified turtles as anapsids both because they assumed this classification (most of them studying what sort of anapsid turtles are) and because they did not sample fossil and extant taxa broadly enough for constructing the cladogram. Testudines were suggested to have diverged from other diapsids between 200 and 279 million years ago, though the debate is far from settled.[54][57][62] Even the traditional placement of turtles outside Diapsida cannot be ruled out at this point. A combined analysis of morphological and molecular data conducted by Lee (2001) found turtles to be anapsids (though a relationship with archosaurs couldn't be statistically rejected).[63] Similarly, a morphological study conducted by Lyson et al.. (2010) recovered them as anapsids most closely related to Eunotosaurus.[64] A molecular analysis of 248 nuclear genes from 16 vertebrate taxa suggests that turtles are a sister group to birds and crocodiles (the Archosauria).[65] The date of separation of turtles and birds and crocodiles was estimated to be 255 million years ago. The most recent common ancestor of living turtles, corresponding to the split between Pleurodira and Cryptodira, was estimated to have occurred around 157 million years ago.[66][67] The oldest definitive crown-group turtle (member of the modern clade Testudines) is the species Caribemys oxfordiensis from the late Jurassic period (Oxfordian stage).[66] Through utilizing the first genomic-scale phylogenetic analysis of ultraconserved elements (UCEs) to investigate the placement of turtles within reptiles, Crawford et al. (2012) also suggest that turtles are a sister group to birds and crocodiles (the Archosauria).[68]
|
101 |
+
|
102 |
+
The first genome-wide phylogenetic analysis was completed by Wang et al. (2013). Using the draft genomes of Chelonia mydas and Pelodiscus sinensis, the team used the largest turtle data set to date in their analysis and concluded that turtles are likely a sister group of crocodilians and birds (Archosauria).[69] This placement within the diapsids suggests that the turtle lineage lost diapsid skull characteristics as it now possesses an anapsid-like skull.
|
103 |
+
|
104 |
+
The earliest known fully shelled member of the turtle lineage is the late Triassic Proganochelys. This genus already possessed many advanced turtle traits, and thus probably indicates many millions of years of preceding turtle evolution; this is further supported by evidence from fossil tracks from the Early Triassic of the United States (Wyoming and Utah) and from the Middle Triassic of Germany, indicating that proto-turtles already existed as early as the Early Triassic.[70] Proganochelys lacked the ability to pull its head into its shell, had a long neck, and had a long, spiked tail ending in a club. While this body form is similar to that of ankylosaurs, it resulted from convergent evolution.
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Turtles are divided into two extant suborders: Cryptodira and Pleurodira. The Cryptodira is the larger of the two groups and includes all the marine turtles, the terrestrial tortoises, and many of the freshwater turtles. The Pleurodira are sometimes known as the side-necked turtles, a reference to the way they retract their heads into their shells. This smaller group consists primarily of various freshwater turtles. Until 3,000 years ago, the family Meiolaniidae was also extant, but this family is outside the Testudines crown group, belonging to Perichelydia.
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Order Testudines Linnaeus, 1758
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Turtle fossils of hatchling and nestling size have been documented in the scientific literature.[71] Paleontologists from North Carolina State University have found the fossilized remains of the world's largest turtle in a coal mine in Colombia. The specimen named as Carbonemys cofrinii is around 60 million years old and nearly 2.4 m (8 ft) long.[72]
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On a few rare occasions, paleontologists have succeeded in unearthing large numbers of Jurassic or Cretaceous turtle skeletons accumulated in a single area (the Nemegt Formation in Mongolia, the Turtle Graveyard in North Dakota, or the Black Mountain Turtle Layer in Wyoming). The most spectacular find of this kind to date occurred in 2009 in Shanshan County in Xinjiang, where over a thousand ancient freshwater turtles apparently died after the last water hole in an area dried out during a major drought.[73][74]
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Though absent from New Zealand in recent times, turtle fossils are known from the Miocene Saint Bathans Fauna, represented by a meiolaniid and pleurodires.[75]
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Turtles possess diverse chromosome numbers (2n = 28–66) and a myriad of chromosomal rearrangements have occurred during evolution.[76]
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Some turtles, particularly small terrestrial and freshwater turtles, are kept as pets. Among the most popular are Russian tortoises, spur-thighed tortoises, and red-eared sliders.[77]
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In the United States, due to the ease of contracting salmonellosis through casual contact with turtles, the U.S. Food and Drug Administration (FDA) established a regulation in 1975 to discontinue the sale of turtles under 4 in (100 mm).[78] It is supposed to be illegal in every US state for anyone to sell any turtles under 4 inches (10 cm) long, but many stores and flea markets still sell small turtles due to a loophole in the FDA regulation which allows turtles under 4 in (100 mm) to be sold for educational purposes.[79][80]
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Some states have other laws and regulations regarding possession of red-eared sliders as pets because they are looked upon as invasive species or pests where they are not native, but have been introduced through the pet trade. As of July 1, 2007, it is illegal in Florida to sell any wild type red-eared slider.[81] Unusual color varieties such as albino and pastel red-eared sliders, which are derived from captive breeding, are still allowed for sale in that state.
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In Europe, turtle and tortoise keeping became popular in the 1960s and 1970s, when large numbers of wild-caught turtles and tortoises were imported. This was especially devastating to the Mediterranean tortoise population. In the 1980s the import of wild-caught tortoises started to be banned in various countries. Most turtles and tortoises for sale in Europe today are captive-bred. [82]
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Turtles and tortoises are seen by some people as cheap pets that need little care. The complexity and expense of proper turtle and tortoise husbandry is often underestimated.[82] Most species of tortoise need a spacious outdoor enclosure with areas at different temperatures so they can thermoregulate. They also need opportunities to climb, dig and forage. Most species of tortoise should be fed dark, leafy greens with calcium and vitamin supplements. [83]
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Turtles require a large tub or aquarium with land areas where they can dry off completely and other areas where they can rest near the water's surface, on a piece of submerged driftwood for example. Like tortoises, turtles need access to UVB lighting and a varied diet rich in calcium. [84]
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The flesh of turtles, calipash or calipee, was and still is considered a delicacy in a number of cultures.[85] Turtle soup has been a prized dish in Anglo-American cuisine,[86] and still remains so in some parts of Asia.[which?] Gopher tortoise stew was popular with some groups in Florida.[87]
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Turtles remain a part of the traditional diet on the island of Grand Cayman, so much so that when wild stocks became depleted, a turtle farm was established specifically to raise sea turtles for their meat. The farm also releases specimens to the wild as part of an effort to repopulate the Caribbean Sea.[88]
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Fat from turtles is also used in the Caribbean and in Mexico as a main ingredient in cosmetics, marketed under its Spanish name crema de tortuga.[89]
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Turtle plastrons (the part of the shell that covers a tortoise from the bottom) are widely used in traditional Chinese medicine; according to statistics, Taiwan imports hundreds of tons of plastrons every year.[90] A popular medicinal preparation based on powdered turtle plastron (and a variety of herbs) is the guilinggao jelly;[91] these days, though, it is typically made with only herbal ingredients.
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In February 2011, the Tortoise and Freshwater Turtle Specialist Group published a report about the top 25 species of turtles most likely to become extinct, with a further 40 species at very high risk of becoming extinct. This list excludes sea turtles, however, both the leatherback and the Kemp's ridley would make the top 25 list. The report is due to be updated in four years time allowing to follow the evolution of the list. Between 48 and 54% of all 328 of their species considered threatened, turtles and tortoises are at a much higher risk of extinction than many other vertebrates. Of the 263 species of freshwater and terrestrial turtles, 117 species are considered Threatened, 73 are either Endangered or Critically Endangered and 1 is Extinct. Of the 58 species belonging to the family Testudinidae, 33 species are Threatened, 18 are either Endangered or Critically Endangered, 1 is Extinct in the wild and 7 species are Extinct. 71% of all tortoise species are either gone or almost gone. Asian species are the most endangered, closely followed by the five endemic species from Madagascar. Turtles face many threats, including habitat destruction, harvesting for consumption, and the pet trade. The high extinction risk for Asian species is primarily due to the long-term unsustainable exploitation of turtles and tortoises for consumption and traditional Chinese medicine, and to a lesser extent for the international pet trade.[92]
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|
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Efforts have been made by Chinese entrepreneurs to satisfy increasing demand for turtle meat as gourmet food and traditional medicine with farmed turtles, instead of wild-caught ones; according to a study published in 2007, over a thousand turtle farms operated in China.[93][94] Turtle farms in Oklahoma and Louisiana raise turtles for export to China.[94]
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|
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+
Nonetheless, wild turtles continue to be caught and sent to market in large number (as well as to turtle farms, to be used as breeding stock[93]), resulting in a situation described by conservationists as "the Asian turtle crisis".[95] In the words of the biologist George Amato, "the amount and the volume of captured turtles ... vacuumed up entire species from areas in Southeast Asia", even as biologists still did not know how many distinct turtle species live in the region.[96] About 75% of Asia's 90 tortoise and freshwater turtle species are estimated to have become threatened.[94]
|
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|
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+
Harvesting wild turtles is legal in a number of states in the USA.[94] In one of these states, Florida, just a single seafood company in Fort Lauderdale was reported in 2008 as buying about 5,000 pounds of softshell turtles a week. The harvesters (hunters) are paid about $2 a pound; some manage to catch as many as 30–40 turtles (500 pounds) on a good day. Some of the catch gets to the local restaurants, while most of it is exported to Asia. The Florida Fish and Wildlife Conservation Commission estimated in 2008 that around 3,000 pounds of softshell turtles were exported each week via Tampa International Airport.[97]
|
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|
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+
Nonetheless, the great majority of turtles exported from the USA are farm raised. According to one estimate by the World Chelonian Trust, about 97% of 31.8 million animals harvested in the U.S. over a three-year period (November 4, 2002 – November 26, 2005) were exported.[94][98] It has been estimated (presumably, over the same 2002–2005 period) that about 47% of the US turtle exports go to People's Republic of China (predominantly to Hong Kong), another 20% to Taiwan, and 11% to Mexico.[99][100]
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|
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+
TurtleSAt is a smartphone app that has been developed in Australia in honor of World Turtle Day to help in the conservation of fresh water turtles in Australia. The app will allow the user to identify turtles with a picture guide and the location of turtles using the phones GPS to record sightings and help find hidden turtle nesting grounds. The app has been developed because there has been a high per cent of decline of fresh water turtles in Australia due to foxes, droughts, and urban development. The aim of the app is to reduce the number of foxes and help with targeting feral animal control.[101]
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|
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Queensland's shark culling program, which has killed roughly 50,000 sharks since 1962, has also killed thousands of turtles as bycatch.[102][103] Over 5,000 marine turtles have been killed in Queensland's "shark control" program (which uses shark nets and drum lines).[103] The program has also killed 719 loggerhead turtles and 33 hawksbill turtles (hawksbill turtles are critically endangered).[103] New South Wales has a "shark control" program which has killed many turtles: its program uses shark nets,[104][105] in which more than 5,000 marine turtles have been caught.[106]
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Tortoises (/ˈtɔːr.təs.ɪz/) are reptile species of the family Testudinidae of the order Testudines (the turtles). They are particularly distinguished from other turtles by being land-dwelling, while many (though not all) other turtle species are at least partly aquatic. Like other turtles, tortoises have a shell to protect from predation and other threats. The shell in tortoises is generally hard, and like other members of the suborder Cryptodira, they retract their necks and heads directly backwards into the shell to protect them.
|
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|
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Tortoises are unique among vertebrates in that the pectoral and pelvic girdles are inside the ribcage rather than outside. Tortoises can vary in dimension from a few centimeters to two meters. They are usually diurnal animals with tendencies to be crepuscular depending on the ambient temperatures. They are generally reclusive animals. Tortoises are the longest-living land animals in the world, although the longest-living species of tortoise is a matter of debate. Galápagos tortoises are noted to live over 150 years, but an Aldabra giant tortoise named Adwaita may have lived an estimated 255 years. In general, most tortoise species can live 80–150 years.
|
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|
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Differences exist in usage of the common terms turtle, tortoise, and terrapin, depending on the variety of English being used; usage is inconsistent and contradictory.[1] These terms are common names and do not reflect precise biological or taxonomic distinctions.[2]
|
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|
9 |
+
The American Society of Ichthyologists and Herpetologists uses "turtle" to describe all species of the order Testudines, regardless of whether they are land-dwelling or sea-dwelling, and uses "tortoise" as a more specific term for slow-moving terrestrial species.[1] General American usage agrees; turtle is often a general term (although some restrict it to aquatic turtles); tortoise is used only in reference to terrestrial turtles or, more narrowly, only those members of Testudinidae, the family of modern land tortoises; and terrapin may refer to turtles that are small and live in fresh and brackish water, in particular the diamondback terrapin (Malaclemys terrapin).[3][4][5][6] In America, for example, the members of the genus Terrapene dwell on land, yet are referred to as box turtles rather than tortoises.[2]
|
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|
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+
British usage, by contrast, tends not to use "turtle" as a generic term for all members of the order, and also applies the term "tortoises" broadly to all land-dwelling members of the order Testudines, regardless of whether they are actually members of the family Testudinidae.[6] In Britain, terrapin is used to refer to a larger group of semiaquatic turtles than the restricted meaning in America.[4][7]
|
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+
|
13 |
+
Australian usage is different from both American and British usage.[6] Land tortoises are not native to Australia, yet traditionally freshwater turtles have been called "tortoises" in Australia.[8] Some Australian experts disapprove of this usage—believing that the term tortoises is "better confined to purely terrestrial animals with very different habits and needs, none of which are found in this country"—and promote the use of the term "freshwater turtle" to describe Australia's primarily aquatic members of the order Testudines because it avoids misleading use of the word "tortoise" and also is a useful distinction from marine turtles.[8]
|
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|
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Most species of tortoises lay small clutch sizes, seldom exceeding 20 eggs, and many species have clutch sizes of only 1–2 eggs. Incubation is characteristically long in most species, the average incubation period are between 100 and 160 days. Egg-laying typically occurs at night, after which the mother tortoise covers her clutch with sand, soil, and organic material. The eggs are left unattended, and depending on the species, take from 60 to 120 days to incubate.[9] The size of the egg depends on the size of the mother and can be estimated by examining the width of the cloacal opening between the carapace and plastron. The plastron of a female tortoise often has a noticeable V-shaped notch below the tail which facilitates passing the eggs. Upon completion of the incubation period, a fully formed hatchling uses an egg tooth to break out of its shell. It digs to the surface of the nest and begins a life of survival on its own. They are hatched with an embryonic egg sac which serves as a source of nutrition for the first three to seven days until they have the strength and mobility to find food. Juvenile tortoises often require a different balance of nutrients than adults, so may eat foods which a more mature tortoise would not. For example, the young of a strictly herbivorous species commonly will consume worms or insect larvae for additional protein.[citation needed]
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|
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The number of concentric rings on the carapace, much like the cross-section of a tree, can sometimes give a clue to how old the animal is, but, since the growth depends highly on the accessibility of food and water, a tortoise that has access to plenty of forage (or is regularly fed by its owner) with no seasonal variation will have no noticeable rings. Moreover, some tortoises grow more than one ring per season, and in some others, due to wear, some rings are no longer visible.[10]
|
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|
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Tortoises generally have one of the longest lifespans of any animal, and some individuals are known to have lived longer than 150 years.[11]
|
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Because of this, they symbolize longevity in some cultures, such as China. The oldest tortoise ever recorded, and one of the oldest individual animals ever recorded, was Tu'i Malila, which was presented to the Tongan royal family by the British explorer James Cook shortly after its birth in 1777. Tu'i Malila remained in the care of the Tongan royal family until its death by natural causes on May 19, 1965, at the age of 188.[12] The record for the longest-lived vertebrate is exceeded only by one other, a koi named Hanako, whose death on July 17, 1977, ended a 226-year lifespan.[13]
|
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|
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The Alipore Zoo in India was the home to Adwaita, which zoo officials claimed was the oldest living animal until its death on March 23, 2006. Adwaita (also spelled Addwaita) was an Aldabra giant tortoise brought to India by Lord Wellesley, who handed it over to the Alipur Zoological Gardens in 1875 when the zoo was set up. West Bengal officials said records showed Adwaita was at least 150 years old, but other evidence pointed to 250. Adwaita was said to be the pet of Robert Clive.[14]
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Harriet was a resident at the Australia Zoo in Queensland from 1987 to her death in 2006; she was believed to have been brought to England by Charles Darwin aboard the Beagle and then on to Australia by John Clements Wickham.[15] Harriet died on June 23, 2006, just shy of her 176th birthday.
|
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|
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Timothy, a female spur-thighed tortoise, lived to be about 165 years old. For 38 years, she was carried as a mascot aboard various ships in Britain's Royal Navy. Then in 1892, at age 53, she retired to the grounds of Powderham Castle in Devon. Up to the time of her death in 2004, she was believed to be the United Kingdom's oldest resident.[16]
|
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|
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Jonathan, a Seychelles giant tortoise living on the island of St Helena, may be as old as 188 years [17] or 184 years.[18]
|
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Many species of tortoises are sexually dimorphic, though the differences between males and females vary from species to species.[19] In some species, males have a longer, more protruding neck plate than their female counterparts, while in others, the claws are longer on the females.
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The male plastron is curved inwards to aid reproduction. The easiest way to determine the sex of a tortoise is to look at the tail. The females, as a general rule, have smaller tails, dropped down, whereas the males have much longer tails which are usually pulled up and to the side of the rear shell.
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|
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The brain of a tortoise is extremely small. Red-footed tortoises, from Central and South America, do not have an area in the brain called the hippocampus, which relates to emotion, learning, memory and spatial navigation. Studies have shown that red-footed tortoises may rely on an area of the brain called the medial cortex for emotional actions, an area that humans use for actions such as decision making.[20]
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In the 17th century, Francesco Redi performed an experiment that involved removing the brain of a land tortoise, which then proceeded to live six months. Freshwater tortoises, when subjected to the same experiment, continued similarly, but did not live so long. Redi also cut the head off a tortoise entirely, and it lived for 23 days.[21][22][23]
|
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|
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Tortoises are found from southern North America to southern South America, around the Mediterranean basin, across Eurasia to Southeast Asia, in sub-Saharan Africa, Madagascar, and some Pacific islands. They are absent from Australasia [24] They live in diverse habitats, including deserts, arid grasslands, and scrub to wet evergreen forests, and from sea level to mountains. Most species, however, occupy semiarid habitats.
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Most land-based tortoises are herbivores, feeding on grasses, weeds, leafy greens, flowers, and some fruits, although some omnivorous species are in this family. Pet tortoises typically require diets based on wild grasses, weeds, leafy greens and certain flowers. Certain species consume worms or insects and carrion in their normal habitats. Too much protein is detrimental in herbivorous species, and has been associated with shell deformities and other medical problems. Different tortoise species vary greatly in their nutritional requirements.
|
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+
This species list largely follows Turtle Taxonomy Working Group (2019)[25] and the Turtle Extinctions Working Group (2015)[26]
|
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+
|
44 |
+
Family Testudinidae Batsch 1788[27]
|
45 |
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|
46 |
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In Hinduism, Kurma (Sanskrit: कुर्म) was the second Avatar of Vishnu. Like the Matsya Avatara, Kurma also belongs to the Satya Yuga. Vishnu took the form of a half-man, half-tortoise, the lower half being a tortoise. He is normally shown as having four arms. He sat on the bottom of the ocean after the Great Flood. A mountain was placed on his back by the other gods so they could churn the sea and find the ancient treasures of the Vedic peoples.
|
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|
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In Judaism and early Christianity tortoises were seen as unclean animals .[43]
|
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|
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Tortoise shells were used by ancient Chinese as oracle bones to make predictions.
|
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|
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The tortoise is a symbol of the Ancient Greek god, Hermes.
|
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|
54 |
+
See also World Turtle and Cultural depictions of turtles.
|
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|
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Baby Testudo marginata emerges from its egg
|
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|
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Baby tortoise, less than a day old
|
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|
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+
Testudo graeca ibera, Testudo hermanni boettgeri, Testudo hermanni hermanni, Testudo marginata sarda
|
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|
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Young (3.5 years) African spurred tortoise, Geochelone sulcata
|
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|
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Young, 20-year-old Tanzanian leopard tortoise feeding on grass
|
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|
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Aldabra giant tortoise, Geochelone gigantea
|
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22-year-old leopard tortoise
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African spurred tortoise from the Oakland Zoo
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Pair of African spurred tortoises mate in a zoo
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Boy rides a tortoise at a zoo
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Video of tortoises mating
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Young Testudo hermanni
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Tortoises (/ˈtɔːr.təs.ɪz/) are reptile species of the family Testudinidae of the order Testudines (the turtles). They are particularly distinguished from other turtles by being land-dwelling, while many (though not all) other turtle species are at least partly aquatic. Like other turtles, tortoises have a shell to protect from predation and other threats. The shell in tortoises is generally hard, and like other members of the suborder Cryptodira, they retract their necks and heads directly backwards into the shell to protect them.
|
4 |
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|
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+
Tortoises are unique among vertebrates in that the pectoral and pelvic girdles are inside the ribcage rather than outside. Tortoises can vary in dimension from a few centimeters to two meters. They are usually diurnal animals with tendencies to be crepuscular depending on the ambient temperatures. They are generally reclusive animals. Tortoises are the longest-living land animals in the world, although the longest-living species of tortoise is a matter of debate. Galápagos tortoises are noted to live over 150 years, but an Aldabra giant tortoise named Adwaita may have lived an estimated 255 years. In general, most tortoise species can live 80–150 years.
|
6 |
+
|
7 |
+
Differences exist in usage of the common terms turtle, tortoise, and terrapin, depending on the variety of English being used; usage is inconsistent and contradictory.[1] These terms are common names and do not reflect precise biological or taxonomic distinctions.[2]
|
8 |
+
|
9 |
+
The American Society of Ichthyologists and Herpetologists uses "turtle" to describe all species of the order Testudines, regardless of whether they are land-dwelling or sea-dwelling, and uses "tortoise" as a more specific term for slow-moving terrestrial species.[1] General American usage agrees; turtle is often a general term (although some restrict it to aquatic turtles); tortoise is used only in reference to terrestrial turtles or, more narrowly, only those members of Testudinidae, the family of modern land tortoises; and terrapin may refer to turtles that are small and live in fresh and brackish water, in particular the diamondback terrapin (Malaclemys terrapin).[3][4][5][6] In America, for example, the members of the genus Terrapene dwell on land, yet are referred to as box turtles rather than tortoises.[2]
|
10 |
+
|
11 |
+
British usage, by contrast, tends not to use "turtle" as a generic term for all members of the order, and also applies the term "tortoises" broadly to all land-dwelling members of the order Testudines, regardless of whether they are actually members of the family Testudinidae.[6] In Britain, terrapin is used to refer to a larger group of semiaquatic turtles than the restricted meaning in America.[4][7]
|
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|
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+
Australian usage is different from both American and British usage.[6] Land tortoises are not native to Australia, yet traditionally freshwater turtles have been called "tortoises" in Australia.[8] Some Australian experts disapprove of this usage—believing that the term tortoises is "better confined to purely terrestrial animals with very different habits and needs, none of which are found in this country"—and promote the use of the term "freshwater turtle" to describe Australia's primarily aquatic members of the order Testudines because it avoids misleading use of the word "tortoise" and also is a useful distinction from marine turtles.[8]
|
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Most species of tortoises lay small clutch sizes, seldom exceeding 20 eggs, and many species have clutch sizes of only 1–2 eggs. Incubation is characteristically long in most species, the average incubation period are between 100 and 160 days. Egg-laying typically occurs at night, after which the mother tortoise covers her clutch with sand, soil, and organic material. The eggs are left unattended, and depending on the species, take from 60 to 120 days to incubate.[9] The size of the egg depends on the size of the mother and can be estimated by examining the width of the cloacal opening between the carapace and plastron. The plastron of a female tortoise often has a noticeable V-shaped notch below the tail which facilitates passing the eggs. Upon completion of the incubation period, a fully formed hatchling uses an egg tooth to break out of its shell. It digs to the surface of the nest and begins a life of survival on its own. They are hatched with an embryonic egg sac which serves as a source of nutrition for the first three to seven days until they have the strength and mobility to find food. Juvenile tortoises often require a different balance of nutrients than adults, so may eat foods which a more mature tortoise would not. For example, the young of a strictly herbivorous species commonly will consume worms or insect larvae for additional protein.[citation needed]
|
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+
|
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+
The number of concentric rings on the carapace, much like the cross-section of a tree, can sometimes give a clue to how old the animal is, but, since the growth depends highly on the accessibility of food and water, a tortoise that has access to plenty of forage (or is regularly fed by its owner) with no seasonal variation will have no noticeable rings. Moreover, some tortoises grow more than one ring per season, and in some others, due to wear, some rings are no longer visible.[10]
|
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|
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+
Tortoises generally have one of the longest lifespans of any animal, and some individuals are known to have lived longer than 150 years.[11]
|
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+
Because of this, they symbolize longevity in some cultures, such as China. The oldest tortoise ever recorded, and one of the oldest individual animals ever recorded, was Tu'i Malila, which was presented to the Tongan royal family by the British explorer James Cook shortly after its birth in 1777. Tu'i Malila remained in the care of the Tongan royal family until its death by natural causes on May 19, 1965, at the age of 188.[12] The record for the longest-lived vertebrate is exceeded only by one other, a koi named Hanako, whose death on July 17, 1977, ended a 226-year lifespan.[13]
|
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+
|
22 |
+
The Alipore Zoo in India was the home to Adwaita, which zoo officials claimed was the oldest living animal until its death on March 23, 2006. Adwaita (also spelled Addwaita) was an Aldabra giant tortoise brought to India by Lord Wellesley, who handed it over to the Alipur Zoological Gardens in 1875 when the zoo was set up. West Bengal officials said records showed Adwaita was at least 150 years old, but other evidence pointed to 250. Adwaita was said to be the pet of Robert Clive.[14]
|
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+
Harriet was a resident at the Australia Zoo in Queensland from 1987 to her death in 2006; she was believed to have been brought to England by Charles Darwin aboard the Beagle and then on to Australia by John Clements Wickham.[15] Harriet died on June 23, 2006, just shy of her 176th birthday.
|
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+
|
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+
Timothy, a female spur-thighed tortoise, lived to be about 165 years old. For 38 years, she was carried as a mascot aboard various ships in Britain's Royal Navy. Then in 1892, at age 53, she retired to the grounds of Powderham Castle in Devon. Up to the time of her death in 2004, she was believed to be the United Kingdom's oldest resident.[16]
|
27 |
+
|
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+
Jonathan, a Seychelles giant tortoise living on the island of St Helena, may be as old as 188 years [17] or 184 years.[18]
|
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+
|
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+
Many species of tortoises are sexually dimorphic, though the differences between males and females vary from species to species.[19] In some species, males have a longer, more protruding neck plate than their female counterparts, while in others, the claws are longer on the females.
|
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|
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The male plastron is curved inwards to aid reproduction. The easiest way to determine the sex of a tortoise is to look at the tail. The females, as a general rule, have smaller tails, dropped down, whereas the males have much longer tails which are usually pulled up and to the side of the rear shell.
|
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|
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+
The brain of a tortoise is extremely small. Red-footed tortoises, from Central and South America, do not have an area in the brain called the hippocampus, which relates to emotion, learning, memory and spatial navigation. Studies have shown that red-footed tortoises may rely on an area of the brain called the medial cortex for emotional actions, an area that humans use for actions such as decision making.[20]
|
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In the 17th century, Francesco Redi performed an experiment that involved removing the brain of a land tortoise, which then proceeded to live six months. Freshwater tortoises, when subjected to the same experiment, continued similarly, but did not live so long. Redi also cut the head off a tortoise entirely, and it lived for 23 days.[21][22][23]
|
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Tortoises are found from southern North America to southern South America, around the Mediterranean basin, across Eurasia to Southeast Asia, in sub-Saharan Africa, Madagascar, and some Pacific islands. They are absent from Australasia [24] They live in diverse habitats, including deserts, arid grasslands, and scrub to wet evergreen forests, and from sea level to mountains. Most species, however, occupy semiarid habitats.
|
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+
|
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Most land-based tortoises are herbivores, feeding on grasses, weeds, leafy greens, flowers, and some fruits, although some omnivorous species are in this family. Pet tortoises typically require diets based on wild grasses, weeds, leafy greens and certain flowers. Certain species consume worms or insects and carrion in their normal habitats. Too much protein is detrimental in herbivorous species, and has been associated with shell deformities and other medical problems. Different tortoise species vary greatly in their nutritional requirements.
|
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|
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+
This species list largely follows Turtle Taxonomy Working Group (2019)[25] and the Turtle Extinctions Working Group (2015)[26]
|
43 |
+
|
44 |
+
Family Testudinidae Batsch 1788[27]
|
45 |
+
|
46 |
+
In Hinduism, Kurma (Sanskrit: कुर्म) was the second Avatar of Vishnu. Like the Matsya Avatara, Kurma also belongs to the Satya Yuga. Vishnu took the form of a half-man, half-tortoise, the lower half being a tortoise. He is normally shown as having four arms. He sat on the bottom of the ocean after the Great Flood. A mountain was placed on his back by the other gods so they could churn the sea and find the ancient treasures of the Vedic peoples.
|
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|
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In Judaism and early Christianity tortoises were seen as unclean animals .[43]
|
49 |
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|
50 |
+
Tortoise shells were used by ancient Chinese as oracle bones to make predictions.
|
51 |
+
|
52 |
+
The tortoise is a symbol of the Ancient Greek god, Hermes.
|
53 |
+
|
54 |
+
See also World Turtle and Cultural depictions of turtles.
|
55 |
+
|
56 |
+
Baby Testudo marginata emerges from its egg
|
57 |
+
|
58 |
+
Baby tortoise, less than a day old
|
59 |
+
|
60 |
+
Testudo graeca ibera, Testudo hermanni boettgeri, Testudo hermanni hermanni, Testudo marginata sarda
|
61 |
+
|
62 |
+
Young (3.5 years) African spurred tortoise, Geochelone sulcata
|
63 |
+
|
64 |
+
Young, 20-year-old Tanzanian leopard tortoise feeding on grass
|
65 |
+
|
66 |
+
Aldabra giant tortoise, Geochelone gigantea
|
67 |
+
|
68 |
+
22-year-old leopard tortoise
|
69 |
+
|
70 |
+
African spurred tortoise from the Oakland Zoo
|
71 |
+
|
72 |
+
Pair of African spurred tortoises mate in a zoo
|
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|
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Boy rides a tortoise at a zoo
|
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|
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Video of tortoises mating
|
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|
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Young Testudo hermanni
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A head is the part of an organism which usually includes the ears, brain, forehead, cheeks, chin, eyes, nose, and mouth, each of which aid in various sensory functions such as sight, hearing, smell, and taste, respectively. Some very simple animals may not have a head, but many bilaterally symmetric forms do, regardless of size.
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Heads develop in animals by an evolutionary trend known as cephalization. In bilaterally symmetrical animals, nervous tissues concentrate at the anterior region, forming structures responsible for information processing. Through biological evolution, sense organs and feeding structures also concentrate into the anterior region; these collectively form the head.
|
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|
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The human head is an anatomical unit that consists of the skull, hyoid bone and cervical vertebrae. The term "skull" collectively denotes the mandible (lower jaw bone) and the cranium (upper portion of the skull that houses the brain). The skull can also be described as being composed of the cranium, which encloses the cranial cavity, and the facial skeleton (includes the mandible). There are eight bones in the cranium and fourteen in the facial skeleton.[1]
|
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|
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Sculptures of human heads are generally based on a skeletal structure that consists of a cranium, jawbone, and cheekbone. Though the number of muscles making up the face is generally consistent between sculptures, the shape of the muscles varies widely based on the function, development, and expressions reflected on the faces of the subjects.[2]
|
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Proponents of identism believe that the mind is identical to the brain. Philosopher John Searle asserts his identist beliefs, stating "the brain is the only thing in the human head". Similarly, Dr. Henry Bennet-Clark has stated that the head encloses billions of "miniagents and microagents (with no single Boss)".[3]
|
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Mid-sagittal section of a human skull, by Leonardo da Vinci, c. 1489
|
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|
15 |
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Transection of a human head, by Thomas Bartholin, 1673
|
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The evolution of a head is associated with the cephalization that occurred in Bilateria some 555 million years ago.
|
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In some arthropods, especially trilobites (pictured at right), the cephalon, or cephalic region, is the region of the head which is a collective of "fused segments".[4]
|
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A typical insect head is composed of eyes, antennae, and components of mouth. As these components differ substantially from insect to insect, they form important identification links. Eyes in the head found, in several types of insects, are in the form of a pair of compound eyes with multiple faces. In many other types of insects the compound eyes are seen in a "single facet or group of single facets". In some case, the eyes may be seen as marks on the dorsal or located near or toward the head, two or three ocelli (single faceted organs).[5]
|
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|
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Antennae on the insect's head is found in the form of segmented attachments, in pairs, that are usually located between the eyes. These are in varying shapes and sizes, in the form of filaments or in different enlarged or clubbed form.[5]
|
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+
|
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Insects have mouth parts in various shapes depending on their feeding habits. Labrum is the "upper lip" which is in the front area of the head and is the most exterior part. A pair of mandible is found on backside of the labrum flanking the side of the mouth, succeeded by a pair of maxillae each of which is known as maxilliary palp. At the back side of the mouth is the labium or lower lip. There is also an extra mouth part in some insects which is termed as hypopharynx which is usually located between the maxillac.[5]
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Though invertebrate chordates – such as the tunicate larvae or the lancelets – have heads, there has been a question of how the vertebrate head, characterized by a bony skull clearly separated from the main body, might have evolved from the head structures of these animals.[6]
|
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|
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According to Hyman (1979), the evolution of the head in the vertebrates has occurred by the fusion of a fixed number of anterior segments, in the same manner as in other "heteronomously segmented animals". In some cases, segments or a portion of the segments disappear. The head segments also lose most of its systems except for the nervous system. With the progressive development of cephalization, "the head incorporates more and more of the adjacent segments into its structure, so that in general it may be said that the higher the degree of cephalization the greater is the number of segments composing the head".[7]
|
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|
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+
In the 1980s, the "new head hypothesis" was proposed, suggesting that the vertebrate head is an evolutionary novelty resulting from the emergence of neural crest and cranial placodes.[8][9] In 2014, a transient larva tissue of the lancelet was found to be virtually indistinguishable from the neural crest-derived cartilage which forms the vertebrate skull, suggesting that persistence of this tissue and expansion into the entire headspace could be a viable evolutionary route to formation of the vertebrate head.[6]
|
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|
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The heads of humans and other animals are commonly recurring charges in heraldry.[10] Heads of humans are sometimes blazoned simply as a "man's head", but are far more frequently described in greater detail, either characteristic of a particular race or nationality (such as Moors' heads, Saxons' heads, Egyptians' heads or Turks' heads), or specifically identified (such as the head of Moses in the crest of Hilton, or the head of St. John the Baptist in the crest of the London Company of Tallowchandlers).[10][11] Several varieties of women's heads also occur, including maidens' heads (often couped under the bust, with hair disheveled), ladies' heads, nuns' heads (often veiled), and occasionally queens' heads. The arms of Devaney of Norfolk include "three nun's heads veiled couped at the shoulders proper," and the bust of a queen occurs in the arms of Queenborough, Kent.[10] Infants' or children's heads are often couped at the shoulders with a snake wrapped around the neck (e.g. "Argent, a boy's head proper, crined or, couped below the shoulders, vested gules, tarnished gold," in the arms of Boyman).[10]
|
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|
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+
The flag of Corsica displays a head in profile view
|
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+
|
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+
One of the ways of drawing sketches of heads—as Jack Hamm advises—is to develop it in six well-defined steps starting with the shape of the head in the shape of an egg. Female head, in particular, is sketched in a double circle design procedure with proportions considered as an ideal of a female head. In the first circle, the division is made of five sections on the diameter, each section of five eyes width. It is then developed over a series of ten defined steps with the smaller circle imposed partially over the larger circle at the lower end at the fourth stage. Eyes and chins are fitted in various shapes to form the head.[12]
|
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|
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Leonardo Da Vinci was one of the world’s greatest artists who drew sketches of human anatomy using grid structures. His image of the face drawn on the grid structure principle is in perfect proportion.[13] In this genre, using the technique of pen and ink, Leonardo created a sketch which is a "Study on the proportions of head and eyes" (pictured).[14]
|
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|
43 |
+
An idiom is a phrase or a fixed expression that has a figurative, or sometimes literal, meaning.
|
44 |
+
|
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The head's function and appearance play an analogous role in the etymology of many technical terms. Cylinder head, pothead, and weatherhead are three such examples.
|
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+
|
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+
Nerves of the human head, from Gray's Anatomy, 1858
|
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+
|
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+
Head of St. John the Baptist by Andrea Vaccaro, oil on canvas, 17th century
|
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|
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Sculpture of the beheaded Sainte Solange, patron saint of the French province of Berry
|
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|
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+
Patron saints of Zürich, fresco, c. 1400–1425
|
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|
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Heraldic depiction of a bison head cabossed
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Basilicata (UK: /bəˌsɪlɪˈkɑːtə/,[4] US: /-ˌzɪl-/,[5] Italian: [baziliˈkaːta]), also known by its ancient name Lucania (/luːˈkeɪniə/, also US: /luːˈkɑːnjə/,[6][7] Italian: [luˈkaːnja]), is an administrative region in Southern Italy, bordering on Campania to the west, Apulia (Puglia) to the north and east, and Calabria to the south. It has two coastlines: a 30-km stretch on the Tyrrhenian Sea between Campania and Calabria, and a longer coastline along the Gulf of Taranto between Calabria and Apulia. The region can be thought of as the "instep" of Italy, with Calabria functioning as the "toe" and Apulia the "heel".
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The region covers about 10,000 km2 (3,900 sq mi). In 2010 the population was slightly under 600,000. The regional capital is Potenza. The region is divided into two provinces: Potenza and Matera.[8][9]
|
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Basilicata is an emerging tourist destination, thanks in particular to the city of Matera, whose historical quarter I Sassi was designated in 1993 as a UNESCO World Heritage Site.[10] In 2019 it was designated as the European Capital of Culture for that year. The New York Times ranked Basilicata third in its list of "52 Places to Go in 2018", describing it as "Italy’s best-kept secret".[11]
|
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|
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The name probably derives from "basilikos" (Greek: βασιλικός), which refers to the basileus, the Byzantine emperor, who ruled the region for 200 years, from 536/552 to 571/590 and from 879 to 1059. Others argue that the name may refer to the Basilica of Acerenza, which held judicial power in the Middle Ages.
|
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|
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During the Greek and Roman Ages, Basilicata was known as Lucania. This was possibly derived from "leukos" (Greek: λευκός), meaning "white", from "lykos" (Greek: λύκος), meaning "wolf", or from Latin word "lucus", meaning "sacred wood".
|
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|
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Basilicata covers an extensive part of the southern Apennine Mountains, between the Ofanto river in the north and the Pollino massif in the south. It is bordered on the east by a large part of the Bradano river depression, which is traversed by numerous streams and declines to the southeastern coastal plains on the Ionian Sea. The region also has a short coastline to the southwest on the Tyrrhenian Sea side of the peninsula.
|
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|
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Basilicata is the most mountainous region in the south of Italy, with 47% of its area of 9,992 km2 (3,858 sq mi) covered by mountains. Of the remaining area, 45% is hilly, and 8% is made up of plains. Notable mountains and ranges include the Pollino massif, the Dolomiti lucane, Monte Vulture, Monte Alpi, Monte Carmine, Monti Li Foj and Toppa Pizzuta.
|
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|
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Geological features of the region include the volcanic formations of Monte Vulture, and the seismic faults in the Melfi and Potenza areas in the north, and around Pollino in the south. Much of the region was devastated in the 1857 Basilicata earthquake. More recently, the 1980 Irpinia earthquake destroyed many towns in the northwest of the region.
|
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The mountainous terrain combined with weak rock and soil types makes landslides prevalent. The lithological structure of the substratum and its chaotic tectonic deformation predispose the slope to landslides, and this problem is compounded by the lack of forested land. In common with many another Mediterranean region, Basilicata was once rich in forests, but they were largely felled and made barren during the time of Roman rule.
|
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The variable climate is influenced by three coastlines (Adriatic, Ionian and Tyrrhenian) and the complexity of the region's physical features. In general, the climate is continental in the mountains and Mediterranean along the coasts.
|
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|
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The first traces of human presence in Basilicata date to the late Paleolithic, with findings of Homo erectus. Late Cenozoic fossils, found at Venosa and other locations, include elephants, rhinoceros and species now extinct such as a saber-toothed cat of the genus Machairodus. Examples of rock art from the Mesolithic have been discovered near Filiano. From the fifth millennium, people stopped living in caves and built settlements of huts up to the rivers leading to the interior (Tolve, Tricarico, Aliano, Melfi, Metaponto). In this period, anatomically modern humans lived by cultivating cereals and animal husbandry (Bovinae and Caprinae). Chalcolithic sites include the grottoes of Latronico and the funerary findings of the Cervaro grotto near Lagonegro.
|
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|
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The first known stable market center of the Apennine culture on the sea, consisting of huts on the promontory of Capo la Timpa [it], near to Maratea, dates to the Bronze Age.
|
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|
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The first indigenous Iron Age communities lived in large villages in plateaus located at the borders of the plains and the rivers, in places fitting their breeding and agricultural activities. Such settlements include that of Anglona, located between the fertile valleys of Agri and Sinni, of Siris and, on the coast of the Ionian Sea, of Incoronata-San Teodoro. The first presence of Greek colonists, coming from the Greek islands and Anatolia, date from the late eighth century BC.
|
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|
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There are virtually no traces of survival of the 11th–8th century BC archaeological sites of the settlements (aside from a necropolis at Castelluccio on the coast of the Tyrrhenian Sea): this was perhaps caused by the increasing presence of Greek colonies, which changed the balance of the trades.
|
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|
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In ancient historical times the region was originally known as Lucania, named for the Lucani, an Oscan-speaking population from central Italy. Their name might be derived from Greek leukos meaning "white", lykos ("gray wolf"), or Latin lucus ("sacred grove"). Or more probably Lucania, as much as the Lucius forename (praenomen) derives from the Latin word Lux (gen. lucis), meaning "light" (<PIE *leuk- "brightness", Latin verb lucere "to shine"), and is a cognate of name Lucas. Another etymology proposed is a derivation from Etruscan Lauchum (or Lauchme) meaning "king", which however was transferred into Latin as Lucumo.[12]
|
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|
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Starting from the late eighth century BC, the Greeks established a settlement first at Siris, founded by fugitives from Colophon. Then with the foundation of Metaponto from Achaean colonists, they started the conquest of the whole Ionian coast. There were also indigenous Oenotrian foundations on the coast, which exploited the nearby presence of Greek settlements, such as Velia and Pyxous, for their maritime trades.
|
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|
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The first contacts between the Lucanians and the Romans date from the latter half of the fourth century BC. After the conquest of Taranto in 272, Roman rule was extended to the whole region: the Appian Way reached Brindisi and the colonies of Potentia (modern Potenza) and Grumentum were founded.
|
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|
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After the fall of the Western Roman Empire in 476, Basilicata fell to Germanic rule, which ended in the mid-6th century when the Byzantines reconquered it from the Ostrogoths between 536 and 552 during the apocalyptic Byzantine-Gothic war under the leadership of Byzantine generals Belisarius and Narses. The region, deeply Christianized since as early as the 5th century, became part of the Lombard Duchy of Benevento founded by the invading Lombards between 571 and 590. In the following centuries, Saracen raids led part of the population to move from the plain and coastal settlements to more protected centers located on hills. The towns of Tricarico and Tursi were under Muslim rule for a short period: later the "Saracen" population would be expelled.[13] The region was conquered once more for Byzantium from the Saracens and the Lombards in the late 9th century, with the campaigns of Nikephoros Phokas the Elder and his successors, and became part of the theme of Longobardia. In 968 the theme of Lucania was established, with the capital at Tursikon (Tursi). In 1059, Basilicata, together with the rest of much of southern Italy, was conquered by the Italo-Normans. Later, it was inherited by the Hohenstaufen, who were ousted in the 13th century by the Capetian House of Anjou.
|
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|
37 |
+
In 1485, Basilicata was the seat of plotters against King Ferdinand I of Naples, the so-called "Conspiracy of the Barons", which included the Sanseverino of Tricarico, the Caracciolo of Melfi, the Gesualdo of Caggiano, the Orsini Del Balzo [it] of Altamura and Venosa and other anti-Aragonese families. Later, Charles V stripped most of the barons of their lands, replacing them with the Carafa, Revertera, Pignatelli and Colonna among others. After the formation of the Neapolitan Republic (1647), Basilicata also rebelled, but the revolt was suppressed. In 1663 a new province was created in Basilicata with its capital in Matera.
|
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|
39 |
+
The region became part of the Kingdom of the Two Sicilies in 1735. Basilicata autonomously declared its annexation to the Kingdom of Italy on August 18, 1860 with the Potenza insurrection. It was during this period that the State confiscated and sold off vast tracts of Basilicata's territory formerly owned by the Catholic Church. As the new owners were a handful of wealthy aristocratic families, the average citizen did not see any immediate economic and social improvements after unification, and poverty continued unabated. This gave rise to the phenomenon of Brigandage in Southern Italy after 1861, whereby the Church encouraged the local people to rise up against the nobility and the new Italian state. This strong opposition movement continued for many years. Carmine Crocco from Rionero in Vulture was the most important chief in the region and the most impressive leader in southern Italy.[14]
|
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|
41 |
+
It was only really after World War II that things slowly began to improve thanks to land reform. In 1952, the inhabitants of the Sassi di Matera were rehoused by the State, but many of Basilicata's population had emigrated or were in the process of emigrating, which led to a demographic crisis from which it is still recovering.
|
42 |
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|
43 |
+
In 1993, UNESCO declared the Sassi di Matera a World Heritage Site. Meanwhile, Fiat Automobiles established a huge factory in Melfi, leading to jobs and an upsurge in the economy. In the same year the Pollino National Park was established.
|
44 |
+
|
45 |
+
Cultivation consists mainly of sowables (especially wheat), which represent 46% of the total land. Potatoes and maize are produced in the mountain areas. Olives and vine production is relatively small with about 31,000 hectares (77,000 acres) under cultivation.[15] The terrain is mountainous and hilly with poor transportation routes that hinders harvesting. Most oils are sold unbranded and only 3% is exported. The main olive cultivars are Ogliarola del Vulture, Ogliarola del Bradano, Majatica di Ferrandina and Farasana with only Ogliarola del Vulture having the Protected Designation of Origin (PDO).[16] Other varieties are the Arnasca, Ascolana, Augellina, Cellina, Frantoio, Leccino, Majatica, Nostrale, Ogliarola (Ogliarola Barese), Palmarola or Fasolina, Rapolese di Lavello, and Sargano (Sargano di Fermo and Sargano di San Benedetto).[17]
|
46 |
+
|
47 |
+
A quality wine called "Aglianico del Vulture" is produced around Rionero and received the d.o.c. classification in 1971. There are several wines such as Vino Spumante Rosso d.o.c., Aglianico di Matera from Matera produced from a mixture of the Montepulciano and other grapes, Lambrusco del Basento made from Lambrusco Maestri grapes, and Malvasia del Vulture from Malvasia grapes.[18] According to the latest Census of Agriculture, there are large herds of cattle (77,711 head in 2000).[19]
|
48 |
+
|
49 |
+
Among industrial activities, the manufacturing sector contributes to the gross value added of the secondary sector with 64% of the total, while the building sector contributes 24%. Within the services sector, the main activities in terms of gross value added are business activities, distributive trade, education and public administration. In the last few years, new productive sectors have developed: manufacturing, automotive, and especially oil extraction. In 2009, Eni employed 230 people in this area (of whom over 50% were from Basilicata), and about 1,800 were employed in activities directly generated by Eni's operations, distributed in 80 companies of which over 50% were from Basilicata.[20] The region produced about 100,000 bbl/d (16,000 m3/d), meeting 11 percent of Italy's domestic oil demand.[21]
|
50 |
+
|
51 |
+
The Gross domestic product (GDP) of the region was 12.6 billion euros in 2018, accounting for 0.7% of Italy's economic output. GDP per capita adjusted for purchasing power was 22,200 euros or 74% of the EU27 average in the same year. The GDP per employee was 95% of the EU average.[22]
|
52 |
+
|
53 |
+
Although Basilicata has never had a large population, there have nevertheless been quite considerable fluctuations in the demographic pattern of the region. In 1881, there were 539,258 inhabitants but by 1911 the population had decreased by 11% to 485,911, mainly as a result of emigration overseas. There was a slow increase in the population until World War II, after which there was a resurgence of emigration to other countries in Europe, which continued until 1971 and the start of another period of steady increase until 1993 (611,000 inhabitants). However, in recent years the population has decreased as a result of a new wave of migration, both towards northern Italy and to other countries in Europe, and a reduction in the birth rate.[23]
|
54 |
+
|
55 |
+
The population density is very low compared to that of Italy as a whole: 59.1 inhabitants per km² compared to 200.4 nationwide in 2010. There is not a great difference between the population densities of the provinces of Matera and Potenza.[23]
|
56 |
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|
57 |
+
Basilicata is divided into two provinces:
|
58 |
+
|
59 |
+
Notable people linked to Basilicata include:
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en/5670.html.txt
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1 |
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A head is the part of an organism which usually includes the ears, brain, forehead, cheeks, chin, eyes, nose, and mouth, each of which aid in various sensory functions such as sight, hearing, smell, and taste, respectively. Some very simple animals may not have a head, but many bilaterally symmetric forms do, regardless of size.
|
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|
5 |
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Heads develop in animals by an evolutionary trend known as cephalization. In bilaterally symmetrical animals, nervous tissues concentrate at the anterior region, forming structures responsible for information processing. Through biological evolution, sense organs and feeding structures also concentrate into the anterior region; these collectively form the head.
|
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|
7 |
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The human head is an anatomical unit that consists of the skull, hyoid bone and cervical vertebrae. The term "skull" collectively denotes the mandible (lower jaw bone) and the cranium (upper portion of the skull that houses the brain). The skull can also be described as being composed of the cranium, which encloses the cranial cavity, and the facial skeleton (includes the mandible). There are eight bones in the cranium and fourteen in the facial skeleton.[1]
|
8 |
+
|
9 |
+
Sculptures of human heads are generally based on a skeletal structure that consists of a cranium, jawbone, and cheekbone. Though the number of muscles making up the face is generally consistent between sculptures, the shape of the muscles varies widely based on the function, development, and expressions reflected on the faces of the subjects.[2]
|
10 |
+
|
11 |
+
Proponents of identism believe that the mind is identical to the brain. Philosopher John Searle asserts his identist beliefs, stating "the brain is the only thing in the human head". Similarly, Dr. Henry Bennet-Clark has stated that the head encloses billions of "miniagents and microagents (with no single Boss)".[3]
|
12 |
+
|
13 |
+
Mid-sagittal section of a human skull, by Leonardo da Vinci, c. 1489
|
14 |
+
|
15 |
+
Transection of a human head, by Thomas Bartholin, 1673
|
16 |
+
|
17 |
+
The evolution of a head is associated with the cephalization that occurred in Bilateria some 555 million years ago.
|
18 |
+
|
19 |
+
In some arthropods, especially trilobites (pictured at right), the cephalon, or cephalic region, is the region of the head which is a collective of "fused segments".[4]
|
20 |
+
|
21 |
+
A typical insect head is composed of eyes, antennae, and components of mouth. As these components differ substantially from insect to insect, they form important identification links. Eyes in the head found, in several types of insects, are in the form of a pair of compound eyes with multiple faces. In many other types of insects the compound eyes are seen in a "single facet or group of single facets". In some case, the eyes may be seen as marks on the dorsal or located near or toward the head, two or three ocelli (single faceted organs).[5]
|
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+
|
23 |
+
Antennae on the insect's head is found in the form of segmented attachments, in pairs, that are usually located between the eyes. These are in varying shapes and sizes, in the form of filaments or in different enlarged or clubbed form.[5]
|
24 |
+
|
25 |
+
Insects have mouth parts in various shapes depending on their feeding habits. Labrum is the "upper lip" which is in the front area of the head and is the most exterior part. A pair of mandible is found on backside of the labrum flanking the side of the mouth, succeeded by a pair of maxillae each of which is known as maxilliary palp. At the back side of the mouth is the labium or lower lip. There is also an extra mouth part in some insects which is termed as hypopharynx which is usually located between the maxillac.[5]
|
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|
27 |
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|
28 |
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|
29 |
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Though invertebrate chordates – such as the tunicate larvae or the lancelets – have heads, there has been a question of how the vertebrate head, characterized by a bony skull clearly separated from the main body, might have evolved from the head structures of these animals.[6]
|
30 |
+
|
31 |
+
According to Hyman (1979), the evolution of the head in the vertebrates has occurred by the fusion of a fixed number of anterior segments, in the same manner as in other "heteronomously segmented animals". In some cases, segments or a portion of the segments disappear. The head segments also lose most of its systems except for the nervous system. With the progressive development of cephalization, "the head incorporates more and more of the adjacent segments into its structure, so that in general it may be said that the higher the degree of cephalization the greater is the number of segments composing the head".[7]
|
32 |
+
|
33 |
+
In the 1980s, the "new head hypothesis" was proposed, suggesting that the vertebrate head is an evolutionary novelty resulting from the emergence of neural crest and cranial placodes.[8][9] In 2014, a transient larva tissue of the lancelet was found to be virtually indistinguishable from the neural crest-derived cartilage which forms the vertebrate skull, suggesting that persistence of this tissue and expansion into the entire headspace could be a viable evolutionary route to formation of the vertebrate head.[6]
|
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+
|
35 |
+
The heads of humans and other animals are commonly recurring charges in heraldry.[10] Heads of humans are sometimes blazoned simply as a "man's head", but are far more frequently described in greater detail, either characteristic of a particular race or nationality (such as Moors' heads, Saxons' heads, Egyptians' heads or Turks' heads), or specifically identified (such as the head of Moses in the crest of Hilton, or the head of St. John the Baptist in the crest of the London Company of Tallowchandlers).[10][11] Several varieties of women's heads also occur, including maidens' heads (often couped under the bust, with hair disheveled), ladies' heads, nuns' heads (often veiled), and occasionally queens' heads. The arms of Devaney of Norfolk include "three nun's heads veiled couped at the shoulders proper," and the bust of a queen occurs in the arms of Queenborough, Kent.[10] Infants' or children's heads are often couped at the shoulders with a snake wrapped around the neck (e.g. "Argent, a boy's head proper, crined or, couped below the shoulders, vested gules, tarnished gold," in the arms of Boyman).[10]
|
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+
|
37 |
+
The flag of Corsica displays a head in profile view
|
38 |
+
|
39 |
+
One of the ways of drawing sketches of heads—as Jack Hamm advises—is to develop it in six well-defined steps starting with the shape of the head in the shape of an egg. Female head, in particular, is sketched in a double circle design procedure with proportions considered as an ideal of a female head. In the first circle, the division is made of five sections on the diameter, each section of five eyes width. It is then developed over a series of ten defined steps with the smaller circle imposed partially over the larger circle at the lower end at the fourth stage. Eyes and chins are fitted in various shapes to form the head.[12]
|
40 |
+
|
41 |
+
Leonardo Da Vinci was one of the world’s greatest artists who drew sketches of human anatomy using grid structures. His image of the face drawn on the grid structure principle is in perfect proportion.[13] In this genre, using the technique of pen and ink, Leonardo created a sketch which is a "Study on the proportions of head and eyes" (pictured).[14]
|
42 |
+
|
43 |
+
An idiom is a phrase or a fixed expression that has a figurative, or sometimes literal, meaning.
|
44 |
+
|
45 |
+
The head's function and appearance play an analogous role in the etymology of many technical terms. Cylinder head, pothead, and weatherhead are three such examples.
|
46 |
+
|
47 |
+
Nerves of the human head, from Gray's Anatomy, 1858
|
48 |
+
|
49 |
+
Head of St. John the Baptist by Andrea Vaccaro, oil on canvas, 17th century
|
50 |
+
|
51 |
+
Sculpture of the beheaded Sainte Solange, patron saint of the French province of Berry
|
52 |
+
|
53 |
+
Patron saints of Zürich, fresco, c. 1400–1425
|
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+
|
55 |
+
Heraldic depiction of a bison head cabossed
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1 |
+
Coordinates: 31°N 100°W / 31°N 100°W / 31; -100
|
2 |
+
|
3 |
+
Texas (/ˈtɛksəs/, also locally /ˈtɛksɪz/;[10] Spanish: Texas or Tejas, pronounced [ˈtexas] (listen)) is a state in the South Central Region of the United States. It is the second largest U.S. state by both area (after Alaska) and population (after California). Texas shares borders with the states of Louisiana to the east, Arkansas to the northeast, Oklahoma to the north, New Mexico to the west, and the Mexican states of Chihuahua, Coahuila, Nuevo León, and Tamaulipas to the southwest, and has a coastline with the Gulf of Mexico to the southeast.
|
4 |
+
|
5 |
+
Houston is the most populous city in Texas and the fourth largest in the U.S., while San Antonio is the second-most populous in the state and seventh largest in the U.S. Dallas–Fort Worth and Greater Houston are the fourth and fifth largest metropolitan statistical areas in the country, respectively. Other major cities include Austin, the second-most populous state capital in the U.S., and El Paso. Texas is nicknamed the "Lone Star State" for its former status as an independent republic, and as a reminder of the state's struggle for independence from Mexico. The "Lone Star" can be found on the Texas state flag and on the Texas state seal.[11] The origin of Texas's name is from the word táyshaʼ, which means "friends" in the Caddo language.[12]
|
6 |
+
|
7 |
+
Due to its size and geologic features such as the Balcones Fault, Texas contains diverse landscapes common to both the U.S. Southern and the Southwestern regions.[13] Although Texas is popularly associated with the U.S. southwestern deserts, less than ten percent of Texas's land area is desert.[14] Most of the population centers are in areas of former prairies, grasslands, forests, and the coastline. Traveling from east to west, one can observe terrain that ranges from coastal swamps and piney woods, to rolling plains and rugged hills, and finally the desert and mountains of the Big Bend.
|
8 |
+
|
9 |
+
The term "six flags over Texas"[note 1] refers to several nations that have ruled over the territory. Spain was the first European country to claim and control the area of Texas. France held a short-lived colony. Mexico controlled the territory until 1836 when Texas won its independence, becoming an independent republic. In 1845,[15] Texas joined the union as the 28th state. The state's annexation set off a chain of events that led to the Mexican–American War in 1846. A slave state before the American Civil War, Texas declared its secession from the U.S. in early 1861, and officially joined the Confederate States of America on March 2 of the same year. After the Civil War and the restoration of its representation in the federal government, Texas entered a long period of economic stagnation.
|
10 |
+
|
11 |
+
Historically four major industries shaped the Texas economy prior to World War II: cattle and bison, cotton, timber, and oil.[16] Before and after the U.S. Civil War the cattle industry, which Texas came to dominate, was a major economic driver for the state, thus creating the traditional image of the Texas cowboy. In the later 19th century cotton and lumber grew to be major industries as the cattle industry became less lucrative. It was ultimately, though, the discovery of major petroleum deposits (Spindletop in particular) that initiated an economic boom which became the driving force behind the economy for much of the 20th century. With strong investments in universities, Texas developed a diversified economy and high tech industry in the mid-20th century. As of 2015, it is second on the list of the most Fortune 500 companies with 54.[17] With a growing base of industry, the state leads in many industries, including tourism, agriculture, petrochemicals, energy, computers and electronics, aerospace, and biomedical sciences. Texas has led the U.S. in state export revenue since 2002 and has the second-highest gross state product. If Texas were a sovereign state, it would be the 10th largest economy in the world.
|
12 |
+
|
13 |
+
The name Texas, based on the Caddo word táyshaʼ (/t'ajʃaʔ/) "friend", was applied, in the spelling Tejas or Texas,[18] by the Spanish to the Caddo themselves, specifically the Hasinai Confederacy,[19] the final -s representing the Spanish plural.[20]
|
14 |
+
The Mission San Francisco de los Tejas was completed near the Hasinai village of Nabedaches in May 1690, in what is now Houston County, East Texas.[21]
|
15 |
+
|
16 |
+
During Spanish colonial rule, in the 18th century, the area was known as Nuevas Filipinas ("New Philippines") and
|
17 |
+
Nuevo Reino de Filipinas ("New Kingdom of the Philippines"),[22] or as provincia de los Tejas ("province of the Tejas"),[23] later also provincia de Texas (or de Tejas), ("province of Texas").[24][22]
|
18 |
+
It was incorporated as provincia de Texas into the Mexican Empire in 1821, and declared a republic in 1836.
|
19 |
+
The Royal Spanish Academy recognizes both spellings, Tejas and Texas, as Spanish-language forms of the name of the U.S. State of Texas.[25]
|
20 |
+
|
21 |
+
The English pronunciation with /ks/ is unetymological, and based in the value of the letter x in historical Spanish orthography. Alternative etymologies of the name advanced in the late 19th century connected the Spanish teja "rooftile", the plural tejas being used to designate indigenous Pueblo settlements.[26] A 1760s map by Jacques-Nicolas Bellin shows a village named Teijas on Trinity River, close to the site of modern Crockett.[26]
|
22 |
+
|
23 |
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Texas is the second-largest U.S. state, after Alaska, with an area of 268,820 square miles (696,200 km2). Though 10% larger than France and almost twice as large as Germany or Japan and more than twice the size of the United Kingdom, it ranks only 27th worldwide amongst country subdivisions by size. If it were an independent country, Texas would be the 40th largest behind Chile and Zambia.
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Texas is in the south central part of the United States of America. Three of its borders are defined by rivers. The Rio Grande forms a natural border with the Mexican states of Chihuahua, Coahuila, Nuevo León, and Tamaulipas to the south. The Red River forms a natural border with Oklahoma and Arkansas to the north. The Sabine River forms a natural border with Louisiana to the east. The Texas Panhandle has an eastern border with Oklahoma at 100° W, a northern border with Oklahoma at 36°30' N and a western border with New Mexico at 103° W. El Paso lies on the state's western tip at 32° N and the Rio Grande.[27]
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With 10 climatic regions, 14 soil regions and 11 distinct ecological regions, regional classification becomes problematic with differences in soils, topography, geology, rainfall, and plant and animal communities.[28] One classification system divides Texas, in order from southeast to west, into the following: Gulf Coastal Plains, Interior Lowlands, Great Plains, and Basin and Range Province.
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The Gulf Coastal Plains region wraps around the Gulf of Mexico on the southeast section of the state. Vegetation in this region consists of thick piney woods. The Interior Lowlands region consists of gently rolling to hilly forested land and is part of a larger pine-hardwood forest.
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The Great Plains region in central Texas spans through the state's panhandle and Llano Estacado to the state's hill country near Austin. This region is dominated by prairie and steppe. "Far West Texas" or the "Trans-Pecos" region is the state's Basin and Range Province. The most varied of the regions, this area includes Sand Hills, the Stockton Plateau, desert valleys, wooded mountain slopes and desert grasslands.
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Texas has 3,700 named streams and 15 major rivers,[29][30] with the Rio Grande as the largest. Other major rivers include the Pecos, the Brazos, Colorado, and Red River. While Texas has few natural lakes, Texans have built more than a hundred artificial reservoirs.[31]
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The size and unique history of Texas make its regional affiliation debatable; it can be fairly considered a Southern or a Southwestern state, or both. The vast geographic, economic, and cultural diversity within the state itself prohibits easy categorization of the whole state into a recognized region of the United States. Notable extremes range from East Texas which is often considered an extension of the Deep South, to Far West Texas which is generally acknowledged to be part of the interior Southwest.[32]
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Texas is the southernmost part of the Great Plains, which ends in the south against the folded Sierra Madre Occidental of Mexico. The continental crust forms a stable Mesoproterozoic craton which changes across a broad continental margin and transitional crust into true oceanic crust of the Gulf of Mexico. The oldest rocks in Texas date from the Mesoproterozoic and are about 1,600 million years old.
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These Precambrian igneous and metamorphic rocks underlie most of the state, and are exposed in three places: Llano uplift, Van Horn, and the Franklin Mountains, near El Paso. Sedimentary rocks overlay most of these ancient rocks. The oldest sediments were deposited on the flanks of a rifted continental margin, or passive margin that developed during Cambrian time.
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This margin existed until Laurasia and Gondwana collided in the Pennsylvanian subperiod to form Pangea. This is the buried crest of the Appalachian Mountains–Ouachita Mountains zone of Pennsylvanian continental collision. This orogenic crest is today buried beneath the Dallas–Waco—Austin–San Antonio trend.
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The late Paleozoic mountains collapsed as rifting in the Jurassic period began to open the Gulf of Mexico. Pangea began to break up in the Triassic, but seafloor spreading to form the Gulf of Mexico occurred only in the mid- and late Jurassic. The shoreline shifted again to the eastern margin of the state and the Gulf of Mexico's passive margin began to form. Today 9 to 12 miles (14 to 19 km) of sediments are buried beneath the Texas continental shelf and a large proportion of remaining US oil reserves are here. At the start of its formation, the incipient Gulf of Mexico basin was restricted and seawater often evaporated completely to form thick evaporite deposits of Jurassic age. These salt deposits formed salt dome diapirs, and are found in East Texas along the Gulf coast.[33]
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East Texas outcrops consist of Cretaceous and Paleogene sediments which contain important deposits of Eocene lignite. The Mississippian and Pennsylvanian sediments in the north; Permian sediments in the west; and Cretaceous sediments in the east, along the Gulf coast and out on the Texas continental shelf contain oil. Oligocene volcanic rocks are found in far west Texas in the Big Bend area. A blanket of Miocene sediments known as the Ogallala formation in the western high plains region is an important aquifer.[34] Located far from an active plate tectonic boundary, Texas has no volcanoes and few earthquakes.[35]
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A wide range of animals and insects live in Texas. It is the home to 65 species of mammals, 213 species of reptiles and amphibians, and the greatest diversity of bird life in the United States—590 native species in all.[36] At least 12 species have been introduced and now reproduce freely in Texas.[37]
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Texas plays host to several species of wasps, including an abundance of Polistes exclamans,[38] and is an important ground for the study of Polistes annularis.
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During the spring Texas wildflowers such as the state flower, the bluebonnet, line highways throughout Texas. During the Johnson Administration the first lady, Lady Bird Johnson, worked to draw attention to Texas wildflowers.
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The large size of Texas and its location at the intersection of multiple climate zones gives the state highly variable weather. The Panhandle of the state has colder winters than North Texas, while the Gulf Coast has mild winters. Texas has wide variations in precipitation patterns. El Paso, on the western end of the state, averages 8.7 inches (220 mm) of annual rainfall,[39] while parts of southeast Texas average as much as 64 inches (1,600 mm) per year.[40] Dallas in the North Central region averages a more moderate 37 inches (940 mm) per year.
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Snow falls multiple times each winter in the Panhandle and mountainous areas of West Texas, once or twice a year in North Texas, and once every few years in Central and East Texas. Snow falls south of San Antonio or on the coast only in rare circumstances. Of note is the 2004 Christmas Eve snowstorm, when 6 inches (150 mm) of snow fell as far south as Kingsville, where the average high temperature in December is 65 °F.[41]
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Maximum temperatures in the summer months average from the 80s °F (26 °C) in the mountains of West Texas and on Galveston Island to around 100 °F (38 °C) in the Rio Grande Valley, but most areas of Texas see consistent summer high temperatures in the 90 °F (32 °C) range.
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Night-time summer temperatures range from the upper 50s °F (14 °C) in the West Texas mountains[42] to 80 °F (27 °C) in Galveston.[43]
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The table below consists of averages for August (generally the warmest month) and January (generally the coldest) in selected cities in various regions of the state. El Paso and Amarillo are exceptions with July and December respectively being the warmest and coldest months respectively, but with August and January being only narrowly different.
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Thunderstorms strike Texas often, especially the eastern and northern portions of the state. Tornado Alley covers the northern section of Texas. The state experiences the most tornadoes in the United States, an average of 139 a year. These strike most frequently in North Texas and the Panhandle.[45] Tornadoes in Texas generally occur in the months of April, May, and June.[46]
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64 |
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Some of the most destructive hurricanes in U.S. history have impacted Texas. A hurricane in 1875 killed about 400 people in Indianola, followed by another hurricane in 1886 that destroyed the town. These events allowed Galveston to take over as the chief port city. The 1900 Galveston hurricane subsequently devastated that city, killing about 8,000 people or possibly as many as 12,000. This makes it the deadliest natural disaster in U.S. history.[47] In 2017, Hurricane Harvey made landfall in Rockport as a Category 4 Hurricane, causing significant damage there. The storm stalled over land for a very long time, allowing it to drop unprecedented amounts of rain over the Greater Houston area and surrounding counties. The result was widespread and catastrophic flooding that inundated hundreds of thousands of homes. Harvey ultimately became the costliest hurricane worldwide, causing an estimated $198.6 billion in damage, surpassing the cost of Hurricane Katrina.[48]
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66 |
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67 |
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Other devastating Texas hurricanes include the 1915 Galveston hurricane, Hurricane Audrey in 1957 which killed more than 600 people, Hurricane Carla in 1961, Hurricane Beulah in 1967, Hurricane Alicia in 1983, Hurricane Rita in 2005, and Hurricane Ike in 2008. Tropical storms have also caused their share of damage: Allison in 1989 and again during 2001, and Claudette in 1979 among them.
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As of 2017[update] Texas emits the most greenhouse gases in the U.S, almost twice the amount of California, the second most polluting state.[49]
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As of 2017[update] the state emits about 1,600 billion pounds (707 million metric tons) of carbon dioxide annually.[49] As an independent nation, Texas would rank as the world's seventh-largest producer of greenhouse gases.[50] Causes of the state's vast greenhouse gas emissions include the state's large number of coal power plants and the state's refining and manufacturing industries.[50] In 2010, there were 2,553 "emission events" which poured 44.6 million pounds (20,200 metric tons) of contaminants into the Texas sky.[51]
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Texas lies between two major cultural spheres of Pre-Columbian North America: the Southwestern and the Plains areas. Archaeologists have found that three major indigenous cultures lived in this territory, and reached their developmental peak before the first European contact. These were:[52]
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When Europeans arrived in the Texas region, there were several races of Native peoples divided into many smaller tribes. They were Caddoan, Atakapan, Athabaskan, Coahuiltecan, and Uto-Aztecan. The Uto-Aztecan Puebloan peoples lived neared the Rio Grande in the western portion of the state, the Athabaskan-speaking Apache tribes lived throughout the interior, the Caddoans controlled much of the Red River region and the Atakapans were mostly centered along the Gulf Coast. At least one tribe of Coahuiltecans, the Aranama, lived in southern Texas. This entire culture group, primarily centered in northeastern Mexico, is now extinct. It is difficult to say who lived in the northwestern region of the state originally. By the time the region came to be explored, it belonged to the fairly well-known Comanche, another Uto-Aztecan people who had transitioned into a powerful horse culture, but it is believed that they came later and did not live there during the 16th century. It may have been claimed by several different peoples, including Uto-Aztecans, Athabaskans, or even Dhegihan Siouans.
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No culture was dominant in the present-day Texas region, and many peoples inhabited the area.[52] Native American tribes who lived inside the boundaries of present-day Texas include the Alabama, Apache, Atakapan, Bidai, Caddo, Aranama, Comanche, Choctaw, Coushatta, Hasinai, Jumano, Karankawa, Kickapoo, Kiowa, Tonkawa, and Wichita.[53][54] The name Texas derives from táyshaʔ, a word in the Caddoan language of the Hasinai, which means "friends" or "allies".[1][55][56][57][58]
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The region was primarily controlled by the Spanish for the first couple centuries of contact, until the Texas Revolution. They were not particularly kind to their native populations—even less so with the Caddoans, who were not trusted as their culture was split between the Spanish and the French. When the Spanish briefly managed to conquer the Louisiana colony, they decided to switch tactics and attempt being exceedingly friendly to the Indians, which they continued even after the French took back the colony. After the 1803 Louisiana Purchase, the United States inherited this odd circumstance. The Caddoans preferred the company of Americans[according to whom?] and almost the entire population of them migrated into the states of Louisiana and Arkansas. The Spanish felt jilted after having spent so much time and effort and began trying to lure the Caddo back, even promising them more land. Seemingly without actually knowing how they came by it,[according to whom?] the United States (who had begun convincing tribes to self-segregate from whites by selling everything and moving west ever since they gained the Louisiana Purchase) faced an overflow of native peoples in Missouri and Arkansas and were able to negotiate with the Caddoans to allow several displaced peoples to settle on unused lands in eastern Texas. They included the Muscogee, Houma Choctaw, Lenape and Mingo Seneca, among others, who all came to view the Caddoans as saviors, making those peoples highly influential.[59][60]
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Whether a Native American tribe was friendly or warlike was critical to the fates of European explorers and settlers in that land.[61] Friendly tribes taught newcomers how to grow indigenous crops, prepare foods, and hunt wild game. Warlike tribes made life difficult and dangerous for Europeans through their attacks and resistance to the newcomers.[62]
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During the Texas Revolution, the U.S. became heavily involved. Prior treaties with the Spanish forbade either side from militarizing its native population in any potential conflict between the two nations. At that time, several sudden outbreaks of violence between Caddoans and Texans started to spread. The Caddoans were always clueless[speculation?] when questioned, The Texan and American authorities in the region could never find hard evidence linking them to it and often it was so far-flung from Caddoan lands, it barely made any sense. It seems most likely that these were false-flag attacks meant to start a cascading effect to force the natives under Caddoan influence into armed conflict without breaking any treaties—preferably on the side of the Spanish. While no proof was found as to who the culprit was, those in charge of Texas at the time attempted multiple times to publicly blame and punish the Caddoans for the incidents with the U.S. government trying to keep them in check. Furthermore, the Caddoans never turned to violence because of it, excepting cases of self-defense.[59]
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By the 1830s, the U.S. had drafted the Indian Removal Act, which was used to facilitate the Trail of Tears. Fearing retribution of other native peoples, Indian Agents all over the eastern U.S. began desperately trying to convince all their native peoples to uproot and move west. This included the Caddoans of Louisiana and Arkansas. Following the Texas Revolution, the Texans chose to make peace with their Native peoples but did not honor former land claims or agreements. This began the movement of Native populations north into what would become Indian Territory—modern-day Oklahoma.[59]
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The first historical document related to Texas was a map of the Gulf Coast, created in 1519 by Spanish explorer Alonso Álvarez de Pineda.[63] Nine years later, shipwrecked Spanish explorer Álvar Núñez Cabeza de Vaca and his cohort became the first Europeans in what is now Texas.[64][65] Cabeza de Vaca reported that in 1528, when the Spanish landed in the area, "half the natives died from a disease of the bowels and blamed us."[66] Cabeza de Vaca also made observations about the way of life of the Ignaces Natives of Texas:
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They went about with a firebrand, setting fire to the plains and timber so as to drive off the mosquitos, and also to get lizards and similar things which they eat, to come out of the soil. In the same manner they kill deer, encircling them with fires, and they do it also to deprive the animals of pasture, compelling them to go for food where the Indians want.[67]
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Francisco Vázquez de Coronado describes his 1541 encounter:
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Two kinds of people travel around these plains with the cows; one is called Querechos and the others Teyas; they are very well built, and painted, and are enemies of each other. They have no other settlement or location than comes from traveling around with the cows. They kill all of these they wish and tan the hides, with which they clothe themselves and make their tents, and they eat the flesh, sometimes even raw, and they also even drink the blood when thirsty. The tents they make are like field tents, and they set them up over poles they have made for this purpose, which come together and are tied at the top, and when they go from one place to another they carry them on some dogs they have, of which they have many, and they load them with the tents and poles and other things, for the country is so level, as I said, that they can make use of these, because they carry the poles dragging along on the ground. The sun is what they worship most.[68]
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European powers ignored the area until accidentally settling there in 1685. Miscalculations by René-Robert Cavelier de La Salle resulted in his establishing the colony of Fort Saint Louis at Matagorda Bay rather than along the Mississippi River.[69] The colony lasted only four years before succumbing to harsh conditions and hostile natives.[70]
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In 1690 Spanish authorities, concerned that France posed a competitive threat, constructed several missions in East Texas.[71] After Native American resistance, the Spanish missionaries returned to Mexico.[72] When France began settling Louisiana, mostly in the southern part of the state, in 1716 Spanish authorities responded by founding a new series of missions in East Texas.[73] Two years later, they created San Antonio as the first Spanish civilian settlement in the area.[74]
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Hostile native tribes and distance from nearby Spanish colonies discouraged settlers from moving to the area. It was one of New Spain's least populated provinces.[76] In 1749, the Spanish peace treaty with the Lipan Apache[77] angered many tribes, including the Comanche, Tonkawa, and Hasinai.[78] The Comanche signed a treaty with Spain in 1785[79] and later helped to defeat the Lipan Apache and Karankawa tribes.[80] With more numerous missions being established, priests led a peaceful conversion of most tribes. By the end of the 18th century only a few nomadic tribes had not converted to Christianity.[81]
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When the United States purchased Louisiana from France in 1803, American authorities insisted the agreement also included Texas. The boundary between New Spain and the United States was finally set at the Sabine River in 1819, at what is now the border between Texas and Louisiana.[82] Eager for new land, many United States settlers refused to recognize the agreement. Several filibusters raised armies to invade the area west of the Sabine River.[83] In 1821, the Mexican War of Independence included the Texas territory, which became part of Mexico.[84] Due to its low population, Mexico made the area part of the state of Coahuila y Tejas.[85]
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Hoping more settlers would reduce the near-constant Comanche raids, Mexican Texas liberalized its immigration policies to permit immigrants from outside Mexico and Spain.[86] Under the Mexican immigration system, large swathes of land were allotted to empresarios, who recruited settlers from the United States, Europe, and the Mexican interior. The first grant, to Moses Austin, was passed to his son Stephen F. Austin after his death.
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Austin's settlers, the Old Three Hundred, made places along the Brazos River in 1822.[87] Twenty-three other empresarios brought settlers to the state, the majority of whom were from the United States.[88] The population of Texas grew rapidly. In 1825, Texas had about 3,500 people, with most of Mexican descent.[89] By 1834, the population had grown to about 37,800 people, with only 7,800 of Mexican descent.[90] Most of these early settlers who arrived with Austin and soon after were persons less than fortunate in life, as Texas was devoid of the comforts found elsewhere in Mexico and the United States during that time. Early Texas settler David B. Edwards described his fellow Texans as being "banished from the pleasures of life".[91]
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Many immigrants openly flouted Mexican law, especially the prohibition against slavery. Combined with United States' attempts to purchase Texas, Mexican authorities decided in 1830 to prohibit continued immigration from the United States.[92] Illegal immigration from the United States into Mexico continued to increase the population of Texas anyway.[93] New laws also called for the enforcement of customs duties angering native Mexican citizens (Tejanos) and recent immigrants alike.[94]
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The Anahuac Disturbances in 1832 were the first open revolt against Mexican rule, and they coincided with a revolt in Mexico against the nation's president.[95] Texians sided with the federalists against the current government and drove all Mexican soldiers out of East Texas.[96] They took advantage of the lack of oversight to agitate for more political freedom. Texians met at the Convention of 1832 to discuss requesting independent statehood, among other issues.[97] The following year, Texians reiterated their demands at the Convention of 1833.[98]
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Within Mexico, tensions continued between federalists and centralists. In early 1835, wary Texians formed Committees of Correspondence and Safety.[99] The unrest erupted into armed conflict in late 1835 at the Battle of Gonzales.[100] This launched the Texas Revolution, and over the next two months the Texians defeated all Mexican troops in the region.[101] Texians elected delegates to the Consultation, which created a provisional government.[102] The provisional government soon collapsed from infighting, and Texas was without clear governance for the first two months of 1836.[103]
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During this time of political turmoil, Mexican President Antonio López de Santa Anna personally led an army to end the revolt.[104] The Mexican expedition was initially successful. General José de Urrea defeated all the Texian resistance along the coast culminating in the Goliad massacre.[105] Santa Anna's forces, after a thirteen-day siege, overwhelmed Texian defenders at the Battle of the Alamo. News of the defeats sparked panic among Texas settlers.[106]
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The newly elected Texian delegates to the Convention of 1836 quickly signed a Declaration of Independence on March 2, forming the Republic of Texas. After electing interim officers, the Convention disbanded.[107] The new government joined the other settlers in Texas in the Runaway Scrape, fleeing from the approaching Mexican army.[106] After several weeks of retreat, the Texian Army commanded by Sam Houston attacked and defeated Santa Anna's forces at the Battle of San Jacinto.[108] Santa Anna was captured and forced to sign the Treaties of Velasco, ending the war.[109] The Constitution of the Republic of Texas prohibited the government from restricting slavery or freeing slaves, required free people of African descent to leave the country, and prohibited Native Americans from becoming citizens.
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While Texas had won its independence, political battles raged between two factions of the new Republic. The nationalist faction, led by Mirabeau B. Lamar, advocated the continued independence of Texas, the expulsion of the Native Americans, and the expansion of the Republic to the Pacific Ocean. Their opponents, led by Sam Houston, advocated the annexation of Texas to the United States and peaceful co-existence with Native Americans. The conflict between the factions was typified by an incident known as the Texas Archive War.[110] With wide popular support, Texas first applied for annexation to the United States in 1836, but its status as a slaveholding country caused its admission to be controversial and it was initially rebuffed. This status, and Mexican diplomacy in support of its claims to the territory, also complicated Texas's ability to form foreign alliances and trade relationships.[111]
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Mexico launched two small expeditions into Texas in 1842. The town of San Antonio was captured twice and Texans were defeated in battle in the Dawson massacre. Despite these successes, Mexico did not keep an occupying force in Texas, and the republic survived.[112] The cotton price crash of the 1840s depressed the country's economy,[111] contributing to the republic's inability to defend itself, and adding momentum to Texas's eventual annexation into the United States.
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As early as 1837, the Republic made several attempts to negotiate annexation with the United States.[113] Opposition within the republic from the nationalist faction, along with strong abolitionist opposition within the United States, slowed Texas's admission into the Union. Texas was finally annexed when the expansionist James K. Polk won the election of 1844.[114] On December 29, 1845, Congress admitted Texas to the U.S. as a constituent state of the Union.[115]
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The population of the new state was quite small at first, and there was a strong mix between the English-speaking American settlers who dominated in the state's eastern/northeastern portions and the Spanish-speaking former Mexicans who dominated in the state's southern and western portions. Statehood brought many new settlers. Because of the long Spanish presence in Mexico and various failed colonization efforts by the Spanish and Mexicans in northern Mexico, there were large herds of Longhorn cattle that roamed the state. Hardy by nature, but also suitable for slaughtering and consumption, they represented an economic opportunity many entrepreneurs seized upon, thus creating the cowboy culture for which Texas is famous. While in the early days of the republic cattle and bison were slaughtered for their hides, soon a beef industry was established with cattle being shipped all over the U.S. and the Caribbean (within a few decades, beef had become a staple of the American diet).
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After Texas's annexation, Mexico broke diplomatic relations with the United States. While the United States claimed Texas's border stretched to the Rio Grande, Mexico claimed it was the Nueces River leaving the Rio Grande Valley under contested Texan sovereignty.[115] While the former Republic of Texas could not enforce its border claims, the United States had the military strength and the political will to do so. President Polk ordered General Zachary Taylor south to the Rio Grande on January 13, 1846. A few months later Mexican troops routed an American cavalry patrol in the disputed area in the Thornton Affair starting the Mexican–American War. The first battles of the war were fought in Texas: the Siege of Fort Texas, Battle of Palo Alto and Battle of Resaca de la Palma. After these decisive victories, the United States invaded Mexican territory, ending the fighting in Texas.[116]
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After a series of United States victories, the Treaty of Guadalupe Hidalgo ended the two-year war. In return, for US$18,250,000, Mexico gave the U.S. undisputed control of Texas, ceded the Mexican Cession in 1848, most of which today is called the American Southwest, and Texas's borders were established at the Rio Grande.[116]
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The Compromise of 1850 set Texas's boundaries at their present form. U.S. Senator James Pearce of Maryland drafted the final proposal[27] where Texas ceded its claims to land which later became half of present-day New Mexico, a third of Colorado, and small portions of Kansas, Oklahoma, and Wyoming to the federal government, in return for the assumption of $10 million of the old republic's debt.[27] Post-war Texas grew rapidly as migrants poured into the cotton lands of the state.[117]
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They also brought or purchased enslaved African Americans, whose numbers tripled in the state from 1850 to 1860, from 58,000 to 182,566.[118]
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Texas was at war again after the election of 1860. At this time, blacks comprised 30 percent of the state's population, and they were overwhelmingly enslaved.[119] When Abraham Lincoln was elected, South Carolina seceded from the Union. Five other Lower South states quickly followed. A State Convention considering secession opened in Austin on January 28, 1861. On February 1, by a vote of 166–8, the Convention adopted an Ordinance of Secession from the United States. Texas voters approved this Ordinance on February 23, 1861. Texas joined the newly created Confederate States of America on March 4, 1861 ratifying the permanent C.S. Constitution on March 23.[1][120]
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Not all Texans favored secession initially, although many of the same would later support the Southern cause. Texas's most notable Unionist was the state Governor, Sam Houston. Not wanting to aggravate the situation, Houston refused two offers from President Lincoln for Union troops to keep him in office. After refusing to swear an oath of allegiance to the Confederacy, Houston was deposed as governor.[121]
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While far from the major battlefields of the American Civil War, Texas contributed large numbers of men and equipment to the rest of the Confederacy.[122] Union troops briefly occupied the state's primary port, Galveston. Texas's border with Mexico was known as the "backdoor of the Confederacy" because trade occurred at the border, bypassing the Union blockade.[123] The Confederacy repulsed all Union attempts to shut down this route,[122] but Texas's role as a supply state was marginalized in mid-1863 after the Union capture of the Mississippi River. The final battle of the Civil War was fought near Brownsville, Texas at Palmito Ranch[124] with a Confederate victory.
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Texas descended into anarchy for two months between the surrender of the Army of Northern Virginia and the assumption of authority by Union General Gordon Granger. Violence marked the early months of Reconstruction.[122] Juneteenth commemorates the announcement of the Emancipation Proclamation in Galveston by General Gordon Granger, almost two and a half years after the original announcement.[125][126] President Johnson, in 1866, declared the civilian government restored in Texas.[127] Despite not meeting reconstruction requirements, Congress resumed allowing elected Texas representatives into the federal government in 1870. Social volatility continued as the state struggled with agricultural depression and labor issues.[128]
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Like most of the South, the Texas economy was devastated by the War. However, since the state had not been as dependent on slaves as other parts of the South, it was able to recover more quickly. The culture in Texas during the later 19th century exhibited many facets of a frontier territory. The state became notorious as a haven for people from other parts of the country who wanted to escape debt, criminal prosecution, or other problems. Indeed, "Gone to Texas" was a common expression for those fleeing the law in other states. Nevertheless, the state also attracted many businessmen and other settlers with more legitimate interests as well.
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The cattle industry continued to thrive, though it gradually became less profitable. Cotton and lumber became major industries creating new economic booms in various regions of the state. Railroad networks grew rapidly as did the port at Galveston as commerce between Texas and the rest of the U.S. (and the rest of the world) expanded. As with some other states before, the lumber industry quickly decimated the forests of Texas such that, by the early 20th century, the majority of the forest population in Texas was gone (later conservation efforts restored some of it, but never to the level it once was).
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In 1900, Texas suffered the deadliest natural disaster in U.S. history during the Galveston hurricane.[47] On January 10, 1901, the first major oil well in Texas, Spindletop, was found south of Beaumont. Other fields were later discovered nearby in East Texas, West Texas, and under the Gulf of Mexico. The resulting "oil boom" transformed Texas.[129] Oil production eventually averaged three million barrels per day at its peak in 1972.[130]
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In 1901, the Democratic-dominated state legislature passed a bill requiring payment of a poll tax for voting, which effectively disenfranchised most blacks and many poor whites and Latinos. In addition, the legislature established white primaries, ensuring minorities were excluded from the formal political process. The number of voters dropped dramatically, and the Democrats crushed competition from the Republican and Populist parties.[131][132] The Socialist Party became the second-largest party in Texas after 1912,[133] coinciding with a large socialist upsurge in the United States during fierce battles in the labor movement and the popularity of national heroes like Eugene V. Debs. The Socialists' popularity soon waned after their vilification by the United States government for their opposition to US involvement in World War I.
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The Great Depression and the Dust Bowl dealt a double blow to the state's economy, which had significantly improved since the Civil War. Migrants abandoned the worst-hit sections of Texas during the Dust Bowl years. Especially from this period on, blacks left Texas in the Great Migration to get work in the Northern United States or California and to escape the oppression of segregation.[119] In 1940, Texas was 74 percent Anglo, 14.4 percent black, and 11.5 percent Hispanic.[134]
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World War II had a dramatic impact on Texas, as federal money poured in to build military bases, munitions factories, POW detention camps and Army hospitals; 750,000 young men left for service; the cities exploded with new industry; the colleges took on new roles; and hundreds of thousands of poor farmers left the fields for much better-paying war jobs, never to return to agriculture.[135][136] Texas manufactured 3.1 percent of total United States military armaments produced during World War II, ranking eleventh among the 48 states.[137]
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Texas modernized and expanded its system of higher education through the 1960s. The state created a comprehensive plan for higher education, funded in large part by oil revenues, and a central state apparatus designed to manage state institutions more efficiently. These changes helped Texas universities receive federal research funds.[138]
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On November 22, 1963, President John F. Kennedy was assassinated in Dallas.[139]
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Beginning around the mid-20th century, Texas began to transform from a rural and agricultural state to one urban and industrialized.[140] The state's population grew quickly during this period, with large levels of migration from outside the state.[140] As a part of the Sun Belt, Texas experienced strong economic growth, particularly during the 1970s and early 1980s.[140] Texas's economy diversified, lessening its reliance on the petroleum industry.[140] By 1990, Hispanics overtook blacks to become the largest minority group in the state.[140]
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During the late 20th century, the Republican Party replaced the Democratic Party as the dominant party in the state, as the latter became more politically liberal and as demographic changes favored the former.[140]
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The current Texas Constitution was adopted in 1876. Like many states, it explicitly provides for a separation of powers. The state's Bill of Rights is much larger than its federal counterpart, and has provisions unique to Texas.[141]
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Texas has a plural executive branch system limiting the power of the governor, which is a weak executive compared to some other states. Except for the Secretary of State, voters elect executive officers independently; thus candidates are directly answerable to the public, not the governor.[142] This election system has led to some executive branches split between parties and reduced the ability of the governor to carry out a program. When Republican President George W. Bush served as Texas's governor, the state had a Democratic lieutenant governor, Bob Bullock. The executive branch positions consist of the Governor, Lieutenant Governor, Comptroller of Public Accounts, Land Commissioner, Attorney General, Agriculture Commissioner, the three-member Texas Railroad Commission, the State Board of Education, and the Secretary of State.[142]
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The bicameral Texas Legislature consists of the House of Representatives, with 150 members, and a Senate, with 31 members. The Speaker of the House leads the House, and the lieutenant governor, the Senate.[143] The Legislature meets in regular session biennially for just over a hundred days, but the governor can call for special sessions as often as desired (notably, the Legislature cannot call itself into session).[144] The state's fiscal year begins September 1.
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The judiciary of Texas is one of the most complex in the United States, with many layers and overlapping jurisdictions. Texas has two courts of last resort: the Texas Supreme Court, for civil cases, and the Texas Court of Criminal Appeals. Except for some municipal benches, partisan elections select judges at all levels of the judiciary; the governor fills vacancies by appointment.[145] Texas is notable for its use of capital punishment, having led the country in executions since capital punishment was reinstated in the Gregg v. Georgia case (see Capital punishment in Texas).
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The Texas Ranger Division of the Texas Department of Public Safety is a law enforcement agency with statewide jurisdiction. Over the years, the Texas Rangers have investigated crimes ranging from murder to political corruption. They have acted as riot police and as detectives, protected the Texas governor, tracked down fugitives, and functioned as a paramilitary force both for the republic and for the state. The Texas Rangers were unofficially created by Stephen F. Austin in 1823 and formally constituted in 1835. The Rangers were integral to several important events of Texas history and some of the best-known criminal cases in the history of the Old West.[146]
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The Texas constitution defines the responsibilities of county governments, which serve as agents of the state. What are called commissioners court and court judges are elected to serve as the administrative arm. Most cities in the state, those over 5,000 in population, have home-rule governments. The vast majority of these have charters for council-manager forms of government, by which voters elect council members, who hire a professional city manager as an operating officer.
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In the 1870s, white Democrats wrested power back in the state legislature from the biracial coalition at the end of Reconstruction. In the early 20th century, the legislature passed bills to impose poll taxes, followed by white primaries; these measures effectively disfranchised most blacks, poor whites and Mexican Americans.[131][132] In the 1890s, 100,000 blacks voted in the state; by 1906, only 5,000 could vote.[147] As a result, the Democratic Party dominated Texas politics from the turn of the century, imposing racial segregation and white supremacy. It held power until after passage in the mid-1960s of national civil rights legislation enforcing constitutional rights of all citizens.
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Although Texas was essentially a one-party state during this time and the Democratic primary was viewed as "the real election", the Democratic Party had conservative and liberal factions, which became more pronounced after the New Deal.[148] Additionally, several factions of the party briefly split during the 1930s and 1940s.[148]
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The state's conservative white voters began to support Republican presidential candidates by the mid-20th century. After this period, they supported Republicans for local and state offices as well, and most whites became Republican Party members.[149] The party also attracted some minorities, but many have continued to vote for Democratic candidates. The shift to the Republican Party is much-attributed to the fact the Democratic Party became increasingly liberal during the 20th century, and thus increasingly out-of-touch with the average Texas voter.[150] As Texas was always a conservative state, voters switched to the GOP, which now more closely reflected their beliefs.[150][151] Commentators have also attributed the shift to Republican political consultant Karl Rove, who managed numerous political campaigns in Texas in the 1980s and 1990s.[151] Other stated reasons included court-ordered redistricting and the demographic shift in relation to the Sun Belt that favored the Republican Party and conservatism.[140]
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The 2003 Texas redistricting of Congressional districts led by Republican Tom DeLay, was called by The New York Times "an extreme case of partisan gerrymandering".[152] A group of Democratic legislators, the "Texas Eleven", fled the state in a quorum-busting effort to prevent the legislature from acting, but was unsuccessful.[153] The state had already redistricted following the 2000 census. Despite these efforts, the legislature passed a map heavily in favor of Republicans, based on 2000 data and ignoring the estimated nearly one million new residents in the state since then. Career attorneys and analysts at the Department of Justice objected to the plan as diluting the votes of African American and Hispanic voters, but political appointees overrode them and approved it.[152] Legal challenges to the redistricting reached the national Supreme Court in the case League of United Latin American Citizens v. Perry (2006), but the court ruled in favor of the state (and Republicans).[154]
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In the 2014 Texas elections, the Tea Party movement made large gains, with numerous Tea Party favorites being elected into office, including Dan Patrick as lieutenant governor,[155][156] Ken Paxton as attorney general,[155][157] in addition to numerous other candidates[157] including conservative Republican Greg Abbott as governor.[158]
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Texas voters lean toward fiscal conservatism, while enjoying the benefits of huge federal investment in the state in military and other facilities achieved by the power of the Solid South in the 20th century. They also tend to have socially conservative values.[159][160]
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Since 1980, most Texas voters have supported Republican presidential candidates. In 2000 and 2004, Republican George W. Bush won Texas with respectively 59.3 and 60.1 percent of the vote, partly due to his "favorite son" status as a former governor of the state. John McCain won the state in 2008, but with a smaller margin of victory compared to Bush at 55 percent of the vote. Austin, Dallas, Houston, and San Antonio consistently lean Democratic in both local and statewide elections.
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The state's changing demographics may result in a change in its overall political alignment, as a majority population of Black and Hispanic/Latino voters support the Democratic Party.[161] Residents of counties along the Rio Grande closer to the Mexico–United States border, where there are many Latino residents, generally vote for Democratic Party candidates, while most other rural and suburban areas of Texas have shifted to voting for Republican Party candidates.[162][163]
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As of the general elections of 2014, a large majority of the members of Texas's U.S. House delegation are Republican, along with both U.S. Senators. In the 114th United States Congress, of the 36 Congressional districts in Texas, 24 are held by Republicans and 11 by Democrats. One seat is vacant. Texas's Senators are John Cornyn and Ted Cruz. Since 1994, Texans have not elected a Democrat to a statewide office. The state's Democratic voters are made up primarily by liberal and minority groups in Austin, Beaumont, Dallas, El Paso, Houston, and San Antonio as well as minority voters in East and South Texas.
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Texas has banned sanctuary cities,[164] but Houston Mayor Sylvester Turner has vowed that the city will not assist ICE agents.[165]
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Texas has 254 counties—the most nationwide. Each county runs on Commissioners' Court system consisting of four elected commissioners (one from each of four precincts in the county, roughly divided according to population) and a county judge elected at large from the entire county. County government runs similar to a "weak" mayor-council system; the county judge has no veto authority, but votes along with the other commissioners.
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Although Texas permits cities and counties to enter "interlocal agreements" to share services, the state does not allow consolidated city-county governments, nor does it have metropolitan governments. Counties are not granted home rule status; their powers are strictly defined by state law. The state does not have townships—areas within a county are either incorporated or unincorporated. Incorporated areas are part of a municipality. The county provides limited services to unincorporated areas and to some smaller incorporated areas. Municipalities are classified either "general law" cities or "home rule".[167] A municipality may elect home rule status once it exceeds 5,000 population with voter approval.
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Texas also permits the creation of "special districts", which provide limited services. The most common is the school district, but can also include hospital districts, community college districts, and utility districts (one utility district near Austin was the plaintiff in a landmark Supreme Court case involving the Voting Rights Act).
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Municipal, school district, and special district elections are nonpartisan,[168] though the party affiliation of a candidate may be well-known. County and state elections are partisan.
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Texas has a reputation of very harsh criminal punishment for criminal offenses. It is one of the 32 states that practice capital punishment, and since the US Supreme Court allowed capital punishment to resume in 1976, 40% of all US executions have taken place in Texas.[169] As of 2008, Texas had the 4th highest incarceration rate in the US.[170] Texas also has strong self defense laws, allowing citizens to use lethal force to defend themselves, their families, or their property.[171]
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As of 2019, Texas had a gross state product (GSP) of $1.9 trillion, the second highest in the U.S.[172] Its GSP is greater than the GDPs of Brazil, Canada, Russia, South Korea and Spain, which are the world's 9th-, 10th-, 11th-, 12th- and 13th-largest economies, respectively.[173] Texas's economy is the second-largest of any country subdivision globally, behind California. Its per capita personal income in 2009 was $36,484, ranking 29th in the nation.[174]
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Texas's large population, an abundance of natural resources, thriving cities and leading centers of higher education have contributed to a large and diverse economy. Since oil was discovered, the state's economy has reflected the state of the petroleum industry. In recent times, urban centers of the state have increased in size, containing two-thirds of the population in 2005. The state's economic growth has led to urban sprawl and its associated symptoms.[174]
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As of May 2020, during the COVID-19 pandemic, the state's unemployment rate was 13 percent.[175]
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In 2010, Site Selection Magazine ranked Texas as the most business-friendly state in the nation, in part because of the state's three-billion-dollar Texas Enterprise Fund.[176] Texas has the joint-highest number of Fortune 500 company headquarters in the United States, along with California.[177][178]
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In 2010, there were 346,000 millionaires in Texas, constituting the second-largest population of millionaires in the nation.[179][180]
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Texas has a "low taxes, low services" reputation.[159] According to the Tax Foundation, Texans' state and local tax burdens rank among the lowest in the nation, 7th lowest nationally; state and local taxes cost $3,580 per capita, or 8.4 percent of resident incomes.[181] Texas is one of seven states that lack a state income tax.[181][182]
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Instead, the state collects revenue from property taxes (though these are collected at the county, city, and school district level; Texas has a state constitutional prohibition against a state property tax) and sales taxes. The state sales tax rate is 6.25 percent,[181][183] but local taxing jurisdictions (cities, counties, special purpose districts, and transit authorities) may also impose sales and use tax up to 2 percent for a total maximum combined rate of 8.25 percent.[184]
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Texas is a "tax donor state"; in 2005, for every dollar Texans paid to the federal government in federal income taxes, the state got back about $0.94 in benefits.[181] To attract business, Texas has incentive programs worth $19 billion per year (2012); more than any other US state.[185][186]
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Texas has the most farms and the highest acreage in the United States. The state is ranked No. 1 for revenue generated from total livestock and livestock products. It is ranked No. 2 for total agricultural revenue, behind California.[187] At $7.4 billion or 56.7 percent of Texas's annual agricultural cash receipts, beef cattle production represents the largest single segment of Texas agriculture. This is followed by cotton at $1.9 billion (14.6 percent), greenhouse/nursery at $1.5 billion (11.4 percent), broilers at $1.3 billion (10 percent), and dairy products at $947 million (7.3 percent).[188]
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Texas leads the nation in the production of cattle, horses, sheep, goats, wool, mohair and hay.[188] The state also leads the nation in production of cotton[187][189] which is the number one crop grown in the state in terms of value.[190] The state grows significant amounts of cereal crops and produce.[187] Texas has a large commercial fishing industry. With mineral resources, Texas leads in creating cement, crushed stone, lime, salt, sand and gravel.[187]
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Texas throughout the 21st century has been hammered by drought. This has cost the state billions of dollars in livestock and crops.[191]
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Ever since the discovery of oil at Spindletop, energy has been a dominant force politically and economically within the state.[192] If Texas were its own country it would be the sixth largest oil producer in the world.[193]
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The Railroad Commission of Texas, contrary to its name, regulates the state's oil and gas industry, gas utilities, pipeline safety, safety in the liquefied petroleum gas industry, and surface coal and uranium mining. Until the 1970s, the commission controlled the price of petroleum because of its ability to regulate Texas's oil reserves. The founders of the Organization of Petroleum Exporting Countries (OPEC) used the Texas agency as one of their models for petroleum price control.[194]
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Texas has known petroleum deposits of about 5 billion barrels (790,000,000 m3), which makes up about one-fourth of the known U.S. reserves.[195] The state's refineries can process 4.6 million barrels (730,000 m3) of oil a day.[195] The Port Arthur Refinery in Southeast Texas is the largest refinery in the U.S.[195] Texas also leads in natural gas production, producing one-fourth of the nation's supply.[195] Several petroleum companies are based in Texas such as: Occidental Petroleum, ConocoPhillips, ExxonMobil, Halliburton, Marathon Oil, Tesoro, Valero Energy, and Western Refining.
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According to the Energy Information Administration, Texans consume, on average, the fifth most energy (of all types) in the nation per capita and as a whole, following behind Wyoming, Alaska, Louisiana, North Dakota, and Iowa.[195]
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Unlike the rest of the nation, most of Texas is on its own alternating current power grid, the Texas Interconnection. Texas has a deregulated electric service. Texas leads the nation in total net electricity production, generating 437,236 MWh in 2014, 89% more MWh than Florida, which ranked second.[196][197] As an independent nation, Texas would rank as the world's eleventh-largest producer of electricity, after South Korea, and ahead of the United Kingdom.
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The state is a leader in renewable energy commercialization; it produces the most wind power in the nation.[195][198] In 2014, 10.6% of the electricity consumed in Texas came from wind turbines.[199] The Roscoe Wind Farm in Roscoe, Texas, is one of the world's largest wind farms with a 781.5 megawatt (MW) capacity.[200] The Energy Information Administration states the state's large agriculture and forestry industries could give Texas an enormous amount biomass for use in biofuels. The state also has the highest solar power potential for development in the nation.[195]
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With large universities systems coupled with initiatives like the Texas Enterprise Fund and the Texas Emerging Technology Fund, a wide array of different high tech industries have developed in Texas. The Austin area is nicknamed the "Silicon Hills" and the north Dallas area the "Silicon Prairie". Many high-tech companies are located in or have their headquarters in Texas (and Austin in particular), including Dell, Inc., Borland, Forcepoint, Indeed.com, Texas Instruments, Perot Systems, Rackspace and AT&T.
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The National Aeronautics and Space Administration's Lyndon B. Johnson Space Center (NASA JSC) in Southeast Houston, sits as the crown jewel of Texas's aeronautics industry. Fort Worth hosts both Lockheed Martin's Aeronautics division and Bell Helicopter Textron.[201][202] Lockheed builds the F-16 Fighting Falcon, the largest Western fighter program, and its successor, the F-35 Lightning II in Fort Worth.[203]
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Texas's affluence stimulates a strong commercial sector consisting of retail, wholesale, banking and insurance, and construction industries. Examples of Fortune 500 companies not based on Texas traditional industries are AT&T, Kimberly-Clark, Blockbuster, J. C. Penney, Whole Foods Market, and Tenet Healthcare.[204]
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Nationally, the Dallas–Fort Worth area, home to the second shopping mall in the United States, has the most shopping malls per capita of any American metropolitan area.[205]
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Mexico, the state's largest trading partner, imports a third of the state's exports because of the North American Free Trade Agreement (NAFTA). NAFTA has encouraged the formation of controversial maquiladoras on the Texas–Mexico border.[206]
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The United States Census Bureau estimates the population of Texas was 28,995,881 on July 1, 2019, a 15.31 percent increase since the 2010 United States Census.[208][6]
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As of 2015, Texas had 4.7 million foreign-born residents, about 17% of the population and 21.6% of the state workforce.[209] The major countries of origin for Texan immigrants were Mexico (55.1% of immigrants), India (5%), El Salvador (4.3%), Vietnam (3.7%), and China (2.3%).[209] Of immigrant residents, some 35.8 percent were naturalized U.S. citizens.[209] In 2014, there were an estimated 1.7 million undocumented immigrants in Texas, making up 35% of the total Texas immigrant population and 6.1% of the total state population.[209] In addition to the state's foreign-born population, an additional 4.1 million Texans (15% of the state's population) were born in the United States and had at least one immigrant parent.[209]
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Texas's Rio Grande Valley has seen significant migration from across the U.S.–Mexico border. During the 2014 crisis, many Central Americans, including unaccompanied minors traveling alone from Guatemala, Honduras, and El Salvador, reached the state, overwhelming Border Patrol resources for a time. Many sought asylum in the United States.[210][211]
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Texas's population density is 90.5 people per square mile (34.9/km2) which is slightly higher than the average population density of the U.S. as a whole, at 80.6 people per square mile (31.1/km2). In contrast, while Texas and France are similarly sized geographically, the European country has a population density of 301.8 people per square mile (116.5/km2).
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Two-thirds of all Texans live in major metropolitan areas such as Houston. The Dallas-Fort Worth Metropolitan Area is the largest in Texas. While Houston is the largest city in Texas and the fourth-largest city in the United States, the Dallas-Fort Worth metropolitan area is larger than Houston.
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As of the 2015 Texas Population Estimate Program, the population of the state was 27,469,114; non-Hispanic whites 11,505,371 (41.9%); Black Americans 3,171,043 (11.5%); other races 1,793,580 (6.5%); and Hispanics and Latinos (of any race) 10,999,120 (40.0%).[212]
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According to the 2010 United States census, the racial composition of Texas was the following:[213]
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In addition, 37.6 percent of the population was Hispanic or Latino (of any race) (31.6 percent Mexican, 0.9 percent Salvadoran, 0.5 percent Puerto Rican, 0.4 percent Honduran, 0.3 percent Guatemalan 0.3 percent Spaniard, 0.2 percent Colombian, 0.2 percent Cuban).[214]
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As of 2011, 69.8% of the population of Texas younger than age 1 were minorities (meaning they had at least one parent who was not non-Hispanic white).[215]
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German, Irish, and English Americans are the three largest European ancestry groups in Texas. German Americans make up 11.3 percent of the population and number over 2.7 million members. Irish Americans make up 8.2 percent of the population and number over 1.9 million. There are roughly 600,000 French Americans and 472,000 Italian Americans residing in Texas; these two ethnic groups make up 2.5 percent and 2.0 percent of the population respectively. In the 1980 United States Census the largest ancestry group reported in Texas was English with 3,083,323 Texans citing they were of English or mostly English ancestry, making them 27 percent of the state at the time.[219] Their ancestry primarily goes back to the original thirteen colonies[citation needed] (the census of 1790 gives 48% of the population of English ancestry, together with 12% Scots and Scots-Irish, 4.5% other Irish, and 3% Welsh, for a total of 67.5% British and Irish; 13% were German, Swiss, Dutch, and French Huguenots; 19% were African-American),[220] thus many of them today identify as "American" in ancestry, though they are of predominantly British stock.[221][222] There are nearly 200,000 Czech Americans living in Texas, the largest number of any state.[223]
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African Americans are the largest racial minority in Texas. Their proportion of the population has declined since the early 20th century after many left the state in the Great Migration. Blacks of both Hispanic and non-Hispanic origin make up 11.5 percent of the population; blacks of non-Hispanic origin form 11.3 percent of the populace. African Americans of both Hispanic and non-Hispanic origin number at roughly 2.7 million individuals.
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Native Americans are a smaller minority in the state. Native Americans make up 0.5 percent of Texas's population and number over 118,000 individuals. Native Americans of non-Hispanic origin make up 0.3 percent of the population and number over 75,000 individuals. Cherokee made up 0.1 percent of the population, and numbered over 19,400. In contrast, only 583 identified as Chippewa.
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Asian Americans are a sizable minority group in Texas. Americans of Asian descent form 3.8 percent of the population, with those of non-Hispanic descent making up 3.7 percent of the populace. They total more than 808,000 individuals. Non-Hispanic Asians number over 795,000. Just over 200,000 Indian Americans make Texas their home. Texas is also home to more than 187,000 Vietnamese and 136,000 Chinese. In addition to 92,000 Filipinos and 62,000 Koreans, there are 18,000 Japanese Americans living in the state. Lastly, more than 111,000 people are of other Asian ancestry groups, such as Cambodian, Thai, and Hmong. Sugar Land, a city within the Houston metropolitan area, and Plano, within the Dallas metropolitan area, both have high concentrations of ethnic Chinese and Korean residents. The Houston and Dallas areas, and to a lesser extent, the Austin metropolitan area, all contain substantial Vietnamese communities.
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Americans with origins from the Pacific Islands are the smallest minority in Texas. According to the survey, only 18,000 Texans are Pacific Islanders; 16,400 are of non-Hispanic descent. There are roughly 5,400 Native Hawaiians, 5,300 Guamanians, and 6,400 people from other groups. Samoan Americans were scant; only 2,920 people were from this group. The city of Euless, a suburb of Fort Worth, contains a sizable population of Tongan Americans, at nearly 900 people, over one percent of the city's population. Killeen has a sufficient population of Samoans and Guamanian, and people of Pacific Islander descent surpass one percent of the city's population.
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Multiracial individuals are also a visible minority in Texas. People identifying as multiracial form 1.9 percent of the population, and number over 448,000 people. Almost 80,000 Texans claim African and European heritage and makeup 0.3 percent of the population. People of European and American Indian ancestry number over 108,800 (close to the number of Native Americans), and makeup 0.5 percent of the population. People of European and Asian ancestry number over 57,600, and form just 0.2 percent of the population. People of African and Native American ancestry were even smaller in number (15,300), and makeup just 0.1 percent of the total population.
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Hispanics and Latinos are the second-largest groups in Texas after non-Hispanic European Americans. More than 8.5 million people claim Hispanic or Latino ethnicity. This group forms over 37 percent of Texas's population. People of Mexican descent alone number over 7.9 million, and make up 31.6 percent of the population. The vast majority of the Hispanic/Latino population in the state is of Mexican descent, the next two largest groups are Salvadorans and Puerto Ricans. There are more than 222,000 Salvadorans and more than 130,000 Puerto Ricans in Texas. Other groups with large numbers in Texas include Hondurans, Guatemalans, Nicaraguans and Cubans, among others.[224][225] The Hispanics in Texas are more likely than in some other states (such as California) to identify as white; according to the 2010 U.S. Census, Texas is home to 6,304,207 White Hispanics and 2,594,206 Hispanics of "some other race" (usually mestizo).
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German descendants inhabit much of central and southeast-central Texas. Over one-third of Texas, residents are of Hispanic origin; while many have recently arrived, some Tejanos have ancestors with multi-generational ties to 18th century Texas. The African American population in Texas is increasing due to the New Great Migration.[226][227] In addition to the descendants of the state's former slave population, many African American college graduates have come to the state for work recently in the New Great Migration.[226] Recently, the Asian population in Texas has grown—primarily in Houston and Dallas. Other communities with a significantly growing Asian American population is in Austin, Corpus Christi, San Antonio, and the Sharyland area next McAllen, Texas. Three federally recognized Native American tribes reside in Texas: the Alabama-Coushatta Tribe, the Kickapoo Traditional Tribe, and the Ysleta del Sur Pueblo.[54]
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In 2010, 49 percent of all births were Hispanics; 35 percent were non-Hispanic whites; 11.5 percent were non-Hispanic blacks, and 4.3 percent were Asians/Pacific Islanders.[228] Based on Census Bureau data released in February 2011, for the first time in recent history, Texas's white population is below 50 percent (45 percent) and Hispanics grew to 38 percent. Between 2000 and 2010, the total population growth by 20.6 percent, but Hispanics growth by 65 percent, whereas non-Hispanic whites grew by only 4.2 percent.[229] Texas has the fifth highest rate of teenage births in the nation and a plurality of these are to Hispanics.[230]
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The state has three cities with populations exceeding one million: Houston, San Antonio, and Dallas.[235] These three rank among the 10 most populous cities of the United States. As of 2010, six Texas cities had populations greater than 600,000 people. Austin, Fort Worth, and El Paso are among the 20 largest U.S. cities. Texas has four metropolitan areas with populations greater than a million: Dallas–Fort Worth–Arlington, Houston–Sugar Land–Baytown, San Antonio–New Braunfels, and Austin–Round Rock–San Marcos. The Dallas–Fort Worth and Houston metropolitan areas number about 6.3 million and 5.7 million residents, respectively.
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Three interstate highways—I-35 to the west (Dallas–Fort Worth to San Antonio, with Austin in between), I-45 to the east (Dallas to Houston), and I-10 to the south (San Antonio to Houston) define the Texas Urban Triangle region. The region of 60,000 square miles (160,000 km2) contains most of the state's largest cities and metropolitan areas as well as 17 million people, nearly 75 percent of Texas's total population.[236] Houston and Dallas have been recognized as beta world cities.[237] These cities are spread out amongst the state. Texas has 254 counties, which is more than any other state by 95 (Georgia).[238]
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In contrast to the cities, unincorporated rural settlements known as colonias often lack basic infrastructure and are marked by poverty.[239] The office of the Texas Attorney General stated, in 2011, that Texas had about 2,294 colonias and estimates about 500,000 lived in the colonias. Hidalgo County, as of 2011, has the largest number of colonias.[240] Texas has the largest number of people of all states, living in colonias.[239]
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The most common accent or dialect spoken by natives throughout Texas is sometimes referred to as Texan English, which itself is a sub-variety of a broader category of American English known as Southern American English.[242][243] Creole language is spoken in East Texas.[244] In some areas of the state—particularly in the large cities—Western American English and General American English, have been on the increase. Chicano English—due to a growing Hispanic population—is widespread in South Texas, while African-American English is especially notable in historically minority areas of urban Texas.
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As of 2010, 65.8% (14,740,304) of Texas residents age 5 and older spoke only English at home, while 29.2% (6,543,702) spoke Spanish, 0.75 percent (168,886) Vietnamese, and Chinese (which includes Cantonese and Mandarin) was spoken by 0.56% (122,921) of the population over five.[245]
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Other languages spoken include German (including Texas German) by 0.33% (73,137), Tagalog with 0.29% (64,272) speakers, and French (including Cajun French) was spoken by 0.25% (55,773) of Texans.[245] Reportedly, Cherokee is the most widely spoken Native American language in Texas.[246]
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In total, 34.2% (7,660,406) of Texas's population aged five and older spoke a language at home other than English.[245]
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The largest denominations by number of adherents in 2010 were the Roman Catholic Church (4,673,500); the Southern Baptist Convention (3,721,318); Non-denominational Churches (1,546,542); and the United Methodist Church with (1,035,168).[248]
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Known as the buckle of the Bible Belt, East Texas is socially conservative.[249] The Dallas–Fort Worth metroplex is home to three major evangelical seminaries and a host of Bible schools. Lakewood Church in Houston, boasts the largest attendance in the nation averaging more than 43,000 weekly.[250]
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Adherents of many other religions reside predominantly in the urban centers of Texas. In 1990, the Islamic population was about 140,000 with more recent figures putting the current number of Muslims between 350,000 and 400,000.[251] The Jewish population is around 128,000.[252] Around 146,000 adherents of religions such as Hinduism and Sikhism live in Texas.[253] It is the fifth-largest Muslim-populated state in the country.[254]
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Historically, Texas culture comes from a blend of Southern (Dixie), Western (frontier), and Southwestern (Mexican/Anglo fusion) influences, varying in degrees of such from one intrastate region to another. Texas is placed in the Southern United States by the United States Census Bureau.[255] A popular food item, the breakfast burrito, draws from all three, having a soft flour tortilla wrapped around bacon and scrambled eggs or other hot, cooked fillings. Adding to Texas's traditional culture, established in the 18th and 19th centuries, immigration has made Texas a melting pot of cultures from around the world.
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Texas has made a strong mark on national and international pop culture. The entire state is strongly associated with the image of the cowboy shown in westerns and in country western music. The state's numerous oil tycoons are also a popular pop culture topic as seen in the hit TV series Dallas.
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The internationally known slogan "Don't Mess with Texas" began as an anti-littering advertisement. Since the campaign's inception in 1986, the phrase has become "an identity statement, a declaration of Texas swagger".[256]
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"Texas-sized" is an expression that can be used in two ways: to describe something that is about the size of the U.S. state of Texas,[257][258] or to describe something (usually but not always originating from Texas) that is large compared to other objects of its type.[259][260][261] Texas was the largest U.S. state, until Alaska became a state in 1959. The phrase "everything is bigger in Texas" has been in regular use since at least 1950;[262] and was used as early as 1913.[263]
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Houston is one of only five American cities with permanent professional resident companies in all the major performing arts disciplines: the Houston Grand Opera, the Houston Symphony Orchestra, the Houston Ballet, and The Alley Theatre.[264] Known for the vibrancy of its visual and performing arts, the Houston Theater District—a 17-block area in the heart of Downtown Houston—ranks second in the country in the number of theater seats in a concentrated downtown area, with 12,948 seats for live performances and 1,480 movie seats.[264]
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Founded in 1892, Modern Art Museum of Fort Worth, also called "The Modern", is Texas's oldest art museum. Fort Worth also has the Kimbell Art Museum, the Amon Carter Museum, the National Cowgirl Museum and Hall of Fame, the Will Rogers Memorial Center, and the Bass Performance Hall downtown. The Arts District of Downtown Dallas has arts venues such as the Dallas Museum of Art, the Morton H. Meyerson Symphony Center, the Margot and Bill Winspear Opera House, the Trammell & Margaret Crow Collection of Asian Art, and the Nasher Sculpture Center.[265]
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The Deep Ellum district within Dallas became popular during the 1920s and 1930s as the prime jazz and blues hotspot in the Southern United States. The name Deep Ellum comes from local people pronouncing "Deep Elm" as "Deep Ellum".[266] Artists such as Blind Lemon Jefferson, Robert Johnson, Huddie "Lead Belly" Ledbetter, and Bessie Smith played in early Deep Ellum clubs.[267]
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Austin, The Live Music Capital of the World, boasts "more live music venues per capita than such music hotbeds as Nashville, Memphis, Los Angeles, Las Vegas or New York City".[268] The city's music revolves around the nightclubs on 6th Street; events like the film, music, and multimedia festival South by Southwest; the longest-running concert music program on American television, Austin City Limits; and the Austin City Limits Music Festival held in Zilker Park.[269]
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Since 1980, San Antonio has evolved into "The Tejano Music Capital Of The World".[270] The Tejano Music Awards have provided a forum to create greater awareness and appreciation for Tejano music and culture.[271]
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The second president of the Republic of Texas, Mirabeau B. Lamar, is the Father of Texas Education. During his term, the state set aside three leagues of land in each county for equipping public schools. An additional 50 leagues of land set aside for the support of two universities would later become the basis of the state's Permanent University Fund.[272] Lamar's actions set the foundation for a Texas-wide public school system.[273]
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Between 2006 and 2007, Texas spent $7,275 per pupil ranking it below the national average of $9,389. The pupil/teacher ratio was 14.9, below the national average of 15.3. Texas paid instructors $41,744, below the national average of $46,593. The Texas Education Agency (TEA) administers the state's public school systems. Texas has over 1,000 school districts; all districts except the Stafford Municipal School District are independent from municipal government and many cross city boundaries.[274] School districts have the power to tax their residents and to assert eminent domain over privately owned property. Due to court-mandated equitable school financing for school districts, the state has a controversial tax redistribution system called the "Robin Hood plan". This plan transfers property tax revenue from wealthy school districts to poor ones.[275] The TEA has no authority over private or home school activities.[276]
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Students in Texas take the State of Texas Assessments of Academic Readiness (STAAR) in primary and secondary school. STAAR assess students' attainment of reading, writing, mathematics, science, and social studies skills required under Texas education standards and the No Child Left Behind Act. The test replaced the Texas Assessment of Knowledge and Skills (TAKS) test in the 2011–2012 school year.[277]
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Generally prohibited in the West at large, school corporal punishment is not unusual in the more conservative, rural areas of the state, with 28,569 public school students[278] paddled at least one time, according to government data for the 2011–2012 school year.[279] The rate of school corporal punishment in Texas is surpassed only by Mississippi, Alabama, and Arkansas.[279]
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The state's two most widely recognized flagship universities are The University of Texas at Austin and Texas A&M University, ranked as the 52nd[280] and 69th[281] best universities in the nation according to the 2014 edition of U.S. News & World Report's "Best Colleges", respectively. Some observers[282] also include the University of Houston and Texas Tech University as tier one flagships alongside UT Austin and A&M.[283][284] The Texas Higher Education Coordinating Board (THECB) ranks the state's public universities into three distinct tiers:[285]
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Texas's controversial alternative affirmative action plan, Texas House Bill 588, guarantees Texas students who graduated in the top 10 percent of their high school class automatic admission to state-funded universities. The bill encourages demographic diversity while avoiding problems stemming from the Hopwood v. Texas (1996) case.
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Thirty-six (36) separate and distinct public universities exist in Texas, of which 32 belong to one of the six state university systems.[288][289] Discovery of minerals on Permanent University Fund land, particularly oil, has helped fund the rapid growth of the state's two largest university systems: the University of Texas System and the Texas A&M System. The four other university systems: the University of Houston System, the University of North Texas System, the Texas State System, and the Texas Tech System are not funded by the Permanent University Fund.
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The Carnegie Foundation classifies three of Texas's universities as Tier One research institutions: The University of Texas at Austin, the Texas A&M University, and the University of Houston. The University of Texas at Austin and Texas A&M University are flagship universities of the state of Texas. Both were established by the Texas Constitution and hold stakes in the Permanent University Fund. The state has been putting effort to expand the number of flagship universities by elevating some of its seven institutions designated as "emerging research universities". The two expected to emerge first are the University of Houston and Texas Tech University, likely in that order according to discussions on the House floor of the 82nd Texas Legislature.[290]
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The state is home to various private institutions of higher learning—ranging from liberal arts colleges to a nationally recognized top-tier research university. Rice University in Houston is one of the leading teaching and research universities of the United States and is ranked the nation's 17th-best overall university by U.S. News & World Report.[291] Trinity University, a private, primarily undergraduate liberal arts university in San Antonio, has ranked first among universities granting primarily bachelor's and select master's degrees in the Western United States for 20 consecutive years by U.S. News.[292] Private universities include Abilene Christian University, Austin College, Baylor University, University of Mary Hardin–Baylor, and Southwestern University.[293][294][295]
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Universities in Texas host three presidential libraries: George Bush Presidential Library at Texas A&M University, the Lyndon Baines Johnson Library and Museum at The University of Texas at Austin, and the George W. Bush Presidential Library at Southern Methodist University.
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Notwithstanding the concentration of elite medical centers in the state, The Commonwealth Fund ranks the Texas healthcare system the third worst in the nation.[296] Texas ranks close to last in access to healthcare, quality of care, avoidable hospital spending, and equity among various groups.[296] Causes of the state's poor rankings include politics, a high poverty rate, and the highest rate of illegal immigration in the nation.[297] In May 2006, Texas initiated the program "code red" in response to the report the state had 25.1 percent of the population without health insurance, the largest proportion in the nation.[298]
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The Trust for America's Health ranked Texas 15th highest in adult obesity, with 27.2 percent of the state's population measured as obese.[299] The 2008 Men's Health obesity survey ranked four Texas cities among the top 25 fattest cities in America; Houston ranked 6th, Dallas 7th, El Paso 8th, and Arlington 14th.[300] Texas had only one city (Austin, ranked 21st) in the top 25 among the "fittest cities" in America.[300] The same survey has evaluated the state's obesity initiatives favorably with a "B+".[300] The state is ranked forty-second in the percentage of residents who engage in regular exercise.[301]
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Texas has the highest maternal mortality rate in the developed world, and the rate by which Texas women died from pregnancy-related complications doubled from 2010 to 2014, to 23.8 per 100,000. A rate unmatched in any other U.S. state or economically developed country.[302]
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Texas has many elite research medical centers. The state has nine medical schools,[303] three dental schools,[304] and two optometry schools.[305] Texas has two Biosafety Level 4 (BSL-4) laboratories: one at The University of Texas Medical Branch (UTMB) in Galveston,[306] and the other at the Southwest Foundation for Biomedical Research in San Antonio—the first privately owned BSL-4 lab in the United States.[307]
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The Texas Medical Center in Houston, holds the world's largest concentration of research and healthcare institutions, with 47 member institutions.[308] Texas Medical Center performs the most heart transplants in the world.[309] The University of Texas M. D. Anderson Cancer Center in Houston is a highly regarded academic institution that centers around cancer patient care, research, education and prevention.[310]
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San Antonio's South Texas Medical Center facilities rank sixth in clinical medicine research impact in the United States.[311] The University of Texas Health Science Center is another highly ranked research and educational institution in San Antonio.[312][313]
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Both the American Heart Association and the University of Texas Southwestern Medical Center call Dallas home. The Southwestern Medical Center ranks "among the top academic medical centers in the world".[314] The institution's medical school employs the most medical school Nobel laureates in the world.[314][315]
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Texans have historically had difficulties traversing Texas due to the state's large size and rough terrain. Texas has compensated by building America's largest highway and railway systems. The regulatory authority, the Texas Department of Transportation (TxDOT), maintains the state's immense highway system, regulates aviation,[316] and public transportation systems.[317]
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Located centrally in North America, the state is an important transportation hub. From the Dallas/Fort Worth area, trucks can reach 93 percent of the nation's population within 48 hours, and 37 percent within 24 hours.[318] Texas has 33 foreign trade zones (FTZ), the most in the nation.[319] In 2004, a combined total of $298 billion of goods passed through Texas FTZs.[319]
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The first Texas freeway was the Gulf Freeway opened in 1948 in Houston.[320] As of 2005, 79,535 miles (127,999 km) of public highway crisscrossed Texas (up from 71,000 miles (114,263 km) in 1984).[321] To fund recent growth in the state highways, Texas has 17 toll roads (see list) with several additional tollways proposed.[322] In central Texas, the southern section of the State Highway 130 toll road has a speed limit of 85 miles per hour (137 km/h), the highest in the nation.[323] All federal and state highways in Texas are paved.
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Texas has 730 airports, second-most of any state in the nation. Largest in Texas by size and passengers served, Dallas/Fort Worth International Airport (DFW) is the second-largest by area in the United States, and fourth in the world with 18,076 acres (73.15 km2).[324] In traffic, DFW airport is the busiest in the state, the fourth busiest in the United States,[325] and sixth worldwide.[326] American Airlines Group's American / American Eagle, the world's largest airline in total passengers-miles transported and passenger fleet size,[327] uses DFW as its largest and main hub. It ranks as the largest airline in the United States by number of passengers carried domestically per year and the largest airline in the world by number of passengers carried.[328]Southwest Airlines, headquartered in Dallas, has its operations at Dallas Love Field.[329]
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Texas's second-largest air facility is Houston's George Bush Intercontinental Airport (IAH). It served as the largest hub for the former Continental Airlines, which was based in Houston; it serves as the largest hub for United Airlines, the world's third-largest airline, by passenger-miles flown.[330][331] IAH offers service to the most Mexican destinations of any U.S. airport.[332][333] The next five largest airports in the state all serve more than three million passengers annually; they include Austin-Bergstrom International Airport, William P. Hobby Airport, San Antonio International Airport, Dallas Love Field and El Paso International Airport. The smallest airport in the state to be designated an international airport is Del Rio International Airport.
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Around 1,150 seaports dot Texas's coast with over 1,000 miles (1,600 km) of channels.[334] Ports employ nearly one-million people and handle an average of 317 million metric tons.[335] Texas ports connect with the rest of the U.S. Atlantic seaboard with the Gulf section of the Intracoastal Waterway.[334] The Port of Houston today is the busiest port in the United States in foreign tonnage, second in overall tonnage, and tenth worldwide in tonnage.[336] The Houston Ship Channel spans 530 feet (160 m) wide by 45 feet (14 m) deep by 50 miles (80 km) long.[337]
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Part of the state's tradition of cowboys is derived from the massive cattle drives which its ranchers organized in the nineteenth century to drive livestock to railroads and markets in Kansas, for shipment to the East. Towns along the way, such as Baxter Springs, the first cow town in Kansas, developed to handle the seasonal workers and tens of thousands of head of cattle being driven.
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The first railroad to operate in Texas was the Buffalo Bayou, Brazos and Colorado Railway, opening in August 1853.[338] The first railroad to enter Texas from the north, completed in 1872, was the Missouri–Kansas–Texas Railroad.[339] With increasing railroad access, the ranchers did not have to take their livestock up to the Midwest and shipped beef out from Texas. This caused a decline in the economies of the cow towns.
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Since 1911, Texas has led the nation in length of railroad miles within the state. Texas railway length peaked in 1932 at 17,078 miles (27,484 km), but declined to 14,006 miles (22,540 km) by 2000. While the Railroad Commission of Texas originally regulated state railroads, in 2005 the state reassigned these duties to TxDOT.[340]
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In the Dallas–Fort Worth area, three public transit agencies provide rail service: Dallas Area Rapid Transit (DART), Denton County Transportation Authority (DCTA), and Trinity Metro. DART began operating the first light rail system in the Southwest United States in 1996.[341]
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The Trinity Railway Express (TRE) commuter rail service, which connects Fort Worth and Dallas, is provided by Trinity Metro and DART.[342] Trinity Metro also operates the TEXRail commuter rail line, connecting downtown Fort Worth and Northeast Tarrant County to DFW Airport.[343] The A-train commuter rail line, operated by DCTA, acts as an extension of the DART Green line into Denton County.[344] In the Austin area, Capital Metropolitan Transportation Authority operates a commuter rail service known as Capital MetroRail to the northwestern suburbs. The Metropolitan Transit Authority of Harris County, Texas (METRO) operates light rail lines in the Houston area.
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Amtrak provides Texas with limited intercity passenger rail service. Three scheduled routes serve the state: the daily Texas Eagle (Chicago–San Antonio); the tri-weekly Sunset Limited (New Orleans–Los Angeles), with stops in Texas; and the daily Heartland Flyer (Fort Worth–Oklahoma City). Texas may get one of the nation's first high-speed rail line. Plans for a privately funded high-speed rail line between Dallas and Houston have been planned by the Texas Central Railway company.
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While American football has long been considered "king" in the state, Texans enjoy a wide variety of sports.[345]
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Texans can cheer for a plethora of professional sports teams. Within the "Big Four" professional leagues, Texas has two NFL teams (the Dallas Cowboys and the Houston Texans), two Major League Baseball teams (the Houston Astros and the Texas Rangers), three NBA teams (the San Antonio Spurs, the Houston Rockets, and the Dallas Mavericks), and one National Hockey League team (the Dallas Stars). The Dallas–Fort Worth metroplex is one of only twelve American metropolitan areas that host sports teams from all the "Big Four" professional leagues. Outside the "Big Four", Texas also has a WNBA team, (the Dallas Wings) and two Major League Soccer teams (the Houston Dynamo and FC Dallas).
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Collegiate athletics have deep significance in Texas culture, especially football. The state has twelve Division I-FBS schools, the most in the nation. Four of the state's universities, the Baylor Bears, Texas Longhorns, TCU Horned Frogs, and Texas Tech Red Raiders, compete in the Big 12 Conference. The Texas A&M Aggies left the Big 12 and joined the Southeastern Conference in 2012, which led the Big 12 to invite TCU to join; TCU was previously in the Mountain West Conference. The Houston Cougars and the SMU Mustangs compete in the American Athletic Conference. The Texas State Bobcats and the UT Arlington Mavericks compete in the Sun Belt Conference. Four of the state's schools claim at least one national championship in football: the Texas Longhorns, the Texas A&M Aggies, the TCU Horned Frogs, and the SMU Mustangs.
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According to a survey of Division I-A coaches the rivalry between the University of Oklahoma and the University of Texas at Austin, the Red River Shootout, ranks the third-best in the nation.[346] The TCU Horned Frogs and SMU Mustangs also share a rivalry and compete annually in the Battle for the Iron Skillet. A fierce rivalry, the Lone Star Showdown, also exists between the state's two largest universities, Texas A&M University and the University of Texas at Austin. The athletics portion of the Lone Star Showdown rivalry has been put on hold after the Texas A&M Aggies joined the Southeastern Conference.
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The University Interscholastic League (UIL) organizes most primary and secondary school competitions. Events organized by UIL include contests in athletics (the most popular being high school football) as well as artistic and academic subjects.[347]
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Texans also enjoy the rodeo. The world's first rodeo was hosted in Pecos, Texas.[348] The annual Houston Livestock Show and Rodeo is the largest rodeo in the world. It begins with trail rides from several points throughout the state that convene at Reliant Park.[349] The Southwestern Exposition and Livestock Show in Fort Worth is the oldest continuously running rodeo incorporating many of the state's most historic traditions into its annual events. Dallas hosts the State Fair of Texas each year at Fair Park.[350]
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Texas Motor Speedway hosts annual NASCAR Cup Series and IndyCar Series auto races since 1997. Since 2012, Austin's Circuit of the Americas plays host to a round of the Formula 1 World Championship—[351] the first at a permanent road circuit in the United States since the 1980 Grand Prix at Watkins Glen International—, as well as Grand Prix motorcycle racing, FIA World Endurance Championship and United SportsCar Championship races.
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1 |
+
Coordinates: 31°N 100°W / 31°N 100°W / 31; -100
|
2 |
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3 |
+
Texas (/ˈtɛksəs/, also locally /ˈtɛksɪz/;[10] Spanish: Texas or Tejas, pronounced [ˈtexas] (listen)) is a state in the South Central Region of the United States. It is the second largest U.S. state by both area (after Alaska) and population (after California). Texas shares borders with the states of Louisiana to the east, Arkansas to the northeast, Oklahoma to the north, New Mexico to the west, and the Mexican states of Chihuahua, Coahuila, Nuevo León, and Tamaulipas to the southwest, and has a coastline with the Gulf of Mexico to the southeast.
|
4 |
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|
5 |
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Houston is the most populous city in Texas and the fourth largest in the U.S., while San Antonio is the second-most populous in the state and seventh largest in the U.S. Dallas–Fort Worth and Greater Houston are the fourth and fifth largest metropolitan statistical areas in the country, respectively. Other major cities include Austin, the second-most populous state capital in the U.S., and El Paso. Texas is nicknamed the "Lone Star State" for its former status as an independent republic, and as a reminder of the state's struggle for independence from Mexico. The "Lone Star" can be found on the Texas state flag and on the Texas state seal.[11] The origin of Texas's name is from the word táyshaʼ, which means "friends" in the Caddo language.[12]
|
6 |
+
|
7 |
+
Due to its size and geologic features such as the Balcones Fault, Texas contains diverse landscapes common to both the U.S. Southern and the Southwestern regions.[13] Although Texas is popularly associated with the U.S. southwestern deserts, less than ten percent of Texas's land area is desert.[14] Most of the population centers are in areas of former prairies, grasslands, forests, and the coastline. Traveling from east to west, one can observe terrain that ranges from coastal swamps and piney woods, to rolling plains and rugged hills, and finally the desert and mountains of the Big Bend.
|
8 |
+
|
9 |
+
The term "six flags over Texas"[note 1] refers to several nations that have ruled over the territory. Spain was the first European country to claim and control the area of Texas. France held a short-lived colony. Mexico controlled the territory until 1836 when Texas won its independence, becoming an independent republic. In 1845,[15] Texas joined the union as the 28th state. The state's annexation set off a chain of events that led to the Mexican–American War in 1846. A slave state before the American Civil War, Texas declared its secession from the U.S. in early 1861, and officially joined the Confederate States of America on March 2 of the same year. After the Civil War and the restoration of its representation in the federal government, Texas entered a long period of economic stagnation.
|
10 |
+
|
11 |
+
Historically four major industries shaped the Texas economy prior to World War II: cattle and bison, cotton, timber, and oil.[16] Before and after the U.S. Civil War the cattle industry, which Texas came to dominate, was a major economic driver for the state, thus creating the traditional image of the Texas cowboy. In the later 19th century cotton and lumber grew to be major industries as the cattle industry became less lucrative. It was ultimately, though, the discovery of major petroleum deposits (Spindletop in particular) that initiated an economic boom which became the driving force behind the economy for much of the 20th century. With strong investments in universities, Texas developed a diversified economy and high tech industry in the mid-20th century. As of 2015, it is second on the list of the most Fortune 500 companies with 54.[17] With a growing base of industry, the state leads in many industries, including tourism, agriculture, petrochemicals, energy, computers and electronics, aerospace, and biomedical sciences. Texas has led the U.S. in state export revenue since 2002 and has the second-highest gross state product. If Texas were a sovereign state, it would be the 10th largest economy in the world.
|
12 |
+
|
13 |
+
The name Texas, based on the Caddo word táyshaʼ (/t'ajʃaʔ/) "friend", was applied, in the spelling Tejas or Texas,[18] by the Spanish to the Caddo themselves, specifically the Hasinai Confederacy,[19] the final -s representing the Spanish plural.[20]
|
14 |
+
The Mission San Francisco de los Tejas was completed near the Hasinai village of Nabedaches in May 1690, in what is now Houston County, East Texas.[21]
|
15 |
+
|
16 |
+
During Spanish colonial rule, in the 18th century, the area was known as Nuevas Filipinas ("New Philippines") and
|
17 |
+
Nuevo Reino de Filipinas ("New Kingdom of the Philippines"),[22] or as provincia de los Tejas ("province of the Tejas"),[23] later also provincia de Texas (or de Tejas), ("province of Texas").[24][22]
|
18 |
+
It was incorporated as provincia de Texas into the Mexican Empire in 1821, and declared a republic in 1836.
|
19 |
+
The Royal Spanish Academy recognizes both spellings, Tejas and Texas, as Spanish-language forms of the name of the U.S. State of Texas.[25]
|
20 |
+
|
21 |
+
The English pronunciation with /ks/ is unetymological, and based in the value of the letter x in historical Spanish orthography. Alternative etymologies of the name advanced in the late 19th century connected the Spanish teja "rooftile", the plural tejas being used to designate indigenous Pueblo settlements.[26] A 1760s map by Jacques-Nicolas Bellin shows a village named Teijas on Trinity River, close to the site of modern Crockett.[26]
|
22 |
+
|
23 |
+
Texas is the second-largest U.S. state, after Alaska, with an area of 268,820 square miles (696,200 km2). Though 10% larger than France and almost twice as large as Germany or Japan and more than twice the size of the United Kingdom, it ranks only 27th worldwide amongst country subdivisions by size. If it were an independent country, Texas would be the 40th largest behind Chile and Zambia.
|
24 |
+
|
25 |
+
Texas is in the south central part of the United States of America. Three of its borders are defined by rivers. The Rio Grande forms a natural border with the Mexican states of Chihuahua, Coahuila, Nuevo León, and Tamaulipas to the south. The Red River forms a natural border with Oklahoma and Arkansas to the north. The Sabine River forms a natural border with Louisiana to the east. The Texas Panhandle has an eastern border with Oklahoma at 100° W, a northern border with Oklahoma at 36°30' N and a western border with New Mexico at 103° W. El Paso lies on the state's western tip at 32° N and the Rio Grande.[27]
|
26 |
+
|
27 |
+
With 10 climatic regions, 14 soil regions and 11 distinct ecological regions, regional classification becomes problematic with differences in soils, topography, geology, rainfall, and plant and animal communities.[28] One classification system divides Texas, in order from southeast to west, into the following: Gulf Coastal Plains, Interior Lowlands, Great Plains, and Basin and Range Province.
|
28 |
+
|
29 |
+
The Gulf Coastal Plains region wraps around the Gulf of Mexico on the southeast section of the state. Vegetation in this region consists of thick piney woods. The Interior Lowlands region consists of gently rolling to hilly forested land and is part of a larger pine-hardwood forest.
|
30 |
+
|
31 |
+
The Great Plains region in central Texas spans through the state's panhandle and Llano Estacado to the state's hill country near Austin. This region is dominated by prairie and steppe. "Far West Texas" or the "Trans-Pecos" region is the state's Basin and Range Province. The most varied of the regions, this area includes Sand Hills, the Stockton Plateau, desert valleys, wooded mountain slopes and desert grasslands.
|
32 |
+
|
33 |
+
Texas has 3,700 named streams and 15 major rivers,[29][30] with the Rio Grande as the largest. Other major rivers include the Pecos, the Brazos, Colorado, and Red River. While Texas has few natural lakes, Texans have built more than a hundred artificial reservoirs.[31]
|
34 |
+
|
35 |
+
The size and unique history of Texas make its regional affiliation debatable; it can be fairly considered a Southern or a Southwestern state, or both. The vast geographic, economic, and cultural diversity within the state itself prohibits easy categorization of the whole state into a recognized region of the United States. Notable extremes range from East Texas which is often considered an extension of the Deep South, to Far West Texas which is generally acknowledged to be part of the interior Southwest.[32]
|
36 |
+
|
37 |
+
Texas is the southernmost part of the Great Plains, which ends in the south against the folded Sierra Madre Occidental of Mexico. The continental crust forms a stable Mesoproterozoic craton which changes across a broad continental margin and transitional crust into true oceanic crust of the Gulf of Mexico. The oldest rocks in Texas date from the Mesoproterozoic and are about 1,600 million years old.
|
38 |
+
|
39 |
+
These Precambrian igneous and metamorphic rocks underlie most of the state, and are exposed in three places: Llano uplift, Van Horn, and the Franklin Mountains, near El Paso. Sedimentary rocks overlay most of these ancient rocks. The oldest sediments were deposited on the flanks of a rifted continental margin, or passive margin that developed during Cambrian time.
|
40 |
+
|
41 |
+
This margin existed until Laurasia and Gondwana collided in the Pennsylvanian subperiod to form Pangea. This is the buried crest of the Appalachian Mountains–Ouachita Mountains zone of Pennsylvanian continental collision. This orogenic crest is today buried beneath the Dallas–Waco—Austin–San Antonio trend.
|
42 |
+
|
43 |
+
The late Paleozoic mountains collapsed as rifting in the Jurassic period began to open the Gulf of Mexico. Pangea began to break up in the Triassic, but seafloor spreading to form the Gulf of Mexico occurred only in the mid- and late Jurassic. The shoreline shifted again to the eastern margin of the state and the Gulf of Mexico's passive margin began to form. Today 9 to 12 miles (14 to 19 km) of sediments are buried beneath the Texas continental shelf and a large proportion of remaining US oil reserves are here. At the start of its formation, the incipient Gulf of Mexico basin was restricted and seawater often evaporated completely to form thick evaporite deposits of Jurassic age. These salt deposits formed salt dome diapirs, and are found in East Texas along the Gulf coast.[33]
|
44 |
+
|
45 |
+
East Texas outcrops consist of Cretaceous and Paleogene sediments which contain important deposits of Eocene lignite. The Mississippian and Pennsylvanian sediments in the north; Permian sediments in the west; and Cretaceous sediments in the east, along the Gulf coast and out on the Texas continental shelf contain oil. Oligocene volcanic rocks are found in far west Texas in the Big Bend area. A blanket of Miocene sediments known as the Ogallala formation in the western high plains region is an important aquifer.[34] Located far from an active plate tectonic boundary, Texas has no volcanoes and few earthquakes.[35]
|
46 |
+
|
47 |
+
A wide range of animals and insects live in Texas. It is the home to 65 species of mammals, 213 species of reptiles and amphibians, and the greatest diversity of bird life in the United States—590 native species in all.[36] At least 12 species have been introduced and now reproduce freely in Texas.[37]
|
48 |
+
|
49 |
+
Texas plays host to several species of wasps, including an abundance of Polistes exclamans,[38] and is an important ground for the study of Polistes annularis.
|
50 |
+
|
51 |
+
During the spring Texas wildflowers such as the state flower, the bluebonnet, line highways throughout Texas. During the Johnson Administration the first lady, Lady Bird Johnson, worked to draw attention to Texas wildflowers.
|
52 |
+
|
53 |
+
The large size of Texas and its location at the intersection of multiple climate zones gives the state highly variable weather. The Panhandle of the state has colder winters than North Texas, while the Gulf Coast has mild winters. Texas has wide variations in precipitation patterns. El Paso, on the western end of the state, averages 8.7 inches (220 mm) of annual rainfall,[39] while parts of southeast Texas average as much as 64 inches (1,600 mm) per year.[40] Dallas in the North Central region averages a more moderate 37 inches (940 mm) per year.
|
54 |
+
|
55 |
+
Snow falls multiple times each winter in the Panhandle and mountainous areas of West Texas, once or twice a year in North Texas, and once every few years in Central and East Texas. Snow falls south of San Antonio or on the coast only in rare circumstances. Of note is the 2004 Christmas Eve snowstorm, when 6 inches (150 mm) of snow fell as far south as Kingsville, where the average high temperature in December is 65 °F.[41]
|
56 |
+
|
57 |
+
Maximum temperatures in the summer months average from the 80s °F (26 °C) in the mountains of West Texas and on Galveston Island to around 100 °F (38 °C) in the Rio Grande Valley, but most areas of Texas see consistent summer high temperatures in the 90 °F (32 °C) range.
|
58 |
+
|
59 |
+
Night-time summer temperatures range from the upper 50s °F (14 °C) in the West Texas mountains[42] to 80 °F (27 °C) in Galveston.[43]
|
60 |
+
|
61 |
+
The table below consists of averages for August (generally the warmest month) and January (generally the coldest) in selected cities in various regions of the state. El Paso and Amarillo are exceptions with July and December respectively being the warmest and coldest months respectively, but with August and January being only narrowly different.
|
62 |
+
|
63 |
+
Thunderstorms strike Texas often, especially the eastern and northern portions of the state. Tornado Alley covers the northern section of Texas. The state experiences the most tornadoes in the United States, an average of 139 a year. These strike most frequently in North Texas and the Panhandle.[45] Tornadoes in Texas generally occur in the months of April, May, and June.[46]
|
64 |
+
|
65 |
+
Some of the most destructive hurricanes in U.S. history have impacted Texas. A hurricane in 1875 killed about 400 people in Indianola, followed by another hurricane in 1886 that destroyed the town. These events allowed Galveston to take over as the chief port city. The 1900 Galveston hurricane subsequently devastated that city, killing about 8,000 people or possibly as many as 12,000. This makes it the deadliest natural disaster in U.S. history.[47] In 2017, Hurricane Harvey made landfall in Rockport as a Category 4 Hurricane, causing significant damage there. The storm stalled over land for a very long time, allowing it to drop unprecedented amounts of rain over the Greater Houston area and surrounding counties. The result was widespread and catastrophic flooding that inundated hundreds of thousands of homes. Harvey ultimately became the costliest hurricane worldwide, causing an estimated $198.6 billion in damage, surpassing the cost of Hurricane Katrina.[48]
|
66 |
+
|
67 |
+
Other devastating Texas hurricanes include the 1915 Galveston hurricane, Hurricane Audrey in 1957 which killed more than 600 people, Hurricane Carla in 1961, Hurricane Beulah in 1967, Hurricane Alicia in 1983, Hurricane Rita in 2005, and Hurricane Ike in 2008. Tropical storms have also caused their share of damage: Allison in 1989 and again during 2001, and Claudette in 1979 among them.
|
68 |
+
|
69 |
+
As of 2017[update] Texas emits the most greenhouse gases in the U.S, almost twice the amount of California, the second most polluting state.[49]
|
70 |
+
As of 2017[update] the state emits about 1,600 billion pounds (707 million metric tons) of carbon dioxide annually.[49] As an independent nation, Texas would rank as the world's seventh-largest producer of greenhouse gases.[50] Causes of the state's vast greenhouse gas emissions include the state's large number of coal power plants and the state's refining and manufacturing industries.[50] In 2010, there were 2,553 "emission events" which poured 44.6 million pounds (20,200 metric tons) of contaminants into the Texas sky.[51]
|
71 |
+
|
72 |
+
Texas lies between two major cultural spheres of Pre-Columbian North America: the Southwestern and the Plains areas. Archaeologists have found that three major indigenous cultures lived in this territory, and reached their developmental peak before the first European contact. These were:[52]
|
73 |
+
|
74 |
+
When Europeans arrived in the Texas region, there were several races of Native peoples divided into many smaller tribes. They were Caddoan, Atakapan, Athabaskan, Coahuiltecan, and Uto-Aztecan. The Uto-Aztecan Puebloan peoples lived neared the Rio Grande in the western portion of the state, the Athabaskan-speaking Apache tribes lived throughout the interior, the Caddoans controlled much of the Red River region and the Atakapans were mostly centered along the Gulf Coast. At least one tribe of Coahuiltecans, the Aranama, lived in southern Texas. This entire culture group, primarily centered in northeastern Mexico, is now extinct. It is difficult to say who lived in the northwestern region of the state originally. By the time the region came to be explored, it belonged to the fairly well-known Comanche, another Uto-Aztecan people who had transitioned into a powerful horse culture, but it is believed that they came later and did not live there during the 16th century. It may have been claimed by several different peoples, including Uto-Aztecans, Athabaskans, or even Dhegihan Siouans.
|
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+
|
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+
No culture was dominant in the present-day Texas region, and many peoples inhabited the area.[52] Native American tribes who lived inside the boundaries of present-day Texas include the Alabama, Apache, Atakapan, Bidai, Caddo, Aranama, Comanche, Choctaw, Coushatta, Hasinai, Jumano, Karankawa, Kickapoo, Kiowa, Tonkawa, and Wichita.[53][54] The name Texas derives from táyshaʔ, a word in the Caddoan language of the Hasinai, which means "friends" or "allies".[1][55][56][57][58]
|
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+
|
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+
The region was primarily controlled by the Spanish for the first couple centuries of contact, until the Texas Revolution. They were not particularly kind to their native populations—even less so with the Caddoans, who were not trusted as their culture was split between the Spanish and the French. When the Spanish briefly managed to conquer the Louisiana colony, they decided to switch tactics and attempt being exceedingly friendly to the Indians, which they continued even after the French took back the colony. After the 1803 Louisiana Purchase, the United States inherited this odd circumstance. The Caddoans preferred the company of Americans[according to whom?] and almost the entire population of them migrated into the states of Louisiana and Arkansas. The Spanish felt jilted after having spent so much time and effort and began trying to lure the Caddo back, even promising them more land. Seemingly without actually knowing how they came by it,[according to whom?] the United States (who had begun convincing tribes to self-segregate from whites by selling everything and moving west ever since they gained the Louisiana Purchase) faced an overflow of native peoples in Missouri and Arkansas and were able to negotiate with the Caddoans to allow several displaced peoples to settle on unused lands in eastern Texas. They included the Muscogee, Houma Choctaw, Lenape and Mingo Seneca, among others, who all came to view the Caddoans as saviors, making those peoples highly influential.[59][60]
|
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|
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Whether a Native American tribe was friendly or warlike was critical to the fates of European explorers and settlers in that land.[61] Friendly tribes taught newcomers how to grow indigenous crops, prepare foods, and hunt wild game. Warlike tribes made life difficult and dangerous for Europeans through their attacks and resistance to the newcomers.[62]
|
81 |
+
|
82 |
+
During the Texas Revolution, the U.S. became heavily involved. Prior treaties with the Spanish forbade either side from militarizing its native population in any potential conflict between the two nations. At that time, several sudden outbreaks of violence between Caddoans and Texans started to spread. The Caddoans were always clueless[speculation?] when questioned, The Texan and American authorities in the region could never find hard evidence linking them to it and often it was so far-flung from Caddoan lands, it barely made any sense. It seems most likely that these were false-flag attacks meant to start a cascading effect to force the natives under Caddoan influence into armed conflict without breaking any treaties—preferably on the side of the Spanish. While no proof was found as to who the culprit was, those in charge of Texas at the time attempted multiple times to publicly blame and punish the Caddoans for the incidents with the U.S. government trying to keep them in check. Furthermore, the Caddoans never turned to violence because of it, excepting cases of self-defense.[59]
|
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+
|
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+
By the 1830s, the U.S. had drafted the Indian Removal Act, which was used to facilitate the Trail of Tears. Fearing retribution of other native peoples, Indian Agents all over the eastern U.S. began desperately trying to convince all their native peoples to uproot and move west. This included the Caddoans of Louisiana and Arkansas. Following the Texas Revolution, the Texans chose to make peace with their Native peoples but did not honor former land claims or agreements. This began the movement of Native populations north into what would become Indian Territory—modern-day Oklahoma.[59]
|
85 |
+
|
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+
The first historical document related to Texas was a map of the Gulf Coast, created in 1519 by Spanish explorer Alonso Álvarez de Pineda.[63] Nine years later, shipwrecked Spanish explorer Álvar Núñez Cabeza de Vaca and his cohort became the first Europeans in what is now Texas.[64][65] Cabeza de Vaca reported that in 1528, when the Spanish landed in the area, "half the natives died from a disease of the bowels and blamed us."[66] Cabeza de Vaca also made observations about the way of life of the Ignaces Natives of Texas:
|
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+
|
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+
They went about with a firebrand, setting fire to the plains and timber so as to drive off the mosquitos, and also to get lizards and similar things which they eat, to come out of the soil. In the same manner they kill deer, encircling them with fires, and they do it also to deprive the animals of pasture, compelling them to go for food where the Indians want.[67]
|
89 |
+
|
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+
Francisco Vázquez de Coronado describes his 1541 encounter:
|
91 |
+
|
92 |
+
Two kinds of people travel around these plains with the cows; one is called Querechos and the others Teyas; they are very well built, and painted, and are enemies of each other. They have no other settlement or location than comes from traveling around with the cows. They kill all of these they wish and tan the hides, with which they clothe themselves and make their tents, and they eat the flesh, sometimes even raw, and they also even drink the blood when thirsty. The tents they make are like field tents, and they set them up over poles they have made for this purpose, which come together and are tied at the top, and when they go from one place to another they carry them on some dogs they have, of which they have many, and they load them with the tents and poles and other things, for the country is so level, as I said, that they can make use of these, because they carry the poles dragging along on the ground. The sun is what they worship most.[68]
|
93 |
+
|
94 |
+
European powers ignored the area until accidentally settling there in 1685. Miscalculations by René-Robert Cavelier de La Salle resulted in his establishing the colony of Fort Saint Louis at Matagorda Bay rather than along the Mississippi River.[69] The colony lasted only four years before succumbing to harsh conditions and hostile natives.[70]
|
95 |
+
|
96 |
+
In 1690 Spanish authorities, concerned that France posed a competitive threat, constructed several missions in East Texas.[71] After Native American resistance, the Spanish missionaries returned to Mexico.[72] When France began settling Louisiana, mostly in the southern part of the state, in 1716 Spanish authorities responded by founding a new series of missions in East Texas.[73] Two years later, they created San Antonio as the first Spanish civilian settlement in the area.[74]
|
97 |
+
|
98 |
+
Hostile native tribes and distance from nearby Spanish colonies discouraged settlers from moving to the area. It was one of New Spain's least populated provinces.[76] In 1749, the Spanish peace treaty with the Lipan Apache[77] angered many tribes, including the Comanche, Tonkawa, and Hasinai.[78] The Comanche signed a treaty with Spain in 1785[79] and later helped to defeat the Lipan Apache and Karankawa tribes.[80] With more numerous missions being established, priests led a peaceful conversion of most tribes. By the end of the 18th century only a few nomadic tribes had not converted to Christianity.[81]
|
99 |
+
|
100 |
+
When the United States purchased Louisiana from France in 1803, American authorities insisted the agreement also included Texas. The boundary between New Spain and the United States was finally set at the Sabine River in 1819, at what is now the border between Texas and Louisiana.[82] Eager for new land, many United States settlers refused to recognize the agreement. Several filibusters raised armies to invade the area west of the Sabine River.[83] In 1821, the Mexican War of Independence included the Texas territory, which became part of Mexico.[84] Due to its low population, Mexico made the area part of the state of Coahuila y Tejas.[85]
|
101 |
+
|
102 |
+
Hoping more settlers would reduce the near-constant Comanche raids, Mexican Texas liberalized its immigration policies to permit immigrants from outside Mexico and Spain.[86] Under the Mexican immigration system, large swathes of land were allotted to empresarios, who recruited settlers from the United States, Europe, and the Mexican interior. The first grant, to Moses Austin, was passed to his son Stephen F. Austin after his death.
|
103 |
+
|
104 |
+
Austin's settlers, the Old Three Hundred, made places along the Brazos River in 1822.[87] Twenty-three other empresarios brought settlers to the state, the majority of whom were from the United States.[88] The population of Texas grew rapidly. In 1825, Texas had about 3,500 people, with most of Mexican descent.[89] By 1834, the population had grown to about 37,800 people, with only 7,800 of Mexican descent.[90] Most of these early settlers who arrived with Austin and soon after were persons less than fortunate in life, as Texas was devoid of the comforts found elsewhere in Mexico and the United States during that time. Early Texas settler David B. Edwards described his fellow Texans as being "banished from the pleasures of life".[91]
|
105 |
+
|
106 |
+
Many immigrants openly flouted Mexican law, especially the prohibition against slavery. Combined with United States' attempts to purchase Texas, Mexican authorities decided in 1830 to prohibit continued immigration from the United States.[92] Illegal immigration from the United States into Mexico continued to increase the population of Texas anyway.[93] New laws also called for the enforcement of customs duties angering native Mexican citizens (Tejanos) and recent immigrants alike.[94]
|
107 |
+
|
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+
The Anahuac Disturbances in 1832 were the first open revolt against Mexican rule, and they coincided with a revolt in Mexico against the nation's president.[95] Texians sided with the federalists against the current government and drove all Mexican soldiers out of East Texas.[96] They took advantage of the lack of oversight to agitate for more political freedom. Texians met at the Convention of 1832 to discuss requesting independent statehood, among other issues.[97] The following year, Texians reiterated their demands at the Convention of 1833.[98]
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Within Mexico, tensions continued between federalists and centralists. In early 1835, wary Texians formed Committees of Correspondence and Safety.[99] The unrest erupted into armed conflict in late 1835 at the Battle of Gonzales.[100] This launched the Texas Revolution, and over the next two months the Texians defeated all Mexican troops in the region.[101] Texians elected delegates to the Consultation, which created a provisional government.[102] The provisional government soon collapsed from infighting, and Texas was without clear governance for the first two months of 1836.[103]
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During this time of political turmoil, Mexican President Antonio López de Santa Anna personally led an army to end the revolt.[104] The Mexican expedition was initially successful. General José de Urrea defeated all the Texian resistance along the coast culminating in the Goliad massacre.[105] Santa Anna's forces, after a thirteen-day siege, overwhelmed Texian defenders at the Battle of the Alamo. News of the defeats sparked panic among Texas settlers.[106]
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The newly elected Texian delegates to the Convention of 1836 quickly signed a Declaration of Independence on March 2, forming the Republic of Texas. After electing interim officers, the Convention disbanded.[107] The new government joined the other settlers in Texas in the Runaway Scrape, fleeing from the approaching Mexican army.[106] After several weeks of retreat, the Texian Army commanded by Sam Houston attacked and defeated Santa Anna's forces at the Battle of San Jacinto.[108] Santa Anna was captured and forced to sign the Treaties of Velasco, ending the war.[109] The Constitution of the Republic of Texas prohibited the government from restricting slavery or freeing slaves, required free people of African descent to leave the country, and prohibited Native Americans from becoming citizens.
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While Texas had won its independence, political battles raged between two factions of the new Republic. The nationalist faction, led by Mirabeau B. Lamar, advocated the continued independence of Texas, the expulsion of the Native Americans, and the expansion of the Republic to the Pacific Ocean. Their opponents, led by Sam Houston, advocated the annexation of Texas to the United States and peaceful co-existence with Native Americans. The conflict between the factions was typified by an incident known as the Texas Archive War.[110] With wide popular support, Texas first applied for annexation to the United States in 1836, but its status as a slaveholding country caused its admission to be controversial and it was initially rebuffed. This status, and Mexican diplomacy in support of its claims to the territory, also complicated Texas's ability to form foreign alliances and trade relationships.[111]
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Mexico launched two small expeditions into Texas in 1842. The town of San Antonio was captured twice and Texans were defeated in battle in the Dawson massacre. Despite these successes, Mexico did not keep an occupying force in Texas, and the republic survived.[112] The cotton price crash of the 1840s depressed the country's economy,[111] contributing to the republic's inability to defend itself, and adding momentum to Texas's eventual annexation into the United States.
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As early as 1837, the Republic made several attempts to negotiate annexation with the United States.[113] Opposition within the republic from the nationalist faction, along with strong abolitionist opposition within the United States, slowed Texas's admission into the Union. Texas was finally annexed when the expansionist James K. Polk won the election of 1844.[114] On December 29, 1845, Congress admitted Texas to the U.S. as a constituent state of the Union.[115]
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The population of the new state was quite small at first, and there was a strong mix between the English-speaking American settlers who dominated in the state's eastern/northeastern portions and the Spanish-speaking former Mexicans who dominated in the state's southern and western portions. Statehood brought many new settlers. Because of the long Spanish presence in Mexico and various failed colonization efforts by the Spanish and Mexicans in northern Mexico, there were large herds of Longhorn cattle that roamed the state. Hardy by nature, but also suitable for slaughtering and consumption, they represented an economic opportunity many entrepreneurs seized upon, thus creating the cowboy culture for which Texas is famous. While in the early days of the republic cattle and bison were slaughtered for their hides, soon a beef industry was established with cattle being shipped all over the U.S. and the Caribbean (within a few decades, beef had become a staple of the American diet).
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After Texas's annexation, Mexico broke diplomatic relations with the United States. While the United States claimed Texas's border stretched to the Rio Grande, Mexico claimed it was the Nueces River leaving the Rio Grande Valley under contested Texan sovereignty.[115] While the former Republic of Texas could not enforce its border claims, the United States had the military strength and the political will to do so. President Polk ordered General Zachary Taylor south to the Rio Grande on January 13, 1846. A few months later Mexican troops routed an American cavalry patrol in the disputed area in the Thornton Affair starting the Mexican–American War. The first battles of the war were fought in Texas: the Siege of Fort Texas, Battle of Palo Alto and Battle of Resaca de la Palma. After these decisive victories, the United States invaded Mexican territory, ending the fighting in Texas.[116]
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After a series of United States victories, the Treaty of Guadalupe Hidalgo ended the two-year war. In return, for US$18,250,000, Mexico gave the U.S. undisputed control of Texas, ceded the Mexican Cession in 1848, most of which today is called the American Southwest, and Texas's borders were established at the Rio Grande.[116]
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The Compromise of 1850 set Texas's boundaries at their present form. U.S. Senator James Pearce of Maryland drafted the final proposal[27] where Texas ceded its claims to land which later became half of present-day New Mexico, a third of Colorado, and small portions of Kansas, Oklahoma, and Wyoming to the federal government, in return for the assumption of $10 million of the old republic's debt.[27] Post-war Texas grew rapidly as migrants poured into the cotton lands of the state.[117]
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They also brought or purchased enslaved African Americans, whose numbers tripled in the state from 1850 to 1860, from 58,000 to 182,566.[118]
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Texas was at war again after the election of 1860. At this time, blacks comprised 30 percent of the state's population, and they were overwhelmingly enslaved.[119] When Abraham Lincoln was elected, South Carolina seceded from the Union. Five other Lower South states quickly followed. A State Convention considering secession opened in Austin on January 28, 1861. On February 1, by a vote of 166–8, the Convention adopted an Ordinance of Secession from the United States. Texas voters approved this Ordinance on February 23, 1861. Texas joined the newly created Confederate States of America on March 4, 1861 ratifying the permanent C.S. Constitution on March 23.[1][120]
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Not all Texans favored secession initially, although many of the same would later support the Southern cause. Texas's most notable Unionist was the state Governor, Sam Houston. Not wanting to aggravate the situation, Houston refused two offers from President Lincoln for Union troops to keep him in office. After refusing to swear an oath of allegiance to the Confederacy, Houston was deposed as governor.[121]
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While far from the major battlefields of the American Civil War, Texas contributed large numbers of men and equipment to the rest of the Confederacy.[122] Union troops briefly occupied the state's primary port, Galveston. Texas's border with Mexico was known as the "backdoor of the Confederacy" because trade occurred at the border, bypassing the Union blockade.[123] The Confederacy repulsed all Union attempts to shut down this route,[122] but Texas's role as a supply state was marginalized in mid-1863 after the Union capture of the Mississippi River. The final battle of the Civil War was fought near Brownsville, Texas at Palmito Ranch[124] with a Confederate victory.
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Texas descended into anarchy for two months between the surrender of the Army of Northern Virginia and the assumption of authority by Union General Gordon Granger. Violence marked the early months of Reconstruction.[122] Juneteenth commemorates the announcement of the Emancipation Proclamation in Galveston by General Gordon Granger, almost two and a half years after the original announcement.[125][126] President Johnson, in 1866, declared the civilian government restored in Texas.[127] Despite not meeting reconstruction requirements, Congress resumed allowing elected Texas representatives into the federal government in 1870. Social volatility continued as the state struggled with agricultural depression and labor issues.[128]
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Like most of the South, the Texas economy was devastated by the War. However, since the state had not been as dependent on slaves as other parts of the South, it was able to recover more quickly. The culture in Texas during the later 19th century exhibited many facets of a frontier territory. The state became notorious as a haven for people from other parts of the country who wanted to escape debt, criminal prosecution, or other problems. Indeed, "Gone to Texas" was a common expression for those fleeing the law in other states. Nevertheless, the state also attracted many businessmen and other settlers with more legitimate interests as well.
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The cattle industry continued to thrive, though it gradually became less profitable. Cotton and lumber became major industries creating new economic booms in various regions of the state. Railroad networks grew rapidly as did the port at Galveston as commerce between Texas and the rest of the U.S. (and the rest of the world) expanded. As with some other states before, the lumber industry quickly decimated the forests of Texas such that, by the early 20th century, the majority of the forest population in Texas was gone (later conservation efforts restored some of it, but never to the level it once was).
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In 1900, Texas suffered the deadliest natural disaster in U.S. history during the Galveston hurricane.[47] On January 10, 1901, the first major oil well in Texas, Spindletop, was found south of Beaumont. Other fields were later discovered nearby in East Texas, West Texas, and under the Gulf of Mexico. The resulting "oil boom" transformed Texas.[129] Oil production eventually averaged three million barrels per day at its peak in 1972.[130]
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In 1901, the Democratic-dominated state legislature passed a bill requiring payment of a poll tax for voting, which effectively disenfranchised most blacks and many poor whites and Latinos. In addition, the legislature established white primaries, ensuring minorities were excluded from the formal political process. The number of voters dropped dramatically, and the Democrats crushed competition from the Republican and Populist parties.[131][132] The Socialist Party became the second-largest party in Texas after 1912,[133] coinciding with a large socialist upsurge in the United States during fierce battles in the labor movement and the popularity of national heroes like Eugene V. Debs. The Socialists' popularity soon waned after their vilification by the United States government for their opposition to US involvement in World War I.
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The Great Depression and the Dust Bowl dealt a double blow to the state's economy, which had significantly improved since the Civil War. Migrants abandoned the worst-hit sections of Texas during the Dust Bowl years. Especially from this period on, blacks left Texas in the Great Migration to get work in the Northern United States or California and to escape the oppression of segregation.[119] In 1940, Texas was 74 percent Anglo, 14.4 percent black, and 11.5 percent Hispanic.[134]
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World War II had a dramatic impact on Texas, as federal money poured in to build military bases, munitions factories, POW detention camps and Army hospitals; 750,000 young men left for service; the cities exploded with new industry; the colleges took on new roles; and hundreds of thousands of poor farmers left the fields for much better-paying war jobs, never to return to agriculture.[135][136] Texas manufactured 3.1 percent of total United States military armaments produced during World War II, ranking eleventh among the 48 states.[137]
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Texas modernized and expanded its system of higher education through the 1960s. The state created a comprehensive plan for higher education, funded in large part by oil revenues, and a central state apparatus designed to manage state institutions more efficiently. These changes helped Texas universities receive federal research funds.[138]
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On November 22, 1963, President John F. Kennedy was assassinated in Dallas.[139]
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Beginning around the mid-20th century, Texas began to transform from a rural and agricultural state to one urban and industrialized.[140] The state's population grew quickly during this period, with large levels of migration from outside the state.[140] As a part of the Sun Belt, Texas experienced strong economic growth, particularly during the 1970s and early 1980s.[140] Texas's economy diversified, lessening its reliance on the petroleum industry.[140] By 1990, Hispanics overtook blacks to become the largest minority group in the state.[140]
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During the late 20th century, the Republican Party replaced the Democratic Party as the dominant party in the state, as the latter became more politically liberal and as demographic changes favored the former.[140]
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The current Texas Constitution was adopted in 1876. Like many states, it explicitly provides for a separation of powers. The state's Bill of Rights is much larger than its federal counterpart, and has provisions unique to Texas.[141]
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Texas has a plural executive branch system limiting the power of the governor, which is a weak executive compared to some other states. Except for the Secretary of State, voters elect executive officers independently; thus candidates are directly answerable to the public, not the governor.[142] This election system has led to some executive branches split between parties and reduced the ability of the governor to carry out a program. When Republican President George W. Bush served as Texas's governor, the state had a Democratic lieutenant governor, Bob Bullock. The executive branch positions consist of the Governor, Lieutenant Governor, Comptroller of Public Accounts, Land Commissioner, Attorney General, Agriculture Commissioner, the three-member Texas Railroad Commission, the State Board of Education, and the Secretary of State.[142]
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The bicameral Texas Legislature consists of the House of Representatives, with 150 members, and a Senate, with 31 members. The Speaker of the House leads the House, and the lieutenant governor, the Senate.[143] The Legislature meets in regular session biennially for just over a hundred days, but the governor can call for special sessions as often as desired (notably, the Legislature cannot call itself into session).[144] The state's fiscal year begins September 1.
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The judiciary of Texas is one of the most complex in the United States, with many layers and overlapping jurisdictions. Texas has two courts of last resort: the Texas Supreme Court, for civil cases, and the Texas Court of Criminal Appeals. Except for some municipal benches, partisan elections select judges at all levels of the judiciary; the governor fills vacancies by appointment.[145] Texas is notable for its use of capital punishment, having led the country in executions since capital punishment was reinstated in the Gregg v. Georgia case (see Capital punishment in Texas).
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The Texas Ranger Division of the Texas Department of Public Safety is a law enforcement agency with statewide jurisdiction. Over the years, the Texas Rangers have investigated crimes ranging from murder to political corruption. They have acted as riot police and as detectives, protected the Texas governor, tracked down fugitives, and functioned as a paramilitary force both for the republic and for the state. The Texas Rangers were unofficially created by Stephen F. Austin in 1823 and formally constituted in 1835. The Rangers were integral to several important events of Texas history and some of the best-known criminal cases in the history of the Old West.[146]
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The Texas constitution defines the responsibilities of county governments, which serve as agents of the state. What are called commissioners court and court judges are elected to serve as the administrative arm. Most cities in the state, those over 5,000 in population, have home-rule governments. The vast majority of these have charters for council-manager forms of government, by which voters elect council members, who hire a professional city manager as an operating officer.
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In the 1870s, white Democrats wrested power back in the state legislature from the biracial coalition at the end of Reconstruction. In the early 20th century, the legislature passed bills to impose poll taxes, followed by white primaries; these measures effectively disfranchised most blacks, poor whites and Mexican Americans.[131][132] In the 1890s, 100,000 blacks voted in the state; by 1906, only 5,000 could vote.[147] As a result, the Democratic Party dominated Texas politics from the turn of the century, imposing racial segregation and white supremacy. It held power until after passage in the mid-1960s of national civil rights legislation enforcing constitutional rights of all citizens.
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Although Texas was essentially a one-party state during this time and the Democratic primary was viewed as "the real election", the Democratic Party had conservative and liberal factions, which became more pronounced after the New Deal.[148] Additionally, several factions of the party briefly split during the 1930s and 1940s.[148]
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The state's conservative white voters began to support Republican presidential candidates by the mid-20th century. After this period, they supported Republicans for local and state offices as well, and most whites became Republican Party members.[149] The party also attracted some minorities, but many have continued to vote for Democratic candidates. The shift to the Republican Party is much-attributed to the fact the Democratic Party became increasingly liberal during the 20th century, and thus increasingly out-of-touch with the average Texas voter.[150] As Texas was always a conservative state, voters switched to the GOP, which now more closely reflected their beliefs.[150][151] Commentators have also attributed the shift to Republican political consultant Karl Rove, who managed numerous political campaigns in Texas in the 1980s and 1990s.[151] Other stated reasons included court-ordered redistricting and the demographic shift in relation to the Sun Belt that favored the Republican Party and conservatism.[140]
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The 2003 Texas redistricting of Congressional districts led by Republican Tom DeLay, was called by The New York Times "an extreme case of partisan gerrymandering".[152] A group of Democratic legislators, the "Texas Eleven", fled the state in a quorum-busting effort to prevent the legislature from acting, but was unsuccessful.[153] The state had already redistricted following the 2000 census. Despite these efforts, the legislature passed a map heavily in favor of Republicans, based on 2000 data and ignoring the estimated nearly one million new residents in the state since then. Career attorneys and analysts at the Department of Justice objected to the plan as diluting the votes of African American and Hispanic voters, but political appointees overrode them and approved it.[152] Legal challenges to the redistricting reached the national Supreme Court in the case League of United Latin American Citizens v. Perry (2006), but the court ruled in favor of the state (and Republicans).[154]
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In the 2014 Texas elections, the Tea Party movement made large gains, with numerous Tea Party favorites being elected into office, including Dan Patrick as lieutenant governor,[155][156] Ken Paxton as attorney general,[155][157] in addition to numerous other candidates[157] including conservative Republican Greg Abbott as governor.[158]
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Texas voters lean toward fiscal conservatism, while enjoying the benefits of huge federal investment in the state in military and other facilities achieved by the power of the Solid South in the 20th century. They also tend to have socially conservative values.[159][160]
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Since 1980, most Texas voters have supported Republican presidential candidates. In 2000 and 2004, Republican George W. Bush won Texas with respectively 59.3 and 60.1 percent of the vote, partly due to his "favorite son" status as a former governor of the state. John McCain won the state in 2008, but with a smaller margin of victory compared to Bush at 55 percent of the vote. Austin, Dallas, Houston, and San Antonio consistently lean Democratic in both local and statewide elections.
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The state's changing demographics may result in a change in its overall political alignment, as a majority population of Black and Hispanic/Latino voters support the Democratic Party.[161] Residents of counties along the Rio Grande closer to the Mexico–United States border, where there are many Latino residents, generally vote for Democratic Party candidates, while most other rural and suburban areas of Texas have shifted to voting for Republican Party candidates.[162][163]
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As of the general elections of 2014, a large majority of the members of Texas's U.S. House delegation are Republican, along with both U.S. Senators. In the 114th United States Congress, of the 36 Congressional districts in Texas, 24 are held by Republicans and 11 by Democrats. One seat is vacant. Texas's Senators are John Cornyn and Ted Cruz. Since 1994, Texans have not elected a Democrat to a statewide office. The state's Democratic voters are made up primarily by liberal and minority groups in Austin, Beaumont, Dallas, El Paso, Houston, and San Antonio as well as minority voters in East and South Texas.
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Texas has banned sanctuary cities,[164] but Houston Mayor Sylvester Turner has vowed that the city will not assist ICE agents.[165]
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Texas has 254 counties—the most nationwide. Each county runs on Commissioners' Court system consisting of four elected commissioners (one from each of four precincts in the county, roughly divided according to population) and a county judge elected at large from the entire county. County government runs similar to a "weak" mayor-council system; the county judge has no veto authority, but votes along with the other commissioners.
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Although Texas permits cities and counties to enter "interlocal agreements" to share services, the state does not allow consolidated city-county governments, nor does it have metropolitan governments. Counties are not granted home rule status; their powers are strictly defined by state law. The state does not have townships—areas within a county are either incorporated or unincorporated. Incorporated areas are part of a municipality. The county provides limited services to unincorporated areas and to some smaller incorporated areas. Municipalities are classified either "general law" cities or "home rule".[167] A municipality may elect home rule status once it exceeds 5,000 population with voter approval.
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Texas also permits the creation of "special districts", which provide limited services. The most common is the school district, but can also include hospital districts, community college districts, and utility districts (one utility district near Austin was the plaintiff in a landmark Supreme Court case involving the Voting Rights Act).
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Municipal, school district, and special district elections are nonpartisan,[168] though the party affiliation of a candidate may be well-known. County and state elections are partisan.
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Texas has a reputation of very harsh criminal punishment for criminal offenses. It is one of the 32 states that practice capital punishment, and since the US Supreme Court allowed capital punishment to resume in 1976, 40% of all US executions have taken place in Texas.[169] As of 2008, Texas had the 4th highest incarceration rate in the US.[170] Texas also has strong self defense laws, allowing citizens to use lethal force to defend themselves, their families, or their property.[171]
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As of 2019, Texas had a gross state product (GSP) of $1.9 trillion, the second highest in the U.S.[172] Its GSP is greater than the GDPs of Brazil, Canada, Russia, South Korea and Spain, which are the world's 9th-, 10th-, 11th-, 12th- and 13th-largest economies, respectively.[173] Texas's economy is the second-largest of any country subdivision globally, behind California. Its per capita personal income in 2009 was $36,484, ranking 29th in the nation.[174]
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Texas's large population, an abundance of natural resources, thriving cities and leading centers of higher education have contributed to a large and diverse economy. Since oil was discovered, the state's economy has reflected the state of the petroleum industry. In recent times, urban centers of the state have increased in size, containing two-thirds of the population in 2005. The state's economic growth has led to urban sprawl and its associated symptoms.[174]
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As of May 2020, during the COVID-19 pandemic, the state's unemployment rate was 13 percent.[175]
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In 2010, Site Selection Magazine ranked Texas as the most business-friendly state in the nation, in part because of the state's three-billion-dollar Texas Enterprise Fund.[176] Texas has the joint-highest number of Fortune 500 company headquarters in the United States, along with California.[177][178]
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In 2010, there were 346,000 millionaires in Texas, constituting the second-largest population of millionaires in the nation.[179][180]
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Texas has a "low taxes, low services" reputation.[159] According to the Tax Foundation, Texans' state and local tax burdens rank among the lowest in the nation, 7th lowest nationally; state and local taxes cost $3,580 per capita, or 8.4 percent of resident incomes.[181] Texas is one of seven states that lack a state income tax.[181][182]
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Instead, the state collects revenue from property taxes (though these are collected at the county, city, and school district level; Texas has a state constitutional prohibition against a state property tax) and sales taxes. The state sales tax rate is 6.25 percent,[181][183] but local taxing jurisdictions (cities, counties, special purpose districts, and transit authorities) may also impose sales and use tax up to 2 percent for a total maximum combined rate of 8.25 percent.[184]
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Texas is a "tax donor state"; in 2005, for every dollar Texans paid to the federal government in federal income taxes, the state got back about $0.94 in benefits.[181] To attract business, Texas has incentive programs worth $19 billion per year (2012); more than any other US state.[185][186]
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Texas has the most farms and the highest acreage in the United States. The state is ranked No. 1 for revenue generated from total livestock and livestock products. It is ranked No. 2 for total agricultural revenue, behind California.[187] At $7.4 billion or 56.7 percent of Texas's annual agricultural cash receipts, beef cattle production represents the largest single segment of Texas agriculture. This is followed by cotton at $1.9 billion (14.6 percent), greenhouse/nursery at $1.5 billion (11.4 percent), broilers at $1.3 billion (10 percent), and dairy products at $947 million (7.3 percent).[188]
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Texas leads the nation in the production of cattle, horses, sheep, goats, wool, mohair and hay.[188] The state also leads the nation in production of cotton[187][189] which is the number one crop grown in the state in terms of value.[190] The state grows significant amounts of cereal crops and produce.[187] Texas has a large commercial fishing industry. With mineral resources, Texas leads in creating cement, crushed stone, lime, salt, sand and gravel.[187]
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Texas throughout the 21st century has been hammered by drought. This has cost the state billions of dollars in livestock and crops.[191]
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Ever since the discovery of oil at Spindletop, energy has been a dominant force politically and economically within the state.[192] If Texas were its own country it would be the sixth largest oil producer in the world.[193]
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The Railroad Commission of Texas, contrary to its name, regulates the state's oil and gas industry, gas utilities, pipeline safety, safety in the liquefied petroleum gas industry, and surface coal and uranium mining. Until the 1970s, the commission controlled the price of petroleum because of its ability to regulate Texas's oil reserves. The founders of the Organization of Petroleum Exporting Countries (OPEC) used the Texas agency as one of their models for petroleum price control.[194]
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Texas has known petroleum deposits of about 5 billion barrels (790,000,000 m3), which makes up about one-fourth of the known U.S. reserves.[195] The state's refineries can process 4.6 million barrels (730,000 m3) of oil a day.[195] The Port Arthur Refinery in Southeast Texas is the largest refinery in the U.S.[195] Texas also leads in natural gas production, producing one-fourth of the nation's supply.[195] Several petroleum companies are based in Texas such as: Occidental Petroleum, ConocoPhillips, ExxonMobil, Halliburton, Marathon Oil, Tesoro, Valero Energy, and Western Refining.
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According to the Energy Information Administration, Texans consume, on average, the fifth most energy (of all types) in the nation per capita and as a whole, following behind Wyoming, Alaska, Louisiana, North Dakota, and Iowa.[195]
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Unlike the rest of the nation, most of Texas is on its own alternating current power grid, the Texas Interconnection. Texas has a deregulated electric service. Texas leads the nation in total net electricity production, generating 437,236 MWh in 2014, 89% more MWh than Florida, which ranked second.[196][197] As an independent nation, Texas would rank as the world's eleventh-largest producer of electricity, after South Korea, and ahead of the United Kingdom.
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The state is a leader in renewable energy commercialization; it produces the most wind power in the nation.[195][198] In 2014, 10.6% of the electricity consumed in Texas came from wind turbines.[199] The Roscoe Wind Farm in Roscoe, Texas, is one of the world's largest wind farms with a 781.5 megawatt (MW) capacity.[200] The Energy Information Administration states the state's large agriculture and forestry industries could give Texas an enormous amount biomass for use in biofuels. The state also has the highest solar power potential for development in the nation.[195]
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With large universities systems coupled with initiatives like the Texas Enterprise Fund and the Texas Emerging Technology Fund, a wide array of different high tech industries have developed in Texas. The Austin area is nicknamed the "Silicon Hills" and the north Dallas area the "Silicon Prairie". Many high-tech companies are located in or have their headquarters in Texas (and Austin in particular), including Dell, Inc., Borland, Forcepoint, Indeed.com, Texas Instruments, Perot Systems, Rackspace and AT&T.
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The National Aeronautics and Space Administration's Lyndon B. Johnson Space Center (NASA JSC) in Southeast Houston, sits as the crown jewel of Texas's aeronautics industry. Fort Worth hosts both Lockheed Martin's Aeronautics division and Bell Helicopter Textron.[201][202] Lockheed builds the F-16 Fighting Falcon, the largest Western fighter program, and its successor, the F-35 Lightning II in Fort Worth.[203]
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Texas's affluence stimulates a strong commercial sector consisting of retail, wholesale, banking and insurance, and construction industries. Examples of Fortune 500 companies not based on Texas traditional industries are AT&T, Kimberly-Clark, Blockbuster, J. C. Penney, Whole Foods Market, and Tenet Healthcare.[204]
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Nationally, the Dallas–Fort Worth area, home to the second shopping mall in the United States, has the most shopping malls per capita of any American metropolitan area.[205]
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Mexico, the state's largest trading partner, imports a third of the state's exports because of the North American Free Trade Agreement (NAFTA). NAFTA has encouraged the formation of controversial maquiladoras on the Texas–Mexico border.[206]
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The United States Census Bureau estimates the population of Texas was 28,995,881 on July 1, 2019, a 15.31 percent increase since the 2010 United States Census.[208][6]
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As of 2015, Texas had 4.7 million foreign-born residents, about 17% of the population and 21.6% of the state workforce.[209] The major countries of origin for Texan immigrants were Mexico (55.1% of immigrants), India (5%), El Salvador (4.3%), Vietnam (3.7%), and China (2.3%).[209] Of immigrant residents, some 35.8 percent were naturalized U.S. citizens.[209] In 2014, there were an estimated 1.7 million undocumented immigrants in Texas, making up 35% of the total Texas immigrant population and 6.1% of the total state population.[209] In addition to the state's foreign-born population, an additional 4.1 million Texans (15% of the state's population) were born in the United States and had at least one immigrant parent.[209]
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Texas's Rio Grande Valley has seen significant migration from across the U.S.–Mexico border. During the 2014 crisis, many Central Americans, including unaccompanied minors traveling alone from Guatemala, Honduras, and El Salvador, reached the state, overwhelming Border Patrol resources for a time. Many sought asylum in the United States.[210][211]
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Texas's population density is 90.5 people per square mile (34.9/km2) which is slightly higher than the average population density of the U.S. as a whole, at 80.6 people per square mile (31.1/km2). In contrast, while Texas and France are similarly sized geographically, the European country has a population density of 301.8 people per square mile (116.5/km2).
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Two-thirds of all Texans live in major metropolitan areas such as Houston. The Dallas-Fort Worth Metropolitan Area is the largest in Texas. While Houston is the largest city in Texas and the fourth-largest city in the United States, the Dallas-Fort Worth metropolitan area is larger than Houston.
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As of the 2015 Texas Population Estimate Program, the population of the state was 27,469,114; non-Hispanic whites 11,505,371 (41.9%); Black Americans 3,171,043 (11.5%); other races 1,793,580 (6.5%); and Hispanics and Latinos (of any race) 10,999,120 (40.0%).[212]
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According to the 2010 United States census, the racial composition of Texas was the following:[213]
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In addition, 37.6 percent of the population was Hispanic or Latino (of any race) (31.6 percent Mexican, 0.9 percent Salvadoran, 0.5 percent Puerto Rican, 0.4 percent Honduran, 0.3 percent Guatemalan 0.3 percent Spaniard, 0.2 percent Colombian, 0.2 percent Cuban).[214]
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As of 2011, 69.8% of the population of Texas younger than age 1 were minorities (meaning they had at least one parent who was not non-Hispanic white).[215]
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German, Irish, and English Americans are the three largest European ancestry groups in Texas. German Americans make up 11.3 percent of the population and number over 2.7 million members. Irish Americans make up 8.2 percent of the population and number over 1.9 million. There are roughly 600,000 French Americans and 472,000 Italian Americans residing in Texas; these two ethnic groups make up 2.5 percent and 2.0 percent of the population respectively. In the 1980 United States Census the largest ancestry group reported in Texas was English with 3,083,323 Texans citing they were of English or mostly English ancestry, making them 27 percent of the state at the time.[219] Their ancestry primarily goes back to the original thirteen colonies[citation needed] (the census of 1790 gives 48% of the population of English ancestry, together with 12% Scots and Scots-Irish, 4.5% other Irish, and 3% Welsh, for a total of 67.5% British and Irish; 13% were German, Swiss, Dutch, and French Huguenots; 19% were African-American),[220] thus many of them today identify as "American" in ancestry, though they are of predominantly British stock.[221][222] There are nearly 200,000 Czech Americans living in Texas, the largest number of any state.[223]
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African Americans are the largest racial minority in Texas. Their proportion of the population has declined since the early 20th century after many left the state in the Great Migration. Blacks of both Hispanic and non-Hispanic origin make up 11.5 percent of the population; blacks of non-Hispanic origin form 11.3 percent of the populace. African Americans of both Hispanic and non-Hispanic origin number at roughly 2.7 million individuals.
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Native Americans are a smaller minority in the state. Native Americans make up 0.5 percent of Texas's population and number over 118,000 individuals. Native Americans of non-Hispanic origin make up 0.3 percent of the population and number over 75,000 individuals. Cherokee made up 0.1 percent of the population, and numbered over 19,400. In contrast, only 583 identified as Chippewa.
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Asian Americans are a sizable minority group in Texas. Americans of Asian descent form 3.8 percent of the population, with those of non-Hispanic descent making up 3.7 percent of the populace. They total more than 808,000 individuals. Non-Hispanic Asians number over 795,000. Just over 200,000 Indian Americans make Texas their home. Texas is also home to more than 187,000 Vietnamese and 136,000 Chinese. In addition to 92,000 Filipinos and 62,000 Koreans, there are 18,000 Japanese Americans living in the state. Lastly, more than 111,000 people are of other Asian ancestry groups, such as Cambodian, Thai, and Hmong. Sugar Land, a city within the Houston metropolitan area, and Plano, within the Dallas metropolitan area, both have high concentrations of ethnic Chinese and Korean residents. The Houston and Dallas areas, and to a lesser extent, the Austin metropolitan area, all contain substantial Vietnamese communities.
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Americans with origins from the Pacific Islands are the smallest minority in Texas. According to the survey, only 18,000 Texans are Pacific Islanders; 16,400 are of non-Hispanic descent. There are roughly 5,400 Native Hawaiians, 5,300 Guamanians, and 6,400 people from other groups. Samoan Americans were scant; only 2,920 people were from this group. The city of Euless, a suburb of Fort Worth, contains a sizable population of Tongan Americans, at nearly 900 people, over one percent of the city's population. Killeen has a sufficient population of Samoans and Guamanian, and people of Pacific Islander descent surpass one percent of the city's population.
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Multiracial individuals are also a visible minority in Texas. People identifying as multiracial form 1.9 percent of the population, and number over 448,000 people. Almost 80,000 Texans claim African and European heritage and makeup 0.3 percent of the population. People of European and American Indian ancestry number over 108,800 (close to the number of Native Americans), and makeup 0.5 percent of the population. People of European and Asian ancestry number over 57,600, and form just 0.2 percent of the population. People of African and Native American ancestry were even smaller in number (15,300), and makeup just 0.1 percent of the total population.
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Hispanics and Latinos are the second-largest groups in Texas after non-Hispanic European Americans. More than 8.5 million people claim Hispanic or Latino ethnicity. This group forms over 37 percent of Texas's population. People of Mexican descent alone number over 7.9 million, and make up 31.6 percent of the population. The vast majority of the Hispanic/Latino population in the state is of Mexican descent, the next two largest groups are Salvadorans and Puerto Ricans. There are more than 222,000 Salvadorans and more than 130,000 Puerto Ricans in Texas. Other groups with large numbers in Texas include Hondurans, Guatemalans, Nicaraguans and Cubans, among others.[224][225] The Hispanics in Texas are more likely than in some other states (such as California) to identify as white; according to the 2010 U.S. Census, Texas is home to 6,304,207 White Hispanics and 2,594,206 Hispanics of "some other race" (usually mestizo).
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German descendants inhabit much of central and southeast-central Texas. Over one-third of Texas, residents are of Hispanic origin; while many have recently arrived, some Tejanos have ancestors with multi-generational ties to 18th century Texas. The African American population in Texas is increasing due to the New Great Migration.[226][227] In addition to the descendants of the state's former slave population, many African American college graduates have come to the state for work recently in the New Great Migration.[226] Recently, the Asian population in Texas has grown—primarily in Houston and Dallas. Other communities with a significantly growing Asian American population is in Austin, Corpus Christi, San Antonio, and the Sharyland area next McAllen, Texas. Three federally recognized Native American tribes reside in Texas: the Alabama-Coushatta Tribe, the Kickapoo Traditional Tribe, and the Ysleta del Sur Pueblo.[54]
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In 2010, 49 percent of all births were Hispanics; 35 percent were non-Hispanic whites; 11.5 percent were non-Hispanic blacks, and 4.3 percent were Asians/Pacific Islanders.[228] Based on Census Bureau data released in February 2011, for the first time in recent history, Texas's white population is below 50 percent (45 percent) and Hispanics grew to 38 percent. Between 2000 and 2010, the total population growth by 20.6 percent, but Hispanics growth by 65 percent, whereas non-Hispanic whites grew by only 4.2 percent.[229] Texas has the fifth highest rate of teenage births in the nation and a plurality of these are to Hispanics.[230]
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The state has three cities with populations exceeding one million: Houston, San Antonio, and Dallas.[235] These three rank among the 10 most populous cities of the United States. As of 2010, six Texas cities had populations greater than 600,000 people. Austin, Fort Worth, and El Paso are among the 20 largest U.S. cities. Texas has four metropolitan areas with populations greater than a million: Dallas–Fort Worth–Arlington, Houston–Sugar Land–Baytown, San Antonio–New Braunfels, and Austin–Round Rock–San Marcos. The Dallas–Fort Worth and Houston metropolitan areas number about 6.3 million and 5.7 million residents, respectively.
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Three interstate highways—I-35 to the west (Dallas–Fort Worth to San Antonio, with Austin in between), I-45 to the east (Dallas to Houston), and I-10 to the south (San Antonio to Houston) define the Texas Urban Triangle region. The region of 60,000 square miles (160,000 km2) contains most of the state's largest cities and metropolitan areas as well as 17 million people, nearly 75 percent of Texas's total population.[236] Houston and Dallas have been recognized as beta world cities.[237] These cities are spread out amongst the state. Texas has 254 counties, which is more than any other state by 95 (Georgia).[238]
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In contrast to the cities, unincorporated rural settlements known as colonias often lack basic infrastructure and are marked by poverty.[239] The office of the Texas Attorney General stated, in 2011, that Texas had about 2,294 colonias and estimates about 500,000 lived in the colonias. Hidalgo County, as of 2011, has the largest number of colonias.[240] Texas has the largest number of people of all states, living in colonias.[239]
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The most common accent or dialect spoken by natives throughout Texas is sometimes referred to as Texan English, which itself is a sub-variety of a broader category of American English known as Southern American English.[242][243] Creole language is spoken in East Texas.[244] In some areas of the state—particularly in the large cities—Western American English and General American English, have been on the increase. Chicano English—due to a growing Hispanic population—is widespread in South Texas, while African-American English is especially notable in historically minority areas of urban Texas.
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As of 2010, 65.8% (14,740,304) of Texas residents age 5 and older spoke only English at home, while 29.2% (6,543,702) spoke Spanish, 0.75 percent (168,886) Vietnamese, and Chinese (which includes Cantonese and Mandarin) was spoken by 0.56% (122,921) of the population over five.[245]
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Other languages spoken include German (including Texas German) by 0.33% (73,137), Tagalog with 0.29% (64,272) speakers, and French (including Cajun French) was spoken by 0.25% (55,773) of Texans.[245] Reportedly, Cherokee is the most widely spoken Native American language in Texas.[246]
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In total, 34.2% (7,660,406) of Texas's population aged five and older spoke a language at home other than English.[245]
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The largest denominations by number of adherents in 2010 were the Roman Catholic Church (4,673,500); the Southern Baptist Convention (3,721,318); Non-denominational Churches (1,546,542); and the United Methodist Church with (1,035,168).[248]
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Known as the buckle of the Bible Belt, East Texas is socially conservative.[249] The Dallas–Fort Worth metroplex is home to three major evangelical seminaries and a host of Bible schools. Lakewood Church in Houston, boasts the largest attendance in the nation averaging more than 43,000 weekly.[250]
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Adherents of many other religions reside predominantly in the urban centers of Texas. In 1990, the Islamic population was about 140,000 with more recent figures putting the current number of Muslims between 350,000 and 400,000.[251] The Jewish population is around 128,000.[252] Around 146,000 adherents of religions such as Hinduism and Sikhism live in Texas.[253] It is the fifth-largest Muslim-populated state in the country.[254]
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Historically, Texas culture comes from a blend of Southern (Dixie), Western (frontier), and Southwestern (Mexican/Anglo fusion) influences, varying in degrees of such from one intrastate region to another. Texas is placed in the Southern United States by the United States Census Bureau.[255] A popular food item, the breakfast burrito, draws from all three, having a soft flour tortilla wrapped around bacon and scrambled eggs or other hot, cooked fillings. Adding to Texas's traditional culture, established in the 18th and 19th centuries, immigration has made Texas a melting pot of cultures from around the world.
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Texas has made a strong mark on national and international pop culture. The entire state is strongly associated with the image of the cowboy shown in westerns and in country western music. The state's numerous oil tycoons are also a popular pop culture topic as seen in the hit TV series Dallas.
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The internationally known slogan "Don't Mess with Texas" began as an anti-littering advertisement. Since the campaign's inception in 1986, the phrase has become "an identity statement, a declaration of Texas swagger".[256]
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"Texas-sized" is an expression that can be used in two ways: to describe something that is about the size of the U.S. state of Texas,[257][258] or to describe something (usually but not always originating from Texas) that is large compared to other objects of its type.[259][260][261] Texas was the largest U.S. state, until Alaska became a state in 1959. The phrase "everything is bigger in Texas" has been in regular use since at least 1950;[262] and was used as early as 1913.[263]
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Houston is one of only five American cities with permanent professional resident companies in all the major performing arts disciplines: the Houston Grand Opera, the Houston Symphony Orchestra, the Houston Ballet, and The Alley Theatre.[264] Known for the vibrancy of its visual and performing arts, the Houston Theater District—a 17-block area in the heart of Downtown Houston—ranks second in the country in the number of theater seats in a concentrated downtown area, with 12,948 seats for live performances and 1,480 movie seats.[264]
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Founded in 1892, Modern Art Museum of Fort Worth, also called "The Modern", is Texas's oldest art museum. Fort Worth also has the Kimbell Art Museum, the Amon Carter Museum, the National Cowgirl Museum and Hall of Fame, the Will Rogers Memorial Center, and the Bass Performance Hall downtown. The Arts District of Downtown Dallas has arts venues such as the Dallas Museum of Art, the Morton H. Meyerson Symphony Center, the Margot and Bill Winspear Opera House, the Trammell & Margaret Crow Collection of Asian Art, and the Nasher Sculpture Center.[265]
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The Deep Ellum district within Dallas became popular during the 1920s and 1930s as the prime jazz and blues hotspot in the Southern United States. The name Deep Ellum comes from local people pronouncing "Deep Elm" as "Deep Ellum".[266] Artists such as Blind Lemon Jefferson, Robert Johnson, Huddie "Lead Belly" Ledbetter, and Bessie Smith played in early Deep Ellum clubs.[267]
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Austin, The Live Music Capital of the World, boasts "more live music venues per capita than such music hotbeds as Nashville, Memphis, Los Angeles, Las Vegas or New York City".[268] The city's music revolves around the nightclubs on 6th Street; events like the film, music, and multimedia festival South by Southwest; the longest-running concert music program on American television, Austin City Limits; and the Austin City Limits Music Festival held in Zilker Park.[269]
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Since 1980, San Antonio has evolved into "The Tejano Music Capital Of The World".[270] The Tejano Music Awards have provided a forum to create greater awareness and appreciation for Tejano music and culture.[271]
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The second president of the Republic of Texas, Mirabeau B. Lamar, is the Father of Texas Education. During his term, the state set aside three leagues of land in each county for equipping public schools. An additional 50 leagues of land set aside for the support of two universities would later become the basis of the state's Permanent University Fund.[272] Lamar's actions set the foundation for a Texas-wide public school system.[273]
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Between 2006 and 2007, Texas spent $7,275 per pupil ranking it below the national average of $9,389. The pupil/teacher ratio was 14.9, below the national average of 15.3. Texas paid instructors $41,744, below the national average of $46,593. The Texas Education Agency (TEA) administers the state's public school systems. Texas has over 1,000 school districts; all districts except the Stafford Municipal School District are independent from municipal government and many cross city boundaries.[274] School districts have the power to tax their residents and to assert eminent domain over privately owned property. Due to court-mandated equitable school financing for school districts, the state has a controversial tax redistribution system called the "Robin Hood plan". This plan transfers property tax revenue from wealthy school districts to poor ones.[275] The TEA has no authority over private or home school activities.[276]
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Students in Texas take the State of Texas Assessments of Academic Readiness (STAAR) in primary and secondary school. STAAR assess students' attainment of reading, writing, mathematics, science, and social studies skills required under Texas education standards and the No Child Left Behind Act. The test replaced the Texas Assessment of Knowledge and Skills (TAKS) test in the 2011–2012 school year.[277]
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Generally prohibited in the West at large, school corporal punishment is not unusual in the more conservative, rural areas of the state, with 28,569 public school students[278] paddled at least one time, according to government data for the 2011–2012 school year.[279] The rate of school corporal punishment in Texas is surpassed only by Mississippi, Alabama, and Arkansas.[279]
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The state's two most widely recognized flagship universities are The University of Texas at Austin and Texas A&M University, ranked as the 52nd[280] and 69th[281] best universities in the nation according to the 2014 edition of U.S. News & World Report's "Best Colleges", respectively. Some observers[282] also include the University of Houston and Texas Tech University as tier one flagships alongside UT Austin and A&M.[283][284] The Texas Higher Education Coordinating Board (THECB) ranks the state's public universities into three distinct tiers:[285]
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Texas's controversial alternative affirmative action plan, Texas House Bill 588, guarantees Texas students who graduated in the top 10 percent of their high school class automatic admission to state-funded universities. The bill encourages demographic diversity while avoiding problems stemming from the Hopwood v. Texas (1996) case.
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Thirty-six (36) separate and distinct public universities exist in Texas, of which 32 belong to one of the six state university systems.[288][289] Discovery of minerals on Permanent University Fund land, particularly oil, has helped fund the rapid growth of the state's two largest university systems: the University of Texas System and the Texas A&M System. The four other university systems: the University of Houston System, the University of North Texas System, the Texas State System, and the Texas Tech System are not funded by the Permanent University Fund.
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The Carnegie Foundation classifies three of Texas's universities as Tier One research institutions: The University of Texas at Austin, the Texas A&M University, and the University of Houston. The University of Texas at Austin and Texas A&M University are flagship universities of the state of Texas. Both were established by the Texas Constitution and hold stakes in the Permanent University Fund. The state has been putting effort to expand the number of flagship universities by elevating some of its seven institutions designated as "emerging research universities". The two expected to emerge first are the University of Houston and Texas Tech University, likely in that order according to discussions on the House floor of the 82nd Texas Legislature.[290]
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The state is home to various private institutions of higher learning—ranging from liberal arts colleges to a nationally recognized top-tier research university. Rice University in Houston is one of the leading teaching and research universities of the United States and is ranked the nation's 17th-best overall university by U.S. News & World Report.[291] Trinity University, a private, primarily undergraduate liberal arts university in San Antonio, has ranked first among universities granting primarily bachelor's and select master's degrees in the Western United States for 20 consecutive years by U.S. News.[292] Private universities include Abilene Christian University, Austin College, Baylor University, University of Mary Hardin–Baylor, and Southwestern University.[293][294][295]
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Universities in Texas host three presidential libraries: George Bush Presidential Library at Texas A&M University, the Lyndon Baines Johnson Library and Museum at The University of Texas at Austin, and the George W. Bush Presidential Library at Southern Methodist University.
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Notwithstanding the concentration of elite medical centers in the state, The Commonwealth Fund ranks the Texas healthcare system the third worst in the nation.[296] Texas ranks close to last in access to healthcare, quality of care, avoidable hospital spending, and equity among various groups.[296] Causes of the state's poor rankings include politics, a high poverty rate, and the highest rate of illegal immigration in the nation.[297] In May 2006, Texas initiated the program "code red" in response to the report the state had 25.1 percent of the population without health insurance, the largest proportion in the nation.[298]
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The Trust for America's Health ranked Texas 15th highest in adult obesity, with 27.2 percent of the state's population measured as obese.[299] The 2008 Men's Health obesity survey ranked four Texas cities among the top 25 fattest cities in America; Houston ranked 6th, Dallas 7th, El Paso 8th, and Arlington 14th.[300] Texas had only one city (Austin, ranked 21st) in the top 25 among the "fittest cities" in America.[300] The same survey has evaluated the state's obesity initiatives favorably with a "B+".[300] The state is ranked forty-second in the percentage of residents who engage in regular exercise.[301]
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Texas has the highest maternal mortality rate in the developed world, and the rate by which Texas women died from pregnancy-related complications doubled from 2010 to 2014, to 23.8 per 100,000. A rate unmatched in any other U.S. state or economically developed country.[302]
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Texas has many elite research medical centers. The state has nine medical schools,[303] three dental schools,[304] and two optometry schools.[305] Texas has two Biosafety Level 4 (BSL-4) laboratories: one at The University of Texas Medical Branch (UTMB) in Galveston,[306] and the other at the Southwest Foundation for Biomedical Research in San Antonio—the first privately owned BSL-4 lab in the United States.[307]
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The Texas Medical Center in Houston, holds the world's largest concentration of research and healthcare institutions, with 47 member institutions.[308] Texas Medical Center performs the most heart transplants in the world.[309] The University of Texas M. D. Anderson Cancer Center in Houston is a highly regarded academic institution that centers around cancer patient care, research, education and prevention.[310]
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San Antonio's South Texas Medical Center facilities rank sixth in clinical medicine research impact in the United States.[311] The University of Texas Health Science Center is another highly ranked research and educational institution in San Antonio.[312][313]
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Both the American Heart Association and the University of Texas Southwestern Medical Center call Dallas home. The Southwestern Medical Center ranks "among the top academic medical centers in the world".[314] The institution's medical school employs the most medical school Nobel laureates in the world.[314][315]
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Texans have historically had difficulties traversing Texas due to the state's large size and rough terrain. Texas has compensated by building America's largest highway and railway systems. The regulatory authority, the Texas Department of Transportation (TxDOT), maintains the state's immense highway system, regulates aviation,[316] and public transportation systems.[317]
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Located centrally in North America, the state is an important transportation hub. From the Dallas/Fort Worth area, trucks can reach 93 percent of the nation's population within 48 hours, and 37 percent within 24 hours.[318] Texas has 33 foreign trade zones (FTZ), the most in the nation.[319] In 2004, a combined total of $298 billion of goods passed through Texas FTZs.[319]
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The first Texas freeway was the Gulf Freeway opened in 1948 in Houston.[320] As of 2005, 79,535 miles (127,999 km) of public highway crisscrossed Texas (up from 71,000 miles (114,263 km) in 1984).[321] To fund recent growth in the state highways, Texas has 17 toll roads (see list) with several additional tollways proposed.[322] In central Texas, the southern section of the State Highway 130 toll road has a speed limit of 85 miles per hour (137 km/h), the highest in the nation.[323] All federal and state highways in Texas are paved.
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Texas has 730 airports, second-most of any state in the nation. Largest in Texas by size and passengers served, Dallas/Fort Worth International Airport (DFW) is the second-largest by area in the United States, and fourth in the world with 18,076 acres (73.15 km2).[324] In traffic, DFW airport is the busiest in the state, the fourth busiest in the United States,[325] and sixth worldwide.[326] American Airlines Group's American / American Eagle, the world's largest airline in total passengers-miles transported and passenger fleet size,[327] uses DFW as its largest and main hub. It ranks as the largest airline in the United States by number of passengers carried domestically per year and the largest airline in the world by number of passengers carried.[328]Southwest Airlines, headquartered in Dallas, has its operations at Dallas Love Field.[329]
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Texas's second-largest air facility is Houston's George Bush Intercontinental Airport (IAH). It served as the largest hub for the former Continental Airlines, which was based in Houston; it serves as the largest hub for United Airlines, the world's third-largest airline, by passenger-miles flown.[330][331] IAH offers service to the most Mexican destinations of any U.S. airport.[332][333] The next five largest airports in the state all serve more than three million passengers annually; they include Austin-Bergstrom International Airport, William P. Hobby Airport, San Antonio International Airport, Dallas Love Field and El Paso International Airport. The smallest airport in the state to be designated an international airport is Del Rio International Airport.
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Around 1,150 seaports dot Texas's coast with over 1,000 miles (1,600 km) of channels.[334] Ports employ nearly one-million people and handle an average of 317 million metric tons.[335] Texas ports connect with the rest of the U.S. Atlantic seaboard with the Gulf section of the Intracoastal Waterway.[334] The Port of Houston today is the busiest port in the United States in foreign tonnage, second in overall tonnage, and tenth worldwide in tonnage.[336] The Houston Ship Channel spans 530 feet (160 m) wide by 45 feet (14 m) deep by 50 miles (80 km) long.[337]
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Part of the state's tradition of cowboys is derived from the massive cattle drives which its ranchers organized in the nineteenth century to drive livestock to railroads and markets in Kansas, for shipment to the East. Towns along the way, such as Baxter Springs, the first cow town in Kansas, developed to handle the seasonal workers and tens of thousands of head of cattle being driven.
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The first railroad to operate in Texas was the Buffalo Bayou, Brazos and Colorado Railway, opening in August 1853.[338] The first railroad to enter Texas from the north, completed in 1872, was the Missouri–Kansas–Texas Railroad.[339] With increasing railroad access, the ranchers did not have to take their livestock up to the Midwest and shipped beef out from Texas. This caused a decline in the economies of the cow towns.
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Since 1911, Texas has led the nation in length of railroad miles within the state. Texas railway length peaked in 1932 at 17,078 miles (27,484 km), but declined to 14,006 miles (22,540 km) by 2000. While the Railroad Commission of Texas originally regulated state railroads, in 2005 the state reassigned these duties to TxDOT.[340]
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In the Dallas–Fort Worth area, three public transit agencies provide rail service: Dallas Area Rapid Transit (DART), Denton County Transportation Authority (DCTA), and Trinity Metro. DART began operating the first light rail system in the Southwest United States in 1996.[341]
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The Trinity Railway Express (TRE) commuter rail service, which connects Fort Worth and Dallas, is provided by Trinity Metro and DART.[342] Trinity Metro also operates the TEXRail commuter rail line, connecting downtown Fort Worth and Northeast Tarrant County to DFW Airport.[343] The A-train commuter rail line, operated by DCTA, acts as an extension of the DART Green line into Denton County.[344] In the Austin area, Capital Metropolitan Transportation Authority operates a commuter rail service known as Capital MetroRail to the northwestern suburbs. The Metropolitan Transit Authority of Harris County, Texas (METRO) operates light rail lines in the Houston area.
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Amtrak provides Texas with limited intercity passenger rail service. Three scheduled routes serve the state: the daily Texas Eagle (Chicago–San Antonio); the tri-weekly Sunset Limited (New Orleans–Los Angeles), with stops in Texas; and the daily Heartland Flyer (Fort Worth–Oklahoma City). Texas may get one of the nation's first high-speed rail line. Plans for a privately funded high-speed rail line between Dallas and Houston have been planned by the Texas Central Railway company.
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While American football has long been considered "king" in the state, Texans enjoy a wide variety of sports.[345]
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Texans can cheer for a plethora of professional sports teams. Within the "Big Four" professional leagues, Texas has two NFL teams (the Dallas Cowboys and the Houston Texans), two Major League Baseball teams (the Houston Astros and the Texas Rangers), three NBA teams (the San Antonio Spurs, the Houston Rockets, and the Dallas Mavericks), and one National Hockey League team (the Dallas Stars). The Dallas–Fort Worth metroplex is one of only twelve American metropolitan areas that host sports teams from all the "Big Four" professional leagues. Outside the "Big Four", Texas also has a WNBA team, (the Dallas Wings) and two Major League Soccer teams (the Houston Dynamo and FC Dallas).
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Collegiate athletics have deep significance in Texas culture, especially football. The state has twelve Division I-FBS schools, the most in the nation. Four of the state's universities, the Baylor Bears, Texas Longhorns, TCU Horned Frogs, and Texas Tech Red Raiders, compete in the Big 12 Conference. The Texas A&M Aggies left the Big 12 and joined the Southeastern Conference in 2012, which led the Big 12 to invite TCU to join; TCU was previously in the Mountain West Conference. The Houston Cougars and the SMU Mustangs compete in the American Athletic Conference. The Texas State Bobcats and the UT Arlington Mavericks compete in the Sun Belt Conference. Four of the state's schools claim at least one national championship in football: the Texas Longhorns, the Texas A&M Aggies, the TCU Horned Frogs, and the SMU Mustangs.
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According to a survey of Division I-A coaches the rivalry between the University of Oklahoma and the University of Texas at Austin, the Red River Shootout, ranks the third-best in the nation.[346] The TCU Horned Frogs and SMU Mustangs also share a rivalry and compete annually in the Battle for the Iron Skillet. A fierce rivalry, the Lone Star Showdown, also exists between the state's two largest universities, Texas A&M University and the University of Texas at Austin. The athletics portion of the Lone Star Showdown rivalry has been put on hold after the Texas A&M Aggies joined the Southeastern Conference.
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The University Interscholastic League (UIL) organizes most primary and secondary school competitions. Events organized by UIL include contests in athletics (the most popular being high school football) as well as artistic and academic subjects.[347]
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Texans also enjoy the rodeo. The world's first rodeo was hosted in Pecos, Texas.[348] The annual Houston Livestock Show and Rodeo is the largest rodeo in the world. It begins with trail rides from several points throughout the state that convene at Reliant Park.[349] The Southwestern Exposition and Livestock Show in Fort Worth is the oldest continuously running rodeo incorporating many of the state's most historic traditions into its annual events. Dallas hosts the State Fair of Texas each year at Fair Park.[350]
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Texas Motor Speedway hosts annual NASCAR Cup Series and IndyCar Series auto races since 1997. Since 2012, Austin's Circuit of the Americas plays host to a round of the Formula 1 World Championship—[351] the first at a permanent road circuit in the United States since the 1980 Grand Prix at Watkins Glen International—, as well as Grand Prix motorcycle racing, FIA World Endurance Championship and United SportsCar Championship races.
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Coordinates: 15°N 101°E / 15°N 101°E / 15; 101
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– in Asia (light green & dark grey)– in ASEAN (light green)
|
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Thailand,[a] officially the Kingdom of Thailand and formerly known as Siam,[b] is a country in Southeast Asia. Located at the centre of the Indochinese Peninsula, it is composed of 76 provinces, and covers an area of 513,120 square kilometres (198,120 sq mi), and a population of over 66 million people.[4] Thailand is the world's 50th-largest country by land area, and the 22nd-most-populous country in the world. The capital and largest city is Bangkok, a special administrative area. Thailand is bordered to the north by Myanmar and Laos, to the east by Laos and Cambodia, to the south by the Gulf of Thailand and Malaysia, and to the west by the Andaman Sea and the southern extremity of Myanmar. Its maritime boundaries include Vietnam in the Gulf of Thailand to the southeast, and Indonesia and India on the Andaman Sea to the southwest. Nominally, Thailand is a constitutional monarchy and parliamentary democracy; however, in recent history, its government has experienced multiple coups and periods of military dictatorships.
|
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|
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Tai peoples migrated from southwestern China to mainland Southeast Asia from the 11th century; the oldest known mention of their presence in the region by the exonym Siamese dates to the 12th century. Various Indianised kingdoms such as the Mon kingdoms, Khmer Empire and Malay states ruled the region, competing with Thai states such as the Kingdoms of Ngoenyang, Sukhothai, Lan Na and Ayutthaya, which rivalled each other. Documented European contact began in 1511 with a Portuguese diplomatic mission to Ayutthaya, which became a regional power by the end of the 15th century. Ayutthaya reached its peak during cosmopolitan Narai's reign (1656–1688), gradually declining thereafter until being ultimately destroyed in the 1767 Burmese–Siamese War. Taksin (r. 1767–1782) quickly reunified the fragmented territory and established the short-lived Thonburi Kingdom. He was succeeded in 1782 by Buddha Yodfa Chulaloke (r. 1782–1809), the first monarch of the current Chakri dynasty.
|
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|
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Through the 18th and 19th centuries, Siam faced imperialist pressure from France and the United Kingdom, including many unequal treaties with Western powers and forced concessions of territory; it nevertheless remained the only Southeast Asian country to avoid direct Western rule. Siamese system of government was centralized and transformed into modern unitary absolute monarchy in the reign of Chulalongkorn (r. 1868–1910). Siam joined World War I siding with the allies, a political decision to amend the unequal treaties. Following a bloodless revolution in 1932, Siam became a constitutional monarchy and changed its official name to "Thailand". Thailand was a satellite of Japan in World War II. In the late 1950s, a military coup under Field Marshal Sarit Thanarat revived the monarchy's historically influential role in politics. Thailand became a major ally of the United States, and played a key anti-communist role in the region as a member of the Southeast Asia Treaty Organization (SEATO). Apart from a brief period of parliamentary democracy in the mid-1970s, Thailand has periodically alternated between democracy and military rule. Since the 2000s, Thailand has been caught in a bitter political conflict between supporters and opponents of Thaksin Shinawatra, which culminated in two coups, most recently in 2014 and the establishment of its current and 20th constitution.
|
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|
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Thailand is a founding member of the Association of Southeast Asian Nations (ASEAN) and remains a major ally of the United States.[12][13] Despite comparatively sporadic changes in leadership, it is considered a regional power in Southeast Asia and a middle power in global affairs.[14] With a high level of human development, the second-largest economy in Southeast Asia, and the 20th-largest in the world by PPP, Thailand is classified as a newly industrialized economy; manufacturing, agriculture, and tourism are leading sectors of the economy.[15][16]
|
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|
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Thailand (/ˈtaɪlænd/ TY-land or /ˈtaɪlənd/ TY-lənd;[17] Thai: ประเทศไทย, RTGS: Prathet Thai, pronounced [pratʰêːt tʰaj] (listen)), officially the Kingdom of Thailand (Thai: ราชอาณาจักรไทย, RTGS: Ratcha-anachak Thai [râːtt͡ɕʰaʔaːnaːt͡ɕàk tʰaj] (listen), Chinese: 泰国), formerly known as Siam (Thai: สยาม, RTGS: Sayam [sajǎːm]), is a country at the centre of the Indochinese peninsula in Southeast Asia.
|
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|
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The country has always been called Mueang Thai by its citizens. By outsiders, prior to 1949, it was usually known by the exonym Siam (Thai: สยาม RTGS: sayam, pronounced [sajǎːm], also spelled Siem, Syâm, or Syâma). The word Siam may have originated from Pali (suvaṇṇabhūmi, 'land of gold') or Sanskrit श्याम (śyāma, 'dark') or Mon ရာမည(rhmañña, 'stranger'). The names Shan and A-hom seem to be variants of the same word. The word Śyâma is possibly not its origin, but a learned and artificial distortion.[clarification needed][18] Another theory is the name derives from Chinese: "Ayutthaya emerged as a dominant centre in the late fourteenth century. The Chinese called this region Xian, which the Portuguese converted into Siam."[19]:8 A further possibility is that Mon-speaking peoples migrating south called themselves syem as do the autochthonous Mon-Khmer-speaking inhabitants of the Malay Peninsula.[citation needed]
|
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|
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The signature of King Mongkut (r. 1851–1868) reads SPPM (Somdet Phra Poramenthra Maha) Mongkut Rex Siamensium (Mongkut King of the Siamese), giving the name Siam official status until 24 June 1939 when it was changed to "Thailand".[20] Thailand was renamed Siam from 1946 to 1948, after which it again reverted to "Thailand".
|
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|
21 |
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According to George Cœdès, the word Thai (ไทย) means 'free man' in the Thai language, "differentiating the Thai from the natives encompassed in Thai society as serfs".[21]:197 A famous Thai scholar argued that Thai (ไท) simply means 'people' or 'human being', since his investigation shows that in some rural areas the word "Thai" was used instead of the usual Thai word khon (คน) for people.[22] According to Michel Ferlus, the ethnonyms Thai-Tai (or Thay-Tay) would have evolved from the etymon *k(ə)ri: 'human being' through the following chain: *kəri: > *kəli: > *kədi:/*kədaj > *di:/*daj > *dajA (Proto-Southwestern Tai) > tʰajA2 (in Siamese and Lao) or > tajA2 (in the other Southwestern and Central Tai languages classified by Li Fangkuei).[23] Michel Ferlus's work is based on some simple rules of phonetic change observable in the Sinosphere and studied for the most part by William H. Baxter (1992).[24]
|
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|
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While Thai people will often refer to their country using the polite form prathet Thai (Thai: ประเทศไทย), they most commonly use the more colloquial term mueang Thai (Thai: เมืองไทย) or simply Thai; the word mueang, archaically referring to a city-state, is commonly used to refer to a city or town as the centre of a region. Ratcha Anachak Thai (Thai: ราชอาณาจักรไทย) means 'kingdom of Thailand' or 'kingdom of Thai'. Etymologically, its components are: ratcha (Sanskrit: राजन्, rājan, 'king, royal, realm'); -ana- (Pali āṇā 'authority, command, power', itself from the Sanskrit आज्ञा, ājñā, of the same meaning) -chak (from Sanskrit चक्र cakra- 'wheel', a symbol of power and rule). The Thai National Anthem (Thai: เพลงชาติ), written by Luang Saranupraphan during the patriotic 1930s, refers to the Thai nation as prathet Thai (Thai: ประเทศไทย). The first line of the national anthem is: prathet thai ruam lueat nuea chat chuea thai (Thai: ประเทศไทยรวมเลือดเนื้อชาติเชื้อไทย), 'Thailand is the unity of Thai flesh and blood'.
|
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There is evidence of continuous human habitation in present-day Thailand from 20,000 years ago to the present day.[26]:4 The earliest evidence of rice growing is dated at 2,000 BCE.[25]:4 Bronze appeared circa 1,250–1,000 BCE.[25]:4 The site of Ban Chiang in northeast Thailand currently ranks as the earliest known centre of copper and bronze production in Southeast Asia.[27] Iron appeared around 500 BCE.[25]:5 The Kingdom of Funan was the first and most powerful Southeast Asian kingdom at the time (2nd century BCE).[26]:5 The Mon people established the principalities of Dvaravati and Kingdom of Hariphunchai in the 6th century. The Khmer people established the Khmer empire, centred in Angkor, in the 9th century.[26]:7 Tambralinga, a Malay state controlling trade through the Malacca Strait, rose in the 10th century.[26]:5 The Indochina peninsula was heavily influenced by the culture and religions of India from the time of the Kingdom of Funan to that of the Khmer Empire.[28]
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The Thai people are of the Tai ethnic group, characterised by common linguistic roots.[29]:2 Chinese chronicles first mention the Tai peoples in the 6th century BCE. While there are many assumptions regarding the origin of Tai peoples, David K. Wyatt, a historian of Thailand, argued that their ancestors which at the present inhabit Laos, Thailand, Myanmar, India, and China came from the Điện Biên Phủ area between the 5th and the 8th century.[29]:6 Thai people began migrating into present-day Thailand around the 11th century, which Mon and Khmer people occupied at the time.[30] Thus Thai culture was influenced by Indian, Mon, and Khmer cultures.[31]
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|
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According to French historian George Cœdès, "The Thai first enter history of Farther India in the eleventh century with the mention of Syam slaves or prisoners of war in Champa epigraphy, and "in the twelfth century, the bas-reliefs of Angkor Wat" where "a group of warriors" are described as Syam.[21]:190–191, 194–195
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After the decline of the Khmer Empire and Kingdom of Pagan in the early-13th century, various states thrived in their place. The domains of Tai people existed from the northeast of present-day India to the north of present-day Laos and to the Malay peninsula.[29]:38–9 During the 13th century, Tai people had already settled in the core land of Dvaravati and Lavo Kingdom to Nakhon Si Thammarat in the south. There are, however, no records detailing the arrival of the Tais.[29]:50–1
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|
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Around 1240, Pho Khun Bang Klang Hao, a local Tai ruler, rallied the people to rebel against the Khmer. He later crowned himself the first king of Sukhothai Kingdom in 1238.[29]:52–3 Mainstream Thai historians count Sukhothai as the first kingdom of Thai people. Sukhothai expanded furthest during the reign of Ram Khamhaeng (r. 1279–1298). However, it was mostly a network of local lords who swore fealty to Sukhothai, not directly controlled by it.[29]:55–6 He is believed have invented Thai script and Thai ceramics were an important export in his era. Sukhothai embraced Theravada Buddhism in the reign of Maha Thammaracha I (1347–1368).
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To the north, Mangrai, who descended from a local ruler lineage of Ngoenyang, founded the kingdom of Lan Na in 1292, centered in Chiang Mai. He unified the surrounding area and his dynasty would rule the kingdom continuously for the next two centuries. He also created a network of states through political alliances to the east and north of the Mekong.[19]:8 While in the port in Lower Chao Phraya Basin, a federation around Phetchaburi, Suphan Buri, Lopburi, and the Ayutthaya area was created in the 11th century.[19]:8
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According to the most widely accepted version of its origin, the Ayutthaya Kingdom rose from the earlier, nearby Lavo Kingdom and Suvarnabhumi with Uthong as its first king. Ayutthaya was a patchwork of self-governing principalities and tributary provinces owing allegiance to the King of Ayutthaya under the mandala system.[32]:355 Its initial expansion was through conquest and political marriage. Before the end of the 15th century, Ayutthaya invaded the Khmer Empire three times and sacked its capital Angkor.[33]:26 Ayutthaya then became a regional power in place of the Khmer. Constant interference of Sukhothai effectively made it a vassal state of Ayutthaya and it was finally incorporated into the kingdom. Borommatrailokkanat brought about bureaucratic reforms which lasted into the 20th century and created a system of social hierarchy called sakdina, where male commoners were conscripted as corvée labourers for six months a year.[34]:107 Ayutthaya was interested in the Malay peninsula, but failed to conquer the Malacca Sultanate which was supported by the Chinese Ming Dynasty.[26]:11, 13
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European contact and trade started in the early-16th century, with the envoy of Portuguese duke Afonso de Albuquerque in 1511, Portugal became an allied and ceded some soldiers to King Rama Thibodi II.[35] The Portuguese were followed in the 17th century by the French, Dutch, and English. Rivalry for supremacy over Chiang Mai and the Mon people pitted Ayutthaya against the Burmese Kingdom. Several wars with its ruling dynasty Taungoo Dynasty starting in the 1540s in the reign of Tabinshwehti and Bayinnaung were ultimately ended with the capture of the capital in 1570.[34]:146–7 Then was a brief period of vassalage to Burma until Naresuan proclaimed independence in 1584.[19]:11
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Ayutthaya then sought to improve relations with European powers for many successive reigns. The kingdom especially prospered during cosmopolitan Narai's reign (1656–1688) when some European travelers regarded Ayutthaya as an Asian great power, alongside China and India.[25]:ix However, growing French influence later in his reign was met with nationalist sentiment and led eventually to the Siamese revolution of 1688.[34]:185–6 However, overall relations remained stable, with French missionaries still active in preaching Christianity.[34]:186
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After a bloody period of dynastic struggle, Ayutthaya entered into what has been called the golden age, a relatively peaceful episode in the second quarter of the 18th century when art, literature, and learning flourished. There were seldom foreign wars, apart from conflict with the Nguyễn Lords for control of Cambodia starting around 1715. The last fifty years of the kingdom witnessed bloody succession crises, where there were purges of court officials and able generals for many consecutive reigns. In 1765, a combined 40,000-strong force of Burmese armies invaded it from the north and west.[36]:250 The Burmese were under the new Alaungpaya dynasty quickly rose to be a new local power by 1759. After a 14-month siege, the capital city's wall fell and the city was burned in April 1767.[37]:218
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The capital and much territories lied in chaos after the war. The former capital was occupied by the Burmese garrison army and five local leaders declared themselves overlords, including the lords of Sakwangburi, Pimai, Chanthaburi, and Nakhon Si Thammarat. Chao Tak, a capable military leader, proceeded to make himself a lord by right of conquest, beginning with the legendary sack of Chanthaburi. Based at Chanthaburi, Chao Tak raised troops and resources, and sent a fleet up the Chao Phraya to take the fort of Thonburi. In the same year, Chao Tak was able to retake Ayutthaya from the Burmese only seven months after the fall of the city.[38]
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Chao Tak then crowned himself as Taksin and proclaimed Thonburi as temporary capital in the same year. He also quickly subdued the other warlords. His forces engaged in wars with Burma, Laos, and Cambodia, which successfully drove the Burmese out of Lan Na in 1775,[34]:225 captured Vientiane in 1778[34]:227–8 and tried to install a pro-Thai king in Cambodia in the 1770s. In his final years there was a coup, caused supposedly by his "insanity", and eventually Taksin and his sons were executed by his longtime companion General Chao Phraya Chakri (the future Rama I). He was the first king of the ruling Chakri Dynasty and founder of the Rattanakosin Kingdom on 6 April 1782.
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Under Rama I (1782–1809), Rattanakosin successfully defended against Burmese attacks and put an end to Burmese incursions. He also created suzerainty over large portions of Laos and Cambodia.[39] In 1821, Briton John Crawfurd was sent to negotiate a new trade agreement with Siam – the first sign of an issue which was to dominate 19th century Siamese politics.[40] Bangkok signed the Burney Treaty in 1826, after the British victory in the First Anglo-Burmese War.[34]:281 Anouvong of Vientiane, who misunderstood that Britain was about to attack Bangkok, started the Lao rebellion in 1826 and was defeated.[34]:283–5 Vientiane was destroyed and a large number of Lao people was relocated to Khorat Plateau as a result.[34]:285–6 Bangkok also waged several wars with Vietnam, where Bangkok successfully regained hegemony over Cambodia.[34]:290–2
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From the late-19th century, Siam tried to rule the ethnic groups in the realm as colonies.[34]:308 In the reign of Mongkut (1851–1868), who recognised the threat of Western powers, his court contacted the British government directly to defuse tensions.[34]:311 A British mission led by Sir John Bowring, Governor of Hong Kong, led to the signing of the Bowring Treaty, the first of many unequal treaties with Western countries. This, however, brought trade and economic development in Bangkok.[41] The unexpected death of Mongkut from malaria led to the reign of underage Prince Chulalongkorn, with Somdet Chaophraya Sri Suriwongse (Chuang Bunnag) acting as regent.[34]:327
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Chulalongkorn (r. 1868–1910) initiated centralisation, set up a privy council, and abolished slavery and the corvée system.[34] The Front Palace crisis of 1874 stalled attempts at further reforms.[34]:331–3 In the 1870s and 1880s, he incorporated the protectorates up north into the kingdom proper, which later expanded to the protectorates in the northeast and the south.[34]:334–5 He established twelve krom in 1888, which were equivalent to present-day ministries.[34]:347 The crisis of 1893 erupted, caused by French demands for Lao territory east of Mekong.[34]:350–3 Thailand is the only Southeast Asian nation not to have been colonised by a Western power,[42] in part because Britain and France agreed in 1896 to make the Chao Phraya valley a buffer state.[43] Not until the 20th century could Siam renegotiate every unequal treaty dating from the Bowring Treaty, including extraterritoriality, but at a price of many territorial exchanges. The advent of the monthon system marked the creation of the modern Thai nation-state.[34]:362–3 In 1905, there were unsuccessful rebellions in the ancient Patani area, Ubon Ratchathani, and Phrae in opposition to an attempt to blunt the power of local lords.[34]:371–3
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The Palace Revolt of 1912 was a failed attempt by Western-educated military officers to overthrow the absolute monarchy.[34]:397 Vajiravudh (r. 1910–1925) responded by propaganda for the entirety of his reign.[34]:402 He promoted the idea of the Thai nation.[34]:404 In 1917, Siam joined World War I on the side of the Allies as there were concerns that the Allies might punish neutral countries and refuse to amend past unequal treaties.[34]:407 In the aftermath Siam joined the Paris Peace Conference, and gained freedom of taxation and the revocation of extraterritoriality.[34]:408
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A bloodless revolution took place in 1932, carried out by a group of military and civilian officials Khana Ratsadon. Prajadhipok was forced to grant the country's first constitution, thereby ending centuries of absolute monarchy. The combined results of economic hardships brought on by the Great Depression, sharply falling rice prices, and a significant reduction in public spending caused discontent among aristocrats.[26]:25 In 1933, A counter-revolutionary rebellion occurred which aimed to reinstate absolute monarchy, but failed.[34]:446–8 Prajadhipok's conflict with the government eventually led to abdication. The government selected Ananda Mahidol, who was studying in Switzerland, to be the new king.[34]:448–9
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Later that decade, the army wing of Khana Ratsadon came to dominate Siamese politics. Plaek Phibunsongkhram who became premier in 1938, started political oppression and took an openly anti-royalist stance.[34]:457 His government adopted nationalism and Westernisation, anti-Chinese and anti-French policies.[26]:28 In 1940, there was a decree changing the name of the country from "Siam" to "Thailand". In 1941, Thailand was in a brief conflict with Vichy France resulting in Thailand gaining some Lao and Cambodian territories.[34]:462 On 8 December 1941, the Empire of Japan launched an invasion of Thailand, and fighting broke out shortly before Phibun ordered an armistice. Japan was granted free passage, and on 21 December Thailand and Japan signed a military alliance with a secret protocol, wherein Tokyo agreed to help Thailand regain territories lost to the British and French.[44] The Thai government declared war on the United States and the United Kingdom.[34]:465 The Free Thai Movement was launched both in Thailand and abroad to oppose the government and Japanese occupation.[34]:465–6 After the war ended in 1945, Thailand signed formal agreements to end the state of war with the Allies. Most Allied powers had not recognised Thailand's declaration of war.
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In June 1946, young King Ananda was found dead under mysterious circumstances. His younger brother Bhumibol Adulyadej ascended to the throne. Thailand joined the Southeast Asia Treaty Organization (SEATO) to become an active ally of the United States in 1954.[34]:493 Field Marshal Sarit Thanarat launched a coup in 1957, which removed Khana Ratsadon from politics. His rule (premiership 1959–1963) was autocratic; he built his legitimacy around the god-like status of the monarch and by channelling the government's loyalty to the king.[34]:511 His government improved the country's infrastructure and education.[34]:514 After the US joined the Vietnam War in 1961, there was a secret agreement wherein the US promised to protect Thailand.[34]:523
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The period brought about increasing modernisation and Westernisation of Thai society. Rapid urbanisation occurred when the rural populace sought work in growing cities. Rural farmers gained class consciousness and were sympathetic to the Communist Party of Thailand.[34]:528 Economic development and education enabled the rise of a middle class in Bangkok and other cities.[34]:534 In October 1971, there was a large demonstration against the dictatorship of Thanom Kittikachorn (premiership 1963–1973), which led to civilian casualties.[34]:541–3 Bhumibol installed Sanya Dharmasakti (premiership 1973–1975) to replace him, making it the first time that the king intervened in Thai politics directly since 1932.[45] The aftermath of the event marked a short-lived parliamentary democracy,[45] often called the "era when democracy blossomed." (ยุคประชาธิปไตยเบ่งบาน)
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Constant unrest and instability, as well as fear of a communist takeover after the fall of Saigon, made some ultra-right groups brand leftist students as communists.[34]:548 This culminated in the Thammasat University massacre in October 1976.[34]:548–9 A coup d'état on that day brought Thailand a new ultra-right government, which cracked down on media outlets, officials, and intellectuals, and fuelled the communist insurgency. Another coup the following year installed a more moderate government, which offered amnesty to communist fighters in 1978.
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Fueled by Indochina refugee crisis, Vietnamese border raids and economic hardships, Prem Tinsulanonda launched a successful coup and became the Prime Minister from 1980 to 1988. The communists abandoned the insurgency by 1983. Prem's premiership was dubbed "semi-democracy" because the Parliament was composed of all elected House and all appointed Senate. The 1980s also saw increasing intervention in politics by the monarch, who rendered two coup attempts against Prem failed. Thailand had its first elected prime minister in 1988.[46]
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Suchinda Kraprayoon, who was the coup leader in 1991 and said he would not seek to become prime minister, was nominated as one by the majority coalition government after the 1992 general election. This caused a popular demonstration in Bangkok, which ended with a military crackdown. Bhumibol intervened in the event and Suchinda then resigned.
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The 1997 Asian financial crisis originated in Thailand and ended the country's 40 years of uninterrupted economic growth.[47]:3 Chuan Leekpai's government took an IMF loan with unpopular provisions.[34]:576 The populist Thai Rak Thai party, led by prime minister Thaksin Shinawatra, governed from 2001 until 2006. His policies were successful in reducing rural poverty[48] and initiated universal healthcare in the country.[49] A South Thailand insurgency escalated starting from 2004. The 2004 Indian Ocean earthquake and tsunami hit the country, mostly in the south. Massive protests against Thaksin led by the People's Alliance for Democracy (PAD) started in his second term as prime minister and his tenure ended with a coup d'état in 2006. The junta installed a military government which lasted a year.
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In 2007, a civilian government led by the Thaksin-allied People's Power Party (PPP) was elected. Another protest led by PAD ended with the dissolution of PPP, and the Democrat Party led a coalition government in its place. The pro-Thaksin United Front for Democracy Against Dictatorship (UDD) protested both in 2009 and in 2010.
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After the general election of 2011, the populist Pheu Thai Party won a majority and Yingluck Shinawatra, Thaksin's younger sister, became prime minister. The People's Democratic Reform Committee organised another anti-Shinawatra protest[c] after the ruling party proposed an amnesty bill which would benefit Thaksin.[50] Yingluck dissolved parliament and a general election was scheduled, but was invalidated by the Constitution Court. The crisis ended with another coup d'état in 2014, the second coup in a decade.[d] The National Council for Peace and Order, a military junta led by General Prayut Chan-o-cha, has led the country since. Civil and political rights were restricted, and the country saw a surge in lèse-majesté cases. Political opponents and dissenters were sent to "attitude adjustment" camps.[51] Bhumibol, the longest-reigning Thai king, died in 2016, and his son Vajiralongkorn ascended to the throne. The referendum and adoption of Thailand's current constitution happened under the junta's rule.[e] In 2019, the junta agreed to schedule a general election in March.[51] Prayut continued his premiership with the support of Palang Pracharath Party-coalition in the House and junta-appointed Senate, amid allegations of election fraud.[53]
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Prior to 1932, Thai kings were absolute monarchs. During Sukhothai Kingdom, the king was seen as a Dharmaraja or 'king who rules in accordance with Dharma'. The system of government was a network of tributaries ruled by local lords. Modern absolute monarchy and statehood was established by Chulalongkorn when he transformed the decentralized protectorate system into a unitary state. On 24 June 1932, Khana Ratsadon (People's Party) carried out a bloodless revolution which marked the beginning of constitutional monarchy.
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Thailand has had 20 constitutions and charters since 1932, including the latest and current 2017 Constitution. Throughout this time, the form of government has ranged from military dictatorship to electoral democracy.[54][55] Thailand has had the fourth-most coups in the world.[56] "Uniformed or ex-military men have led Thailand for 55 of the 83 years" between 1932 and 2009.[57] Most recently, the National Council for Peace and Order ruled the country between 2014 and 2019.
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The politics of Thailand is conducted within the framework of a constitutional monarchy, whereby a hereditary monarch is head of state. The current King of Thailand is Vajiralongkorn (or Rama X), who has reigned since October 2016. The powers of the king are limited by the constitution and he is primarily a symbolic figurehead. The monarch is head of the armed forces and is required to be Buddhist as well as the Defender of the Faith. He has the power to appoint his heirs, the power to grant pardons, and the royal assent. The king is aided in his duties by the Privy Council of Thailand. However, the monarch still occasionally intervenes in Thai politics, as all constitutions pave the way for customary royal rulings. The monarchy is widely revered and lèse majesté is a severe crime in Thailand.
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Government is separated into three branches:
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Military and bureaucratic aristocrats fully controlled political parties between 1946 and 1980s.[59]:16 Most parties in Thailand are short-lived.[60]:246 Between 1992 and 2006, Thailand had a two-party system.[60]:245 Since 2000, two political parties dominated Thai general elections: one was the Pheu Thai Party (which was a successor of People's Power Party and the Thai Rak Thai Party), and the other was the Democrat Party. The political parties which support Thaksin Shinawatra won the most representatives every general election since 2001. Later constitutions created a multi-party system where a single party cannot gain a majority in the house.
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The 2007 constitution was partially abrogated by the military dictatorship that came to power in May 2014.[61]
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Thailand's kings are protected by lèse-majesté laws which allow critics to be jailed for three to fifteen years.[62] After the 2014 Thai coup d'état, Thailand had the highest number of lèse-majesté prisoners in the nation's history.[63][64] In 2017, the military court in Thailand sentenced a man to 35 years in prison for violating the country's lèse-majesté law.[64] Thailand has been rated not free on the Freedom House Index since 2014.[65] Thai activist and magazine editor Somyot Prueksakasemsuk, who was sentenced to eleven years' imprisonment for lèse-majesté in 2013,[66] is a designated prisoner of conscience by Amnesty International.[67]
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Totalling 513,120 square kilometres (198,120 sq mi), Thailand is the 50th-largest country by total area. It is slightly smaller than Yemen and slightly larger than Spain.[1]
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Thailand comprises several distinct geographic regions, partly corresponding to the provincial groups. The north of the country is the mountainous area of the Thai highlands, with the highest point being Doi Inthanon in the Thanon Thong Chai Range at 2,565 metres (8,415 ft) above sea level. The northeast, Isan, consists of the Khorat Plateau, bordered to the east by the Mekong River. The centre of the country is dominated by the predominantly flat Chao Phraya river valley, which runs into the Gulf of Thailand.
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Southern Thailand consists of the narrow Kra Isthmus that widens into the Malay Peninsula. Politically, there are six geographical regions which differ from the others in population, basic resources, natural features, and level of social and economic development. The diversity of the regions is the most pronounced attribute of Thailand's physical setting.
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The Chao Phraya and the Mekong River are the indispensable water courses of rural Thailand. Industrial scale production of crops use both rivers and their tributaries. The Gulf of Thailand covers 320,000 square kilometres (124,000 sq mi) and is fed by the Chao Phraya, Mae Klong, Bang Pakong, and Tapi Rivers. It contributes to the tourism sector owing to its clear shallow waters along the coasts in the southern region and the Kra Isthmus. The eastern shore of the Gulf of Thailand is an industrial centre of Thailand with the kingdom's premier deepwater port in Sattahip and its busiest commercial port, Laem Chabang.
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The Andaman Sea is a precious natural resource as it hosts popular and luxurious resorts. Phuket, Krabi, Ranong, Phang Nga and Trang, and their islands, all lay along the coasts of the Andaman Sea and, despite the 2004 tsunami, they remain a tourist magnet.
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Thailand's climate is influenced by monsoon winds that have a seasonal character (the southwest and northeast monsoon).[68]:2 Most of the country is classified as Köppen's tropical savanna climate.[69] The majority of the south as well as the eastern tip of the east have a tropical monsoon climate. Parts of the south also have a tropical rainforest climate.
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Thailand is divided into three seasons.[68]:2 The first is the rainy or southwest monsoon season (mid–May to mid–October), which is caused by southwestern wind from Indian Ocean.[68]:2 Rainfall is also contributed by Intertropical Convergence Zone (ITCZ) and tropical cyclones.[68]:2 August and September being the wettest period of the year.[68]:2 The country receives a mean annual rainfall of 1,200 to 1,600 mm (47 to 63 in).[68]:4 Winter or the northeast monsoon starts from mid–October until mid–February.[68]:2 Most of Thailand experiences dry weather with mild temperatures.[68]:2,4 The exception is southern Thailand where it receives abundant rainfall, particularly during October to November.[68]:2 Summer or the pre–monsoon season runs from mid–February until mid–May.[68]:3 Due to its inland nature and latitude, the north, northeast, central and eastern parts of Thailand experience a long period of warm weather, where emperatures can reach up to 40 °C (104 °F) during March to May,[68]:3 in contrast to close to or below 0 °C (32 °F) in some areas in winter.[68]:3 Southern Thailand is characterised by mild weather year-round with less diurnal and seasonal variations in temperatures due to maritime influences.[68]:3
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Thailand is among the world's ten countries that are most exposed to climate change; in particular, it is highly vulnerable to rising sea levels and extreme weather events.[70][71]
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Thailand has a mediocre but improving performance in the global Environmental Performance Index (EPI) with an overall ranking of 91 out of 180 countries in 2016. The environmental areas where Thailand performs worst (i.e., highest ranking) are air quality (167), environmental effects of the agricultural industry (106), and the climate and energy sector (93), the later mainly because of a high CO2 emission per KWh produced. Thailand performs best (i.e., lowest ranking) in water resource management (66), with some major improvements expected for the future, and sanitation (68).[73][74]
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The population of elephants, the country's national symbol, has fallen from fallen from 100,000 in 1850 to an estimated 2,000.[72] Poachers have long hunted elephants for ivory and hides, and now increasingly for meat.[75] Young elephants are often captured for use in tourist attractions or as work animals, which there were claims of mistreatment.[76] Although their use has declined since the government banned logging in 1989.
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Poaching of protected species remains a major problem. Tigers, leopards, and other large cats are hunted for their pelts. Many are farmed or hunted for their meat, supposedly has medicinal properties. Although such trade is illegal, the well-known Bangkok market Chatuchak is still known for the sale of endangered species.[77] The practice of keeping wild animals as pets affects species such as Asiatic black bear, Malayan sun bear, white-handed lar, pileated gibbon, and binturong.[78]
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Thailand is a unitary state; the administrative services of the executive branch are divided into three levels by National Government Organisation Act, BE 2534 (1991): central, provincial and local. Thailand is composed of 76 provinces (จังหวัด, changwat), which are first-level administrative divisions.[79] There are also two specially governed districts: the capital Bangkok and Pattaya. Bangkok is at provincial level and thus often counted as a province. Each province is divided into districts (อำเภอ, amphoe) and the districts are further divided into sub-districts (ตำบล, tambons). The name of each province's capital city (เมือง, mueang) is the same as that of the province. For example, the capital of Chiang Mai Province (Changwat Chiang Mai) is Mueang Chiang Mai or Chiang Mai. All provincial governors and district chiefs, which are administrators of provinces and districts respectively, are appointed by the central government.[80] Thailand's provinces are sometimes grouped into four to six regions, depending on the source.
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The foreign relations of Thailand are handled by the Minister of Foreign Affairs.
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Thailand participates fully in international and regional organisations. It is a major non-NATO ally and Priority Watch List Special 301 Report of the United States. The country remains an active member of ASEAN Association of Southeast Asian Nations. Thailand has developed increasingly close ties with other ASEAN members: Indonesia, Malaysia, the Philippines, Singapore, Brunei, Laos, Cambodia, Myanmar, and Vietnam, whose foreign and economic ministers hold annual meetings. Regional co-operation is progressing in economic, trade, banking, political, and cultural matters. In 2003, Thailand served as APEC (Asia Pacific Economic Cooperation) host. Dr. Supachai Panitchpakdi, the former Deputy Prime Minister of Thailand, currently serves as Secretary-General of the United Nations Conference on Trade and Development (UNCTAD). In 2005 Thailand attended the inaugural East Asia Summit.
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In recent years, Thailand has taken an increasingly active role on the international stage. When East Timor gained independence from Indonesia, Thailand, for the first time in its history, contributed troops to the international peacekeeping effort. Its troops remain there today as part of a UN peacekeeping force. As part of its effort to increase international ties, Thailand has reached out to such regional organisations as the Organization of American States (OAS) and the Organisation for Security and Cooperation in Europe (OSCE). Thailand has contributed troops to reconstruction efforts in Afghanistan and Iraq.
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Thaksin initiated negotiations for several free trade agreements with China, Australia, Bahrain, India, and the US. The latter especially was criticised, with claims that uncompetitive Thai industries could be wiped out.[81]
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Thaksin also announced that Thailand would forsake foreign aid, and work with donor countries to assist in the development of neighbours in the Greater Mekong Sub-region.[82] Thaksin sought to position Thailand as a regional leader, initiating various development projects in poorer neighbouring countries like Laos. More controversially, he established close, friendly ties with the Burmese dictatorship.[83]
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Thailand joined the US-led invasion of Iraq, sending a 423-strong humanitarian contingent.[84] It withdrew its troops on 10 September 2004. Two Thai soldiers died in Iraq in an insurgent attack.
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Abhisit appointed Peoples Alliance for Democracy leader Kasit Piromya as foreign minister. In April 2009, fighting broke out between Thai and Cambodian troops on territory immediately adjacent to the 900-year-old ruins of Cambodia's Preah Vihear Hindu temple near the border. The Cambodian government claimed its army had killed at least four Thais and captured 10 more, although the Thai government denied that any Thai soldiers were killed or injured. Two Cambodian and three Thai soldiers were killed. Both armies blamed the other for firing first and denied entering the other's territory.[85][86]
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The Royal Thai Armed Forces (กองทัพไทย; RTGS: Kong Thap Thai) constitute the military of the Kingdom of Thailand. It consists of the Royal Thai Army (กองทัพบกไทย), the Royal Thai Navy (กองทัพเรือไทย), and the Royal Thai Air Force (กองทัพอากาศไทย). It also incorporates various paramilitary forces.
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The Thai Armed Forces have a combined manpower of 306,000 active duty personnel and another 245,000 active reserve personnel.[87] The head of the Thai Armed Forces (จอมทัพไทย, Chom Thap Thai) is the king,[88] although this position is only nominal. The armed forces are managed by the Ministry of Defence of Thailand, which is headed by the Minister of Defence (a member of the cabinet of Thailand) and commanded by the Royal Thai Armed Forces Headquarters, which in turn is headed by the Chief of Defence Forces of Thailand.[89] Thai annual defense budget almost tripled from 78 billion baht in 2005 to 207 billion baht in 2016, accounting for approximately 1.5% of 2019 Thai GDP.[90][91] Thailand ranked 16th worldwide in the Military Strength Index based on the Credit Suisse report in September 2015.
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The military is also tasked with humanitarian missions, such as escorting Rohingya to Malaysia or Indonesia,[92] ensuring security and welfare for refugees during Indochina refugee crisis.[93]
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According to the constitution, serving in the armed forces is a duty of all Thai citizens.[94] Thailand still use active draft system for males over the age of 21. They are subjected to varying lengths of active service depending on the duration of reserve training as Territorial Defence Student and their level of education. Those who have completed three years or more of reserve training will be exempted entirely. The practice has long been criticized, as some media question its efficacy and value.[95][96] It is alleged that conscripts end up as servants to senior officers[97] or clerks in military cooperative shops.[98][99] In a report issued in March 2020, Amnesty International charged that Thai military conscripts face institutionalised abuse systematically hushed up by military authorities.[100]
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Critics observed that Thai military's main objective is to deal with internal rather than external threats.[101] Internal Security Operations Command is called the political arm of the Thai military, which has overlapping social and political functions with civilian bureaucracy. It also has anti-democracy mission.[101] The military is also notorious for numerous corruption incidents, such as accusation of human trafficking,[102] and nepotism in promotion of high-ranking officers.[103] The military is deeply entrenched in politics. Most recently, the appointed senators include more than 100 active and retired military.[104]
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In 2018 the literacy rate was 93.8%. The youth literacy rate was 98.1% in 2015.[105] Education is provided by a well-organised school system of kindergartens, primary, lower secondary and upper secondary schools, numerous vocational colleges, and universities. The private sector of education is well developed and significantly contributes to the overall provision of education which the government would not be able to meet with public establishments. Education is compulsory up to and including age 14, with the government providing free education through to age 17. Thailand is the 3rd most popular study destination in Asean. The number of international degree students in Thailand increased by fully 979% between 1999 and 2012, from 1,882 to 20,309 students. The most of international students come from Asian neighbor countries[106] from China, Myanmar, Cambodia and Vietnam.[107] The number of higher education institutions in Thailand has grown strongly over the past decades from just a handful of universities in the 1970s to 156 officially. The two top-ranking universities in Thailand are Chulalongkorn University and Mahidol University.[108] Thai universities research output still relatively low by international ranking comparison, Recent initiatives, such as the National Research University from 9 universities around the country[109] and Graduate research intensive university: VISTEC, designed to strengthen Thailand's national research universities, however, appear to be gaining traction. Thailand's research output, as measured by journal publications, increased by 20% between 2011 and 2016.[110]
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Teaching relies heavily on rote learning rather than on student-centred methodology. The establishment of reliable and coherent curricula for its primary and secondary schools is subject to such rapid changes that schools and their teachers are not always sure what they are supposed to be teaching, and authors and publishers of textbooks are unable to write and print new editions quickly enough to keep up with the volatility. Issues concerning university entrance has been in constant upheaval for a number of years. Nevertheless, Thai education has seen its greatest progress in the years since 2001. Most of the present generation of students are computer literate. Thailand was ranked 74th out of 100 countries globally for English proficiency.[111] Thailand has the second highest number of English-medium private international schools in Southeast Asian Nations, according to the International School Consultancy Group 181 schools around the country in 2017 compared to just 10 international schools for expatriate children in 1992.[106]
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Students in ethnic minority areas score consistently lower in standardised national and international tests.[112][113][114]
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This is likely due to unequal allocation of educational resources, weak teacher training, poverty, and low Thai language skill, the language of the tests.[112]
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[115]
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[116]
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Extensive nationwide IQ tests were administered to 72,780 Thai students from December 2010 to January 2011. The average IQ was found to be 98.59, which is higher than previous studies have found. IQ levels were found to be inconsistent throughout the country, with the lowest average of 88.07 found in the southern region of Narathiwat Province and the highest average of 108.91 reported in Nonthaburi Province. The Ministry of Public Health blames the discrepancies on iodine deficiency, and as of 2011[update] steps were being taken to require that iodine be added to table salt, a practice common in many Western countries.[117]
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In 2013, the Ministry of Information and Communication Technology announced that 27,231 schools would receive classroom-level access to high-speed internet.[118]
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In modern times, Thai scientists have made many significant contributions in various fields of study. For example, In chemistry, Krisana Kraisintu as known as the "Gypsy pharmacist".[119] She developed one of the first generic ARV fixed-dose combinations and dedicated her life to making medicines more affordable and accessible. Her efforts have saved countless lives in Africa,GPO-VIR has now been chosen by World Health Organization as the first regimen treatment for HIV/AIDS patients in poor countries.[120] In Thailand, this drug (GPO-VIR) is used in the national HIV/AIDS treatment programme, making it free of charge for 100,000 patients.[121] while Pongrama Ramasoota, He discoveries production of therapeutic human monoclonal antibodies against dengue virus and the world's first Dengue fever medication, include DNA vaccine development for dengue and Canine parvovirus.[122]
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Thailand has also made significant advances technology in the development of Medical Robotics. Medical robots have been used and promoted in Thailand in many areas, including surgery, diagnosis, rehabilitation and services.[123] and their use has been increasing. such as, an elderly care robot made by Thai manufacturer that Japanese nursing homes are widely using.[124] In surgery, back in 2019, The Medical Services Department has unveiled Thailand's robot created to help surgeons in brain surgery on patients afflicted with epilepsy.[125] back in 2017, Ramathibodi Hospital, a leading government hospital in Bangkok and a reputable medical school, successfully performed the first robot-assisted brain surgery in Asia.[126] For rehabilitation and therapy robots, were developed to help patients with arm and leg injuries perform practiced movements aided by the robots is the first prize winner of the i-MEDBOT Innovation Contest 2018 held by Thailand Center of Excellence for Life Sciences (TCELS).[127]
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According to the UNESCO Institute for Statistics, Thailand devoted 1% of its GDP to science research and development in 2017.[128] Between 2014 and 2016, Research and development workforce in Thailand increased from 84,216 people to 112,386 people.[129] The Thai government is developing new
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growth hubs by starting with the
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Eastern Economic Corridor of Innovation (EECi) to accelerating human resource and research development.[130] The National Science and Technology Development Agency is an agency of the government of Thailand which supports research in science and technology and its application in the Thai economy.[131]
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Thailand is an emerging economy and is considered a newly industrialised country. Thailand had a 2017 GDP of US$1.236 trillion (on a purchasing power parity basis).[137] Thailand is the 2nd largest economy in Southeast Asia after Indonesia. Thailand ranks midway in the wealth spread in Southeast Asia as it is the 4th richest nation according to GDP per capita, after Singapore, Brunei, and Malaysia.
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Thailand functions as an anchor economy for the neighbouring developing economies of Laos, Myanmar, and Cambodia. In the third quarter of 2014, the unemployment rate in Thailand stood at 0.84% according to Thailand's National Economic and Social Development Board (NESDB).[138]
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Thailand experienced the world's highest economic growth rate from 1985 to 1996 – averaging 12.4% annually. In 1997 increased pressure on the baht, a year in which the economy contracted by 1.9%, led to a crisis that uncovered financial sector weaknesses and forced the Chavalit Yongchaiyudh administration to float the currency. Prime Minister Chavalit Yongchaiyudh was forced to resign after his cabinet came under fire for its slow response to the economic crisis. The baht was pegged at 25 to the US dollar from 1978 to 1997. The baht reached its lowest point of 56 to the US dollar in January 1998 and the economy contracted by 10.8% that year, triggering the Asian financial crisis.
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Thailand's economy started to recover in 1999, expanding 4.2–4.4% in 2000, thanks largely to strong exports. Growth (2.2%) was dampened by the softening of the global economy in 2001, but picked up in the subsequent years owing to strong growth in Asia, a relatively weak baht encouraging exports, and increased domestic spending as a result of several mega projects and incentives of Prime Minister Thaksin Shinawatra, known as Thaksinomics. Growth in 2002, 2003, and 2004 was 5–7% annually.
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Growth in 2005, 2006, and 2007 hovered around 4–5%. Due both to the weakening of the US dollar and an increasingly strong Thai currency, by March 2008 the dollar was hovering around the 33 baht mark. While Thaksinomics has received criticism, official economic data reveals that between 2001 and 2011, Isan's GDP per capita more than doubled to US$1,475, while, over the same period, GDP in the Bangkok area increased from US$7,900 to nearly US$13,000.[139]
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With the instability surrounding major 2010 protests, the GDP growth of Thailand settled at around 4–5%, from highs of 5–7% under the previous civilian administration. Political uncertainty was identified as the primary cause of a decline in investor and consumer confidence. The IMF predicted that the Thai economy would rebound strongly from the low 0.1% GDP growth in 2011, to 5.5% in 2012 and then 7.5% in 2013, due to the monetary policy of the Bank of Thailand, as well as a package of fiscal stimulus measures introduced by the former Yingluck Shinawatra government.[140]
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Following the Thai military coup of 22 May 2014. In 2017, Concluded with information on the Thai economy's grew an inflation-adjusted 3.9%, up from 3.3% in 2016, marking its fastest expansion since 2012.[141]
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Thais have median wealth per one adult person of $1,469 in 2016,[142]:98 increasing from $605 in 2010.[142]:34 In 2016, Thailand was ranked 87th in Human Development Index, and 70th in the inequality-adjusted HDI.[143]
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In 2017, Thailand's median household income was ฿26,946 per month.[144]:1 Top quintile households had a 45.0% share of all income, while bottom quintile households had 7.1%.[144]:4 There were 26.9 million persons who had the bottom 40% of income earning less than ฿5,344 per person per month.[145]:5 During 2013–2014 Thai political crisis, a survey found that anti-government PDRC mostly (32%) had a monthly income of more than ฿50,000, while pro-government UDD mostly (27%) had between ฿10,000 and ฿20,000.[146]:7
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In 2014, Credit Suisse reported that Thailand was the world's third most unequal country, behind Russia and India.[147] Top 10% richest held 79% of the country's asset.[147] Top 1% richest held 58% worth of the economy.[147] Thai 50 richest families had a total net worth accounting to 30% of GDP.[147]
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In 2016, 5.81 million people lived in poverty, or 11.6 million people (17.2% of population) if "near poor" is included.[145]:1 Proportion of the poor relative to total population in each region was 12.96% in the Northeast, 12.35% in the South, and 9.83% in the North.[145]:2 In 2017, there were 14 million people who applied for social welfare (yearly income of less than ฿100,000 was required).[147] At the end of 2017, Thailand's total household debt was ฿11.76 trillion.[134]:5 In 2010, 3% of all household were bankrupt.[136]:5 In 2016, there were estimated 30,000 homeless persons in the country.[148]
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The economy of Thailand is heavily export-dependent, with exports accounting for more than two-thirds of gross domestic product (GDP). Thailand exports over US$105 billion worth of goods and services annually.[1] Major exports include cars, computers, electrical appliances, rice, textiles and footwear, fishery products, rubber, and jewellery.[1]
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Substantial industries include electric appliances, components, computer components, and vehicles. Thailand's recovery from the 1997–1998 Asian financial crisis depended mainly on exports, among various other factors. As of 2012[update], the Thai automotive industry was the largest in Southeast Asia and the 9th largest in the world.[149][150][151] The Thailand industry has an annual output of near 1.5 million vehicles, mostly commercial vehicles.[151]
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Most of the vehicles built in Thailand are developed and licensed by foreign producers, mainly Japanese and American. The Thai car industry takes advantage of the ASEAN Free Trade Area (AFTA) to find a market for many of its products. Eight manufacturers, five Japanese, two US, and Tata of India, produce pick-up trucks in Thailand.[152] As of 2012, Thailand was the second largest consumer of pick-up trucks in the world, after the US.[153] In 2014, pick-ups accounted for 42% of all new vehicle sales in Thailand.[152]
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Tourism makes up about 6% of the country's economy. Thailand was the most visited country in Southeast Asia in 2013, according to the World Tourism Organisation. Estimates of tourism receipts directly contributing to the Thai GDP of 12 trillion baht range from 9 percent (1 trillion baht) (2013) to 16 percent.[154] When including the indirect effects of tourism, it is said to account for 20.2 percent (2.4 trillion baht) of Thailand's GDP.[155]:1
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Asian tourists primarily visit Thailand for Bangkok and the historical, natural, and cultural sights in its vicinity. Western tourists not only visit Bangkok and surroundings, but in addition many travel to the southern beaches and islands. The north is the chief destination for trekking and adventure travel with its diverse ethnic minority groups and forested mountains. The region hosting the fewest tourists is Isan. To accommodate foreign visitors, a separate tourism police with offices were set up in the major tourist areas and an emergency telephone number.[156]
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Thailand ranks 5th biggest medical tourism destination of inbound medical tourism spending, according to World Travel and Tourism Council, attracting over 2.5 million visitors in 2018.[157] The country is also Asia's number one.[158] The country is popular for the growing practice of sex reassignment surgery (SRS) and cosmetic surgery. In 2010–2012, more than 90% of the visitors travelled to Thailand for SRS.[159]
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Prostitution in Thailand and sex tourism also form a de facto part of the economy. Campaigns promote Thailand as exotic to attract tourists.[160] One estimate published in 2003 placed the trade at US$4.3 billion per year or about 3% of the Thai economy.[161] It is believed that at least 10% of tourist dollars are spent on the sex trade.[162]
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Forty-nine per cent of Thailand's labour force is employed in agriculture.[163] This is down from 70% in 1980.[163] Rice is the most important crop in the country and Thailand had long been the world's leading exporter of rice, until recently falling behind both India and Vietnam.[164] Thailand has the highest percentage of arable land, 27.25%, of any nation in the Greater Mekong Subregion.[165] About 55% of the arable land area is used for rice production.[166]
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Agriculture has been experiencing a transition from labour-intensive and transitional methods to a more industrialised and competitive sector.[163] Between 1962 and 1983, the agricultural sector grew by 4.1% per year on average and continued to grow at 2.2% between 1983 and 2007.[163] The relative contribution of agriculture to GDP has declined while exports of goods and services have increased.
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Furthermore, access to biocapacity in Thailand is lower than world average. In 2016, Thailand had 1.2 global hectares[167] of biocapacity per person within its territory, a little less than world average of 1.6 global hectares per person.[168] In contrast, in 2016, they used 2.5 global hectares of biocapacity – their ecological footprint of consumption. This means they use about twice as much biocapacity as Thailand contains. As a result, Thailand is running a biocapacity deficit.[167]
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The State Railway of Thailand (SRT) operates all of Thailand's national rail lines. Bangkok Railway Station (Hua Lamphong Station) is the main terminus of all routes. Phahonyothin and ICD Lat Krabang are the main freight terminals. As of 2017[update] SRT had 4,507 km (2,801 mi) of track, all of it meter gauge except the Airport Link. Nearly all is single-track (4,097 km), although some important sections around Bangkok are double (303 km or 188 mi) or triple-tracked (107 km or 66 mi) and there are plans to extend this.[169] Rail transport in Bangkok includes long-distance services, and some daily commuter trains running from and to the outskirts of the city during the rush hour, but passenger numbers have remained low. There are also three rapid transit rail systems in the capital.
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Thailand has 390,000 km (242,335 miles) of highways.[170] According to the BBC Thailand has 462,133 roads and many multi-lane highways. As of 2017[update] Thailand has 37 million registered vehicles, 20 million of them motorbikes. A number of undivided two-lane highways have been converted into divided four-lane highways. A Bangkok – Chon Buri motorway (Route 7) now links to the new airport and Eastern Seaboard. There are 4,125 public vans operating on 114 routes from Bangkok alone.[171] Other forms of road transport includes tuk-tuks, taxis—as of November 2018, Thailand has 80,647 registered taxis nationwide[172]—vans (minibus), motorbike taxis and songthaews.
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As of 2012, Thailand had 103 airports with 63 paved runways, in addition to 6 heliports. The busiest airport in the county is Bangkok's Suvarnabhumi Airport.
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75% of Thailand's electrical generation is powered by natural gas in 2014.[173] Coal-fired power plants produce an additional 20% of electricity, with the remainder coming from biomass, hydro, and biogas.[173]
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Thailand produces roughly one-third of the oil it consumes. It is the second largest importer of oil in SE Asia. Thailand is a large producer of natural gas, with reserves of at least 10 trillion cubic feet. After Indonesia, it is the largest coal producer in SE Asia, but must import additional coal to meet domestic demand.
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Thailand has a diverse and robust informal labour sector—in 2012, it was estimated that informal workers comprised 62.6% of the Thai workforce. The Ministry of Labour defines informal workers to be individuals who work in informal economies and do not have employee status under a given country's Labour Protection Act (LPA). The informal sector in Thailand has grown significantly over the past 60 years over the course of Thailand's gradual transition from an agriculture-based economy to becoming more industrialised and service-oriented.[174] Between 1993 and 1995, ten percent of the Thai labour force moved from the agricultural sector to urban and industrial jobs, especially in the manufacturing sector. It is estimated that between 1988 and 1995, the number of factory workers in the country doubled from two to four million, as Thailand's GDP tripled.[175] While the Asian Financial Crisis that followed in 1997 hit the Thai economy hard, the industrial sector continued to expand under widespread deregulation, as Thailand was mandated to adopt a range of structural adjustment reforms upon receiving funding from the IMF and World Bank. These reforms implemented an agenda of increased privatisation and trade liberalisation in the country, and decreased federal subsidisation of public goods and utilities, agricultural price supports, and regulations on fair wages and labour conditions.[176] These changes put further pressure on the agricultural sector, and prompted continued migration from the rural countryside to the growing cities. Many migrant farmers found work in Thailand's growing manufacturing industry, and took jobs in sweatshops and factories with few labour regulations and often exploitative conditions.[177]
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Those that could not find formal factory work, including illegal migrants and the families of rural Thai migrants that followed their relatives to the urban centres, turned to the informal sector to provide the extra support needed for survival—under the widespread regulation imposed by the structural adjustment programs, one family member working in a factory or sweatshop made very little. Scholars argue that the economic consequences and social costs of Thailand's labour reforms in the wake of the 1997 Asian Financial Crisis fell on individuals and families rather than the state. This can be described as the "externalisation of market risk", meaning that as the country's labour market became increasingly deregulated, the burden and responsibility of providing an adequate livelihood shifted from employers and the state to the workers themselves, whose families had to find jobs in the informal sector to make up for the losses and subsidise the wages being made by their relatives in the formal sector. The weight of these economic changes hit migrants and the urban poor especially hard, and the informal sector expanded rapidly as a result.[176]
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Today, informal labour in Thailand is typically broken down into three main groups: subcontracted/self employed/home-based workers, service workers (including those that are employed in restaurants, as street vendors, masseuses, taxi drivers, and as domestic workers), and agricultural workers. Not included in these categories are those that work in entertainment, nightlife, and the sex industry. Individuals employed in these facets of the informal labour sector face additional vulnerabilities, including recruitment into circles of sexual exploitation and human trafficking.[174]
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In general, education levels are low in the informal sector. A 2012 study found that 64% of informal workers had not completed education beyond primary school. Many informal workers are also migrants, only some of which have legal status in the country. Education and citizenship are two main barriers to entry for those looking to work in formal industries, and enjoy the labour protections and social security benefits that come along with formal employment. Because the informal labour sector is not recognised under the Labour Protection Act (LPA), informal workers are much more vulnerable labour to exploitation and unsafe working conditions than those employed in more formal and federally recognised industries. While some Thai labour laws provide minimal protections to domestic and agricultural workers, they are often weak and difficult to enforce. Furthermore, Thai social security policies fail to protect against the risks many informal workers face, including workplace accidents and compensation as well as unemployment and retirement insurance. Many informal workers are not legally contracted for their employment, and many do not make a living wage.[174] As a result, labour trafficking is common in the region, affecting children and adults, men and women, and migrants and Thai citizens alike.
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Thailand had a population of 66,558,935 as of 2019.[4] Thailand's population is largely rural, concentrated in the rice-growing areas of the central, northeastern and northern regions. About 45.7% of Thailand's population lived in urban areas as of 2010[update], concentrated mostly in and around the Bangkok Metropolitan Area.
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Thailand's government-sponsored family planning program resulted in a dramatic decline in population growth from 3.1% in 1960 to around 0.4% today. In 1970, an average of 5.7 people lived in a Thai household. At the time of the 2010 census, the average Thai household size was 3.2 people.
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Thai nationals make up the majority of Thailand's population, 95.9% in 2010. The remaining 4.1% of the population are Burmese (2.0%), others 1.3%, and unspecified 0.9%.[1]
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According to the Royal Thai Government's 2011 Country Report to the UN Committee responsible for the International Convention for the Elimination of All Forms of Racial Discrimination, available from the Department of Rights and Liberties Promotion of the Thai Ministry of Justice,:3 62 ethnic communities are officially recognised in Thailand. Twenty million Central Thai (together with approximately 650,000 Khorat Thai) make up approximately 20,650,000 (34.1 percent) of the nation's population of 60,544,937[180] at the time of completion of the Mahidol University Ethnolinguistic Maps of Thailand data (1997).[181]
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The 2011 Thailand Country Report provides population numbers for mountain peoples ('hill tribes') and ethnic communities in the Northeast and is explicit about its main reliance on the Mahidol University Ethnolinguistic Maps of Thailand data.[181] Thus, though over 3.288 million people in the Northeast alone could not be categorised, the population and percentages of other ethnic communities circa 1997 are known for all of Thailand and constitute minimum populations. In descending order, the largest (equal to or greater than 400,000) are a) 15,080,000 Lao (24.9 percent) consisting of the Thai Lao (14 million) and other smaller Lao groups, namely the Thai Loei (400–500,000), Lao Lom (350,000), Lao Wiang/Klang (200,000), Lao Khrang (90,000), Lao Ngaew (30,000), and Lao Ti (10,000; b) six million Khon Muang (9.9 percent, also called Northern Thais); c) 4.5 million Pak Tai (7.5 percent, also called Southern Thais); d) 1.4 million Khmer Leu (2.3 percent, also called Northern Khmer); e) 900,000 Malay (1.5%); f) 500,000 Nyaw (0.8 percent); g) 470,000 Phu Thai (0.8 percent); h) 400,000 Kuy/Kuay (also known as Suay) (0.7 percent), and i) 350,000 Karen (0.6 percent).:7–13 Thai Chinese, those of significant Chinese heritage, are 14% of the population, while Thais with partial Chinese ancestry comprise up to 40% of the population.[182] Thai Malays represent 3% of the population, with the remainder consisting of Mons, Khmers and various "hill tribes". The country's official language is Thai and the primary religion is Theravada Buddhism, which is practised by around 95% of the population.
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Increasing numbers of migrants from neighbouring Myanmar, Laos, and Cambodia, as well as from Nepal and India, have pushed the total number of non-national residents to around 3.5 million as of 2009[update], up from an estimated 2 million in 2008, and about 1.3 million in 2000.[183] Some 41,000 Britons and 20,000 Australians live in Thailand.[184][185]
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The official language of Thailand is Thai, a Kra–Dai language closely related to Lao, Shan in Myanmar, and numerous smaller languages spoken in an arc from Hainan and Yunnan south to the Chinese border. It is the principal language of education and government and spoken throughout the country. The standard is based on the dialect of the central Thai people, and it is written in the Thai alphabet, an abugida script that evolved from the Khmer alphabet.
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Sixty-two languages were recognised by the Royal Thai Government in the 2011 Country Report to the UN Committee responsible for the International Convention for the Elimination of All Forms of Racial Discrimination, which employed an ethnolinguistic approach and is available from the Department of Rights and Liberties Promotion of the Thai Ministry of Justice.:3 Southern Thai is spoken in the southern provinces, and Northern Thai is spoken in the provinces that were formerly part of the independent kingdom of Lan Na. For the purposes of the national census, which does not recognise all 62 languages recognised by the Royal Thai Government in the 2011 Country Report, four dialects of Thai exist; these partly coincide with regional designations.
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The largest of Thailand's minority languages is the Lao dialect of Isan spoken in the northeastern provinces. Although sometimes considered a Thai dialect, it is a Lao dialect, and the region where it is traditionally spoken was historically part of the Lao kingdom of Lan Xang.[citation needed] In the far south, Kelantan-Pattani Malay is the primary language of Malay Muslims. Varieties of Chinese are also spoken by the large Thai Chinese population, with the Teochew dialect best-represented.
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Numerous tribal languages are also spoken, including many Austroasiatic languages such as Mon, Khmer, Viet, Mlabri and Orang Asli; Austronesian languages such as Cham and Moken; Sino-Tibetan languages like Lawa, Akha, and Karen; and other Tai languages such as Tai Yo, Phu Thai, and Saek. Hmong is a member of the Hmong–Mien languages, which is now regarded as a language family of its own.
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English is a mandatory school subject, but the number of fluent speakers remains low, especially outside cities.
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Thailand's prevalent religion is Theravada Buddhism, which is an integral part of Thai identity and culture. Active participation in Buddhism is among the highest in the world. According to the 2000 census, 94.6% and 93.58% in 2010 of the country's population self-identified as Buddhists of the Theravada tradition. Muslims constitute the second largest religious group in Thailand, comprising 4.29% of the population in 2015.[187]
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Islam is concentrated mostly in the country's southernmost provinces: Pattani, Yala, Satun, Narathiwat, and part of Songkhla Chumphon, which are predominantly Malay, most of whom are Sunni Muslims. Christians represented 1.17% (2015) of the population in 2015, with the remaining population consisting of Hindus and Sikhs, who live mostly in the country's cities. There is also a small but historically significant Jewish community in Thailand dating back to the 17th century.
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The constitution does not name official state religion, and provides for freedom of religion. Even the authority formally does not register new religious groups that have not been accepted and limit the number of missionaries, unregistered religious organisations as well as missionaries who are allowed to operate freely. There have been no widespread reports of societal abuses or discrimination based on religious belief or practice.[188]
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Thailand ranks world's 6th, and Asia's 1st in the 2019 Global Health Security Index of global health security capabilities in 195 countries,[189] making it the only developing country on the world's top ten. Thailand had 62 hospitals accredited by Joint Commission International.[190] In 2002, Bumrungrad became the first hospital in Asia to meet the standard.
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Health and medical care is overseen by the Ministry of Public Health (MOPH), along with several other non-ministerial government agencies, with total national expenditures on health amounting to 4.3 percent of GDP in 2009. Non-communicable diseases form the major burden of morbidity and mortality, while infectious diseases including malaria and tuberculosis, as well as traffic accidents, are also important public health issues. The current Minister for Public Health is Anutin Charnvirakul.
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In December 2018 the interim parliament voted to legalise the use of cannabis for medical reasons. Recreational use remained unlawful. The National Legislative Assembly had 166 votes in favour of the amendment to the Narcotics Bill, while there were no nay votes and 13 abstentions. The vote makes Thailand the first Southeast Asian country to allow the use of medical cannabis.[191]
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Thai culture and traditions incorporate a great deal of influence from India, China, Cambodia, and the rest of Southeast Asia. Thailand's national religion, Theravada Buddhism, is central to modern Thai identity. Thai Buddhism has evolved over time to include many regional beliefs originating from Hinduism, animism, as well as ancestor worship. The official calendar in Thailand is based on the Eastern version of the Buddhist Era (BE). Thai identity today is a social construct of Phibun regime in 1940s.
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Several ethnic groups mediated change between their traditional local culture, national Thai, and global cultural influences. Overseas Chinese also form a significant part of Thai society, particularly in and around Bangkok. Their successful integration into Thai society has allowed them to hold positions of economic and political power. Thai Chinese businesses prosper as part of the larger bamboo network.[192]
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Wai is traditional Thai greeting, and is generally offered first by person who is younger or lower in social status and position. Respects for elderly (by age, position, monks, or certain professions) is Thai mores. As with other Asian cultures, respect towards ancestors is an essential part of Thai spiritual practice. Thais have strong sense of social hierarchy, reflecting in many classes of honorifics. Seniority is paramount in Thai culture. Elders have by tradition ruled in family decisions or ceremonies. Older siblings have duties to younger ones.
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Thais have a strong sense of hospitality and generosity.
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Taboos in Thailand include touching someone's head or pointing with the feet, as the head is considered the most sacred and the foot the lowest part of the body.
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The origins of Thai art were very much influenced by Buddhist art and by scenes from the Indian epics. Traditional Thai sculpture almost exclusively depicts images of the Buddha, being very similar with the other styles from Southeast Asia. Traditional Thai paintings usually consist of book illustrations, and painted ornamentation of buildings such as palaces and temples. Thai art was influenced by indigenous civilisations of the Mon and other civilisations. By the Sukothai and Ayutthaya period, thai had developed into its own unique style and was later further influenced by the other Asian styles, mostly by Sri Lankan and Chinese. Thai sculpture and painting, and the royal courts provided patronage, erecting temples and other religious shrines as acts of merit or to commemorate important events.[193]
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Traditional Thai paintings showed subjects in two dimensions without perspective. The size of each element in the picture reflected its degree of importance. The primary technique of composition is that of apportioning areas: the main elements are isolated from each other by space transformers. This eliminated the intermediate ground, which would otherwise imply perspective. Perspective was introduced only as a result of Western influence in the mid-19th century. Monk artist Khrua In Khong is well known as the first artist to introduce linear perspective to Thai traditional art.[194]
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The most frequent narrative subjects for paintings were or are: the Jataka stories, episodes from the life of the Buddha, the Buddhist heavens and hells, themes derived from the Thai versions of the Ramayana and Mahabharata, not to mention scenes of daily life. Some of the scenes are influenced by Thai folklore instead of following strict Buddhist iconography.[193]
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Architecture is the preeminent medium of the country's cultural legacy and reflects both the challenges of living in Thailand's sometimes extreme climate as well as, historically, the importance of architecture to the Thai people's sense of community and religious beliefs. Influenced by the architectural traditions of many of Thailand's neighbours, it has also developed significant regional variation within its vernacular and religious buildings.
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The Ayutthaya Kingdom movement, which went from approximately 1350 to 1767, was one of the most fruitful and creative periods in Thai architecture The identity of architecture in Ayutthaya period is designed to display might and riches so it has great size and appearance. The temples in Ayutthaya seldom built eaves stretching from the masterhead. The dominant feature of this style is sunlight shining into buildings. During the latter part of the Ayutthaya period, architecture was regarded as a peak achievement that responded to the requirements of people and expressed the gracefulness of Thainess.[195]
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Buddhist temples in Thailand are known as "wats", from the Pāḷi vāṭa, meaning an enclosure. A temple has an enclosing wall that divides it from the secular world. Wat architecture has seen many changes in Thailand in the course of history. Although there are many differences in layout and style, they all adhere to the same principles.[196]
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Thai literature has had a long history. Even before the establishment of the Sukhothai Kingdom there existed oral and written works.
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During the Sukhothai, Most literary works were written in simple prose with certain alliteration schemes. Major works include King Ram Khamhaeng Inscription. King Ram Khamhaeng's Stone Inscription is considered the first Thai literary work in Thai script. It gives an account of the life of King Ramkhamhaeng the Great, the way of life of Thai people in general, laws, religion, economic and political stability. Trai Phum Phra Ruang, was written in 1345 by King Maha Thammaracha I, the fifth king of Sukhothai. It expounds Buddhist philosophy based on a profound and extensive study with reference to over 30 sacred texts. The work could be considered the nation's first piece of research dissertation. It was written in beautiful prose rich in allusions and imagery. It is a treatise on Buddhist cosmology, ethics, biology and belief system.[197]
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During the Ayutthaya, The period produced a variety of forms on diverse subjects. New poetic forms were created, with different rhyme schemes and metres. It is common to find a combination of different poetic forms in one poetic work. Lilit Yuan Phai is a narrative poem describing the war between King Borommatrailokkanat of Ayutthaya and Prince Tilokkarat of Lan Na. One of the most beautiful literary works is Kap He Ruea composed by Prince Thammathibet comparing the scenic beauty to that of his beloved lady on a boat journey in the nirat tradition. Traditionally, the verse is sung during the colourful royal barge procession. It has been the model for subsequent poets to emulate. The same prince also composed the greatly admired Kap Ho Khlong on the Visit to Than Thongdaeng and Kap Ho Khlong Nirat Phrabat.[198]
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Despite its short period of 15 years, Thon Buri Period produced Ramakian, a verse drama to which King Taksin the Great contributed his poetic talent. The revival of literature at this time is remarkable since the country had not quite recovered from the aftermath of war. Some poets who later became a major force in the early Rattanakosin Period had already begun writing at this time.
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During the 18th century Rattanakosin Period. After sporadic fighting at the beginning of the period, the country gradually returned to normal. It is only natural that many of the early Rattanakosin works should deal with war and military strategy. Some examples are Nirat Rop Phama Thi Tha Din Daeng, Phleng Yao Rop Phama Thi Nakhon Si Thammarat.In the performing arts, perhaps the most important dramatic achievement is the complete work of Ramakian by King Rama I. In addition, There were also verse recitals with musical accompaniment, such as Mahori telling the story of Kaki, Sepha relating the story of Khun Chang Khun Phaen. Other recitals include Sri Thanonchai. The most important Thai poet in this period was Sunthorn Phu (สุนทรภู่) (1786–1855), widely known as "the bard of Rattanakosin" (Thai: กวีเอกแห่งกรุงรัตนโกสินทร์). Sunthorn Phu is best known for his epic poem Phra Aphai Mani (Thai: พระอภัยมณี), which he started in 1822 (while in jail) and finished in 1844. Phra Aphai Mani is a versified fantasy-adventure novel, a genre of Siamese literature known as nithan kham klon (Thai: นิทานคำกลอน).[198]
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Aside from folk and regional dances (southern Thailand's Menora (dance) and Ramwong, for example), the two major forms of Thai classical dance drama are Khon and Lakhon nai. In the beginning, both were exclusively court entertainments and it was not until much later that a popular style of dance theatre, likay, evolved as a diversion for common folk who had no access to royal performances.[199]
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Folk dance forms include dance theater forms like likay, numerous regional dances (ram), the ritual dance ram muay, and homage to the teacher, wai khru. Both ram muay and wai khru take place before all traditional muay Thai matches. The wai is also an annual ceremony performed by Thai classical dance groups to honor their artistic ancestors.
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Thai classical music is synonymous with those stylized court ensembles and repertoires that emerged in their present form within the royal centers of Central Thailand some 800 years ago. These ensembles, while being influenced by older practices are today uniquely Thai expressions. While the three primary classical ensembles, the Piphat, Khrueang sai and Mahori differ in significant ways, they all share a basic instrumentation and theoretical approach. Each employs small ching hand cymbals and krap wooden sticks to mark the primary beat reference. Thai classical music has had a wide influence on the musical traditions of neighboring countries. The traditional music of Myanmar was strongly influenced by the Thai music repertoire, called Yodaya (ယိုးဒယား), which was brought over from the Ayutthaya Kingdom. As Siam expanded its political and cultural influence to Laos and Cambodia during the early Rattanakosin period, its music was quickly absorbed by the Cambodia and Lao courts.
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Thai films are exported and exhibited in Southeast Asia.[200] Thai cinema has developed its own unique identity and now being internationally recognized for their culture-driven.[201] Films such as Ong-Bak: Muay Thai Warrior (2003) and 'Tom-Yum-Goong (2005), starred Tony Jaa, feature distinctive aspects of Thai martial arts "Muay Thai".
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Thai horror has always had a significant cult following, unique take on tales from beyond the grave. More recently, horror films such as Shutter (2004), was one of the best-known Thai horror movies and recognized worldwide.[202] Other examples include The Unseeable (2006), Alone (2007), Body (2007), Coming Soon (2008), 4bia (2008), Phobia 2 (2009), Ladda Land (2011), Pee Mak (2013), and The Promise (2017).
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Thai heist thriller film Bad Genius (2017), was one of the most internationally successful Thai film, It broke Thai film earning records in several Asian countries,[203][204] Bad Genius won in 12 categories at the 27th Suphannahong National Film Awards, and also won the Jury Award at the 16th New York Asian Film Festival with a worldwide collection of more than $42 million.[205]
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Thailand television dramas, known as Lakorn, Lakorn have become popular in Thailand and its neighbors.[206] Many dramas tend to have a romantic focus, such as Khluen Chiwit, U-Prince, Ugly Duckling, The Crown Princess and teen dramas television series, such as 2gether: The Series, The Gifted, Girl From Nowhere, Hormones: The Series.
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The Entertainment industries (film and television) are estimated to have directly contributed $2.1 billion in gross domestic product (GDP) to the Thai economy in 2011. They also directly supported 86,600 jobs.[207] Amongst several Dance-pop artists who have made internationally successful can be mentioned "Lisa" Lalisa Manoban[208] and Tata Young.
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Thai cuisine is one of the most popular in the world.[209][210] Thai food blends five fundamental tastes: sweet, spicy, sour, bitter, and salty. The herbs and spices most used in Thai cooking themselves have medicinal qualities such as garlic, lemongrass, Kaffir lime, galangal, turmeric, coriander, coconut milk.[211] Each region of Thailand has its specialities: kaeng khiao wan (green curry) in the central region, som tam (green papaya salad) in the northeast, khao soi in the north, Massaman curry in the south.
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In 2017, seven Thai dishes appeared on a list of the "World's 50 Best Foods"— an online poll of worldwide by CNN Travel. Thailand had more dishes on the list than any other country. They were: tom yam goong (4th), pad Thai (5th), som tam (6th), Massaman curry (10th), green curry (19th), Thai fried rice (24th) and mu nam tok (36th).[212]
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The staple food in Thailand is rice, particularly jasmine rice (also known as hom Mali) which forms part of almost every meal. Thailand is a leading exporter of rice, and Thais consume over 100 kg of milled rice per person per year.[213]
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Thailand generally uses the metric system, but traditional units of measurement for land area are used, and imperial units of measurement are occasionally used for building materials, such as wood and plumbing fixtures. Years are numbered as B.E. (Buddhist Era) in educational settings, civil service, government, contracts, and newspaper datelines. However, in banking, and increasingly in industry and commerce, standard Western year (Christian or Common Era) counting is the standard practice.[214]
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Muay Thai (Thai: มวยไทย, RTGS: Muai Thai, [muaj tʰaj], lit. "Thai boxing") is a combat sport of Thailand that uses stand-up striking along with various clinching techniques. Muay Thai became widespread internationally in the late-20th to 21st century, when Westernized practitioners from Thailand began competing in kickboxing and mixed rules matches as well as matches under muay Thai rules around the world, Famous practitioners such as Buakaw Banchamek, Samart Payakaroon, Dieselnoi Chor Thanasukarn and Apidej Sit-Hirun. Buakaw Banchamek has probably brought more international interest in Muay Thai than any other Muay Thai fighters ever had.[215]
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Association football has overtaken muay Thai as the most widely followed sport in contemporary Thai society. Thailand national football team has played the AFC Asian Cup six times and reached the semifinals in 1972. The country has hosted the Asian Cup twice, in 1972 and in 2007. The 2007 edition was co-hosted together with Indonesia, Malaysia and Vietnam. It is not uncommon to see Thais cheering their favourite English Premier League teams on television and walking around in replica kit. Another widely enjoyed pastime, and once a competitive sport, is kite flying.
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Volleyball is rapidly growing as one of the most popular sports. The women's team has often participated in the World Championship, World Cup, and World Grand Prix Asian Championship. They have won the Asian Championship twice and Asian Cup once. By the success of the women's team, the men team has been growing as well.
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Takraw (Thai: ตะกร้อ) is a sport native to Thailand, in which the players hit a rattan ball and are only allowed to use their feet, knees, chest, and head to touch the ball. Sepak takraw is a form of this sport which is similar to volleyball. The players must volley a ball over a net and force it to hit the ground on the opponent's side. It is also a popular sport in other countries in Southeast Asia. A rather similar game but played only with the feet is buka ball.
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Snooker has enjoyed increasing popularity in Thailand in recent years, with interest in the game being stimulated by the success of Thai snooker player James Wattana in the 1990s.[216] Other notable players produced by the country include Ratchayothin Yotharuck, Noppon Saengkham and Dechawat Poomjaeng.[217]
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Rugby is also a growing sport in Thailand with the Thailand national rugby union team rising to be ranked 61st in the world.[218] Thailand became the first country in the world to host an international 80 welterweight rugby tournament in 2005.[219] The national domestic Thailand Rugby Union (TRU) competition includes several universities and services teams such as Chulalongkorn University, Mahasarakham University, Kasetsart University, Prince of Songkla University, Thammasat University, Rangsit University, the Thai Police, the Thai Army, the Thai Navy and the Royal Thai Air Force. Local sports clubs which also compete in the TRU include the British Club of Bangkok, the Southerners Sports Club (Bangkok) and the Royal Bangkok Sports Club.
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Thailand has been called the golf capital of Asia[220] as it is a popular destination for golf. The country attracts a large number of golfers from Japan, Korea, Singapore, South Africa, and Western countries who come to play golf in Thailand every year.[221] The growing popularity of golf, especially among the middle classes and immigrants, is evident as there are more than 200 world-class golf courses nationwide,[222] and some of them are chosen to host PGA and LPGA tournaments, such as Amata Spring Country Club, Alpine Golf and Sports Club, Thai Country Club, and Black Mountain Golf Club.
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Basketball is a growing sport in Thailand, especially on the professional sports club level. The Chang Thailand Slammers won the 2011 ASEAN Basketball League Championship.[223] The Thailand national basketball team had its most successful year at the 1966 Asian Games where it won the silver medal.[224]
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Other sports in Thailand are slowly growing as the country develops its sporting infrastructure. The success in sports like weightlifting and taekwondo at the last two summer Olympic Games has demonstrated that boxing is no longer the only medal option for Thailand.
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The well-known Lumpinee Boxing Stadium originally sited at Rama IV Road near Lumphini Park hosted its final Muay Thai boxing matches on 8 February 2014 after the venue first opened in December 1956. Managed by the Royal Thai Army, the stadium was officially selected for the purpose of muay Thai bouts following a competition that was staged on 15 March 1956. From 11 February 2014, the stadium will relocate to Ram Intra Road, due to the new venue's capacity to accommodate audiences of up to 3,500. Foreigners typically pay between 1,000–2,000 baht to view a match, with prices depending on the location of the seating.[225]
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Thammasat Stadium is a multi-purpose stadium in Bangkok. It is currently used mostly for football matches. The stadium holds 25,000. It is on Thammasat University's Rangsit campus. It was built for the 1998 Asian Games by construction firm Christiani and Nielsen, the same company that constructed the Democracy Monument in Bangkok.
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Rajamangala National Stadium is the biggest sporting arena in Thailand. It currently has a capacity of 65,000. It is in Bang Kapi, Bangkok. The stadium was built in 1998 for the 1998 Asian Games and is the home stadium of the Thailand national football team.
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Coordinates: 15°N 101°E / 15°N 101°E / 15; 101
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– in Asia (light green & dark grey)– in ASEAN (light green)
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Thailand,[a] officially the Kingdom of Thailand and formerly known as Siam,[b] is a country in Southeast Asia. Located at the centre of the Indochinese Peninsula, it is composed of 76 provinces, and covers an area of 513,120 square kilometres (198,120 sq mi), and a population of over 66 million people.[4] Thailand is the world's 50th-largest country by land area, and the 22nd-most-populous country in the world. The capital and largest city is Bangkok, a special administrative area. Thailand is bordered to the north by Myanmar and Laos, to the east by Laos and Cambodia, to the south by the Gulf of Thailand and Malaysia, and to the west by the Andaman Sea and the southern extremity of Myanmar. Its maritime boundaries include Vietnam in the Gulf of Thailand to the southeast, and Indonesia and India on the Andaman Sea to the southwest. Nominally, Thailand is a constitutional monarchy and parliamentary democracy; however, in recent history, its government has experienced multiple coups and periods of military dictatorships.
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Tai peoples migrated from southwestern China to mainland Southeast Asia from the 11th century; the oldest known mention of their presence in the region by the exonym Siamese dates to the 12th century. Various Indianised kingdoms such as the Mon kingdoms, Khmer Empire and Malay states ruled the region, competing with Thai states such as the Kingdoms of Ngoenyang, Sukhothai, Lan Na and Ayutthaya, which rivalled each other. Documented European contact began in 1511 with a Portuguese diplomatic mission to Ayutthaya, which became a regional power by the end of the 15th century. Ayutthaya reached its peak during cosmopolitan Narai's reign (1656–1688), gradually declining thereafter until being ultimately destroyed in the 1767 Burmese–Siamese War. Taksin (r. 1767–1782) quickly reunified the fragmented territory and established the short-lived Thonburi Kingdom. He was succeeded in 1782 by Buddha Yodfa Chulaloke (r. 1782–1809), the first monarch of the current Chakri dynasty.
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Through the 18th and 19th centuries, Siam faced imperialist pressure from France and the United Kingdom, including many unequal treaties with Western powers and forced concessions of territory; it nevertheless remained the only Southeast Asian country to avoid direct Western rule. Siamese system of government was centralized and transformed into modern unitary absolute monarchy in the reign of Chulalongkorn (r. 1868–1910). Siam joined World War I siding with the allies, a political decision to amend the unequal treaties. Following a bloodless revolution in 1932, Siam became a constitutional monarchy and changed its official name to "Thailand". Thailand was a satellite of Japan in World War II. In the late 1950s, a military coup under Field Marshal Sarit Thanarat revived the monarchy's historically influential role in politics. Thailand became a major ally of the United States, and played a key anti-communist role in the region as a member of the Southeast Asia Treaty Organization (SEATO). Apart from a brief period of parliamentary democracy in the mid-1970s, Thailand has periodically alternated between democracy and military rule. Since the 2000s, Thailand has been caught in a bitter political conflict between supporters and opponents of Thaksin Shinawatra, which culminated in two coups, most recently in 2014 and the establishment of its current and 20th constitution.
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Thailand is a founding member of the Association of Southeast Asian Nations (ASEAN) and remains a major ally of the United States.[12][13] Despite comparatively sporadic changes in leadership, it is considered a regional power in Southeast Asia and a middle power in global affairs.[14] With a high level of human development, the second-largest economy in Southeast Asia, and the 20th-largest in the world by PPP, Thailand is classified as a newly industrialized economy; manufacturing, agriculture, and tourism are leading sectors of the economy.[15][16]
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Thailand (/ˈtaɪlænd/ TY-land or /ˈtaɪlənd/ TY-lənd;[17] Thai: ประเทศไทย, RTGS: Prathet Thai, pronounced [pratʰêːt tʰaj] (listen)), officially the Kingdom of Thailand (Thai: ราชอาณาจักรไทย, RTGS: Ratcha-anachak Thai [râːtt͡ɕʰaʔaːnaːt͡ɕàk tʰaj] (listen), Chinese: 泰国), formerly known as Siam (Thai: สยาม, RTGS: Sayam [sajǎːm]), is a country at the centre of the Indochinese peninsula in Southeast Asia.
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The country has always been called Mueang Thai by its citizens. By outsiders, prior to 1949, it was usually known by the exonym Siam (Thai: สยาม RTGS: sayam, pronounced [sajǎːm], also spelled Siem, Syâm, or Syâma). The word Siam may have originated from Pali (suvaṇṇabhūmi, 'land of gold') or Sanskrit श्याम (śyāma, 'dark') or Mon ရာမည(rhmañña, 'stranger'). The names Shan and A-hom seem to be variants of the same word. The word Śyâma is possibly not its origin, but a learned and artificial distortion.[clarification needed][18] Another theory is the name derives from Chinese: "Ayutthaya emerged as a dominant centre in the late fourteenth century. The Chinese called this region Xian, which the Portuguese converted into Siam."[19]:8 A further possibility is that Mon-speaking peoples migrating south called themselves syem as do the autochthonous Mon-Khmer-speaking inhabitants of the Malay Peninsula.[citation needed]
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The signature of King Mongkut (r. 1851–1868) reads SPPM (Somdet Phra Poramenthra Maha) Mongkut Rex Siamensium (Mongkut King of the Siamese), giving the name Siam official status until 24 June 1939 when it was changed to "Thailand".[20] Thailand was renamed Siam from 1946 to 1948, after which it again reverted to "Thailand".
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According to George Cœdès, the word Thai (ไทย) means 'free man' in the Thai language, "differentiating the Thai from the natives encompassed in Thai society as serfs".[21]:197 A famous Thai scholar argued that Thai (ไท) simply means 'people' or 'human being', since his investigation shows that in some rural areas the word "Thai" was used instead of the usual Thai word khon (คน) for people.[22] According to Michel Ferlus, the ethnonyms Thai-Tai (or Thay-Tay) would have evolved from the etymon *k(ə)ri: 'human being' through the following chain: *kəri: > *kəli: > *kədi:/*kədaj > *di:/*daj > *dajA (Proto-Southwestern Tai) > tʰajA2 (in Siamese and Lao) or > tajA2 (in the other Southwestern and Central Tai languages classified by Li Fangkuei).[23] Michel Ferlus's work is based on some simple rules of phonetic change observable in the Sinosphere and studied for the most part by William H. Baxter (1992).[24]
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While Thai people will often refer to their country using the polite form prathet Thai (Thai: ประเทศไทย), they most commonly use the more colloquial term mueang Thai (Thai: เมืองไทย) or simply Thai; the word mueang, archaically referring to a city-state, is commonly used to refer to a city or town as the centre of a region. Ratcha Anachak Thai (Thai: ราชอาณาจักรไทย) means 'kingdom of Thailand' or 'kingdom of Thai'. Etymologically, its components are: ratcha (Sanskrit: राजन्, rājan, 'king, royal, realm'); -ana- (Pali āṇā 'authority, command, power', itself from the Sanskrit आज्ञा, ājñā, of the same meaning) -chak (from Sanskrit चक्र cakra- 'wheel', a symbol of power and rule). The Thai National Anthem (Thai: เพลงชาติ), written by Luang Saranupraphan during the patriotic 1930s, refers to the Thai nation as prathet Thai (Thai: ประเทศไทย). The first line of the national anthem is: prathet thai ruam lueat nuea chat chuea thai (Thai: ประเทศไทยรวมเลือดเนื้อชาติเชื้อไทย), 'Thailand is the unity of Thai flesh and blood'.
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There is evidence of continuous human habitation in present-day Thailand from 20,000 years ago to the present day.[26]:4 The earliest evidence of rice growing is dated at 2,000 BCE.[25]:4 Bronze appeared circa 1,250–1,000 BCE.[25]:4 The site of Ban Chiang in northeast Thailand currently ranks as the earliest known centre of copper and bronze production in Southeast Asia.[27] Iron appeared around 500 BCE.[25]:5 The Kingdom of Funan was the first and most powerful Southeast Asian kingdom at the time (2nd century BCE).[26]:5 The Mon people established the principalities of Dvaravati and Kingdom of Hariphunchai in the 6th century. The Khmer people established the Khmer empire, centred in Angkor, in the 9th century.[26]:7 Tambralinga, a Malay state controlling trade through the Malacca Strait, rose in the 10th century.[26]:5 The Indochina peninsula was heavily influenced by the culture and religions of India from the time of the Kingdom of Funan to that of the Khmer Empire.[28]
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The Thai people are of the Tai ethnic group, characterised by common linguistic roots.[29]:2 Chinese chronicles first mention the Tai peoples in the 6th century BCE. While there are many assumptions regarding the origin of Tai peoples, David K. Wyatt, a historian of Thailand, argued that their ancestors which at the present inhabit Laos, Thailand, Myanmar, India, and China came from the Điện Biên Phủ area between the 5th and the 8th century.[29]:6 Thai people began migrating into present-day Thailand around the 11th century, which Mon and Khmer people occupied at the time.[30] Thus Thai culture was influenced by Indian, Mon, and Khmer cultures.[31]
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According to French historian George Cœdès, "The Thai first enter history of Farther India in the eleventh century with the mention of Syam slaves or prisoners of war in Champa epigraphy, and "in the twelfth century, the bas-reliefs of Angkor Wat" where "a group of warriors" are described as Syam.[21]:190–191, 194–195
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After the decline of the Khmer Empire and Kingdom of Pagan in the early-13th century, various states thrived in their place. The domains of Tai people existed from the northeast of present-day India to the north of present-day Laos and to the Malay peninsula.[29]:38–9 During the 13th century, Tai people had already settled in the core land of Dvaravati and Lavo Kingdom to Nakhon Si Thammarat in the south. There are, however, no records detailing the arrival of the Tais.[29]:50–1
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Around 1240, Pho Khun Bang Klang Hao, a local Tai ruler, rallied the people to rebel against the Khmer. He later crowned himself the first king of Sukhothai Kingdom in 1238.[29]:52–3 Mainstream Thai historians count Sukhothai as the first kingdom of Thai people. Sukhothai expanded furthest during the reign of Ram Khamhaeng (r. 1279–1298). However, it was mostly a network of local lords who swore fealty to Sukhothai, not directly controlled by it.[29]:55–6 He is believed have invented Thai script and Thai ceramics were an important export in his era. Sukhothai embraced Theravada Buddhism in the reign of Maha Thammaracha I (1347–1368).
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To the north, Mangrai, who descended from a local ruler lineage of Ngoenyang, founded the kingdom of Lan Na in 1292, centered in Chiang Mai. He unified the surrounding area and his dynasty would rule the kingdom continuously for the next two centuries. He also created a network of states through political alliances to the east and north of the Mekong.[19]:8 While in the port in Lower Chao Phraya Basin, a federation around Phetchaburi, Suphan Buri, Lopburi, and the Ayutthaya area was created in the 11th century.[19]:8
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According to the most widely accepted version of its origin, the Ayutthaya Kingdom rose from the earlier, nearby Lavo Kingdom and Suvarnabhumi with Uthong as its first king. Ayutthaya was a patchwork of self-governing principalities and tributary provinces owing allegiance to the King of Ayutthaya under the mandala system.[32]:355 Its initial expansion was through conquest and political marriage. Before the end of the 15th century, Ayutthaya invaded the Khmer Empire three times and sacked its capital Angkor.[33]:26 Ayutthaya then became a regional power in place of the Khmer. Constant interference of Sukhothai effectively made it a vassal state of Ayutthaya and it was finally incorporated into the kingdom. Borommatrailokkanat brought about bureaucratic reforms which lasted into the 20th century and created a system of social hierarchy called sakdina, where male commoners were conscripted as corvée labourers for six months a year.[34]:107 Ayutthaya was interested in the Malay peninsula, but failed to conquer the Malacca Sultanate which was supported by the Chinese Ming Dynasty.[26]:11, 13
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European contact and trade started in the early-16th century, with the envoy of Portuguese duke Afonso de Albuquerque in 1511, Portugal became an allied and ceded some soldiers to King Rama Thibodi II.[35] The Portuguese were followed in the 17th century by the French, Dutch, and English. Rivalry for supremacy over Chiang Mai and the Mon people pitted Ayutthaya against the Burmese Kingdom. Several wars with its ruling dynasty Taungoo Dynasty starting in the 1540s in the reign of Tabinshwehti and Bayinnaung were ultimately ended with the capture of the capital in 1570.[34]:146–7 Then was a brief period of vassalage to Burma until Naresuan proclaimed independence in 1584.[19]:11
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Ayutthaya then sought to improve relations with European powers for many successive reigns. The kingdom especially prospered during cosmopolitan Narai's reign (1656–1688) when some European travelers regarded Ayutthaya as an Asian great power, alongside China and India.[25]:ix However, growing French influence later in his reign was met with nationalist sentiment and led eventually to the Siamese revolution of 1688.[34]:185–6 However, overall relations remained stable, with French missionaries still active in preaching Christianity.[34]:186
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After a bloody period of dynastic struggle, Ayutthaya entered into what has been called the golden age, a relatively peaceful episode in the second quarter of the 18th century when art, literature, and learning flourished. There were seldom foreign wars, apart from conflict with the Nguyễn Lords for control of Cambodia starting around 1715. The last fifty years of the kingdom witnessed bloody succession crises, where there were purges of court officials and able generals for many consecutive reigns. In 1765, a combined 40,000-strong force of Burmese armies invaded it from the north and west.[36]:250 The Burmese were under the new Alaungpaya dynasty quickly rose to be a new local power by 1759. After a 14-month siege, the capital city's wall fell and the city was burned in April 1767.[37]:218
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The capital and much territories lied in chaos after the war. The former capital was occupied by the Burmese garrison army and five local leaders declared themselves overlords, including the lords of Sakwangburi, Pimai, Chanthaburi, and Nakhon Si Thammarat. Chao Tak, a capable military leader, proceeded to make himself a lord by right of conquest, beginning with the legendary sack of Chanthaburi. Based at Chanthaburi, Chao Tak raised troops and resources, and sent a fleet up the Chao Phraya to take the fort of Thonburi. In the same year, Chao Tak was able to retake Ayutthaya from the Burmese only seven months after the fall of the city.[38]
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Chao Tak then crowned himself as Taksin and proclaimed Thonburi as temporary capital in the same year. He also quickly subdued the other warlords. His forces engaged in wars with Burma, Laos, and Cambodia, which successfully drove the Burmese out of Lan Na in 1775,[34]:225 captured Vientiane in 1778[34]:227–8 and tried to install a pro-Thai king in Cambodia in the 1770s. In his final years there was a coup, caused supposedly by his "insanity", and eventually Taksin and his sons were executed by his longtime companion General Chao Phraya Chakri (the future Rama I). He was the first king of the ruling Chakri Dynasty and founder of the Rattanakosin Kingdom on 6 April 1782.
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Under Rama I (1782–1809), Rattanakosin successfully defended against Burmese attacks and put an end to Burmese incursions. He also created suzerainty over large portions of Laos and Cambodia.[39] In 1821, Briton John Crawfurd was sent to negotiate a new trade agreement with Siam – the first sign of an issue which was to dominate 19th century Siamese politics.[40] Bangkok signed the Burney Treaty in 1826, after the British victory in the First Anglo-Burmese War.[34]:281 Anouvong of Vientiane, who misunderstood that Britain was about to attack Bangkok, started the Lao rebellion in 1826 and was defeated.[34]:283–5 Vientiane was destroyed and a large number of Lao people was relocated to Khorat Plateau as a result.[34]:285–6 Bangkok also waged several wars with Vietnam, where Bangkok successfully regained hegemony over Cambodia.[34]:290–2
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From the late-19th century, Siam tried to rule the ethnic groups in the realm as colonies.[34]:308 In the reign of Mongkut (1851–1868), who recognised the threat of Western powers, his court contacted the British government directly to defuse tensions.[34]:311 A British mission led by Sir John Bowring, Governor of Hong Kong, led to the signing of the Bowring Treaty, the first of many unequal treaties with Western countries. This, however, brought trade and economic development in Bangkok.[41] The unexpected death of Mongkut from malaria led to the reign of underage Prince Chulalongkorn, with Somdet Chaophraya Sri Suriwongse (Chuang Bunnag) acting as regent.[34]:327
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Chulalongkorn (r. 1868–1910) initiated centralisation, set up a privy council, and abolished slavery and the corvée system.[34] The Front Palace crisis of 1874 stalled attempts at further reforms.[34]:331–3 In the 1870s and 1880s, he incorporated the protectorates up north into the kingdom proper, which later expanded to the protectorates in the northeast and the south.[34]:334–5 He established twelve krom in 1888, which were equivalent to present-day ministries.[34]:347 The crisis of 1893 erupted, caused by French demands for Lao territory east of Mekong.[34]:350–3 Thailand is the only Southeast Asian nation not to have been colonised by a Western power,[42] in part because Britain and France agreed in 1896 to make the Chao Phraya valley a buffer state.[43] Not until the 20th century could Siam renegotiate every unequal treaty dating from the Bowring Treaty, including extraterritoriality, but at a price of many territorial exchanges. The advent of the monthon system marked the creation of the modern Thai nation-state.[34]:362–3 In 1905, there were unsuccessful rebellions in the ancient Patani area, Ubon Ratchathani, and Phrae in opposition to an attempt to blunt the power of local lords.[34]:371–3
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The Palace Revolt of 1912 was a failed attempt by Western-educated military officers to overthrow the absolute monarchy.[34]:397 Vajiravudh (r. 1910–1925) responded by propaganda for the entirety of his reign.[34]:402 He promoted the idea of the Thai nation.[34]:404 In 1917, Siam joined World War I on the side of the Allies as there were concerns that the Allies might punish neutral countries and refuse to amend past unequal treaties.[34]:407 In the aftermath Siam joined the Paris Peace Conference, and gained freedom of taxation and the revocation of extraterritoriality.[34]:408
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A bloodless revolution took place in 1932, carried out by a group of military and civilian officials Khana Ratsadon. Prajadhipok was forced to grant the country's first constitution, thereby ending centuries of absolute monarchy. The combined results of economic hardships brought on by the Great Depression, sharply falling rice prices, and a significant reduction in public spending caused discontent among aristocrats.[26]:25 In 1933, A counter-revolutionary rebellion occurred which aimed to reinstate absolute monarchy, but failed.[34]:446–8 Prajadhipok's conflict with the government eventually led to abdication. The government selected Ananda Mahidol, who was studying in Switzerland, to be the new king.[34]:448–9
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Later that decade, the army wing of Khana Ratsadon came to dominate Siamese politics. Plaek Phibunsongkhram who became premier in 1938, started political oppression and took an openly anti-royalist stance.[34]:457 His government adopted nationalism and Westernisation, anti-Chinese and anti-French policies.[26]:28 In 1940, there was a decree changing the name of the country from "Siam" to "Thailand". In 1941, Thailand was in a brief conflict with Vichy France resulting in Thailand gaining some Lao and Cambodian territories.[34]:462 On 8 December 1941, the Empire of Japan launched an invasion of Thailand, and fighting broke out shortly before Phibun ordered an armistice. Japan was granted free passage, and on 21 December Thailand and Japan signed a military alliance with a secret protocol, wherein Tokyo agreed to help Thailand regain territories lost to the British and French.[44] The Thai government declared war on the United States and the United Kingdom.[34]:465 The Free Thai Movement was launched both in Thailand and abroad to oppose the government and Japanese occupation.[34]:465–6 After the war ended in 1945, Thailand signed formal agreements to end the state of war with the Allies. Most Allied powers had not recognised Thailand's declaration of war.
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In June 1946, young King Ananda was found dead under mysterious circumstances. His younger brother Bhumibol Adulyadej ascended to the throne. Thailand joined the Southeast Asia Treaty Organization (SEATO) to become an active ally of the United States in 1954.[34]:493 Field Marshal Sarit Thanarat launched a coup in 1957, which removed Khana Ratsadon from politics. His rule (premiership 1959–1963) was autocratic; he built his legitimacy around the god-like status of the monarch and by channelling the government's loyalty to the king.[34]:511 His government improved the country's infrastructure and education.[34]:514 After the US joined the Vietnam War in 1961, there was a secret agreement wherein the US promised to protect Thailand.[34]:523
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The period brought about increasing modernisation and Westernisation of Thai society. Rapid urbanisation occurred when the rural populace sought work in growing cities. Rural farmers gained class consciousness and were sympathetic to the Communist Party of Thailand.[34]:528 Economic development and education enabled the rise of a middle class in Bangkok and other cities.[34]:534 In October 1971, there was a large demonstration against the dictatorship of Thanom Kittikachorn (premiership 1963–1973), which led to civilian casualties.[34]:541–3 Bhumibol installed Sanya Dharmasakti (premiership 1973–1975) to replace him, making it the first time that the king intervened in Thai politics directly since 1932.[45] The aftermath of the event marked a short-lived parliamentary democracy,[45] often called the "era when democracy blossomed." (ยุคประชาธิปไตยเบ่งบาน)
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Constant unrest and instability, as well as fear of a communist takeover after the fall of Saigon, made some ultra-right groups brand leftist students as communists.[34]:548 This culminated in the Thammasat University massacre in October 1976.[34]:548–9 A coup d'état on that day brought Thailand a new ultra-right government, which cracked down on media outlets, officials, and intellectuals, and fuelled the communist insurgency. Another coup the following year installed a more moderate government, which offered amnesty to communist fighters in 1978.
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Fueled by Indochina refugee crisis, Vietnamese border raids and economic hardships, Prem Tinsulanonda launched a successful coup and became the Prime Minister from 1980 to 1988. The communists abandoned the insurgency by 1983. Prem's premiership was dubbed "semi-democracy" because the Parliament was composed of all elected House and all appointed Senate. The 1980s also saw increasing intervention in politics by the monarch, who rendered two coup attempts against Prem failed. Thailand had its first elected prime minister in 1988.[46]
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Suchinda Kraprayoon, who was the coup leader in 1991 and said he would not seek to become prime minister, was nominated as one by the majority coalition government after the 1992 general election. This caused a popular demonstration in Bangkok, which ended with a military crackdown. Bhumibol intervened in the event and Suchinda then resigned.
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The 1997 Asian financial crisis originated in Thailand and ended the country's 40 years of uninterrupted economic growth.[47]:3 Chuan Leekpai's government took an IMF loan with unpopular provisions.[34]:576 The populist Thai Rak Thai party, led by prime minister Thaksin Shinawatra, governed from 2001 until 2006. His policies were successful in reducing rural poverty[48] and initiated universal healthcare in the country.[49] A South Thailand insurgency escalated starting from 2004. The 2004 Indian Ocean earthquake and tsunami hit the country, mostly in the south. Massive protests against Thaksin led by the People's Alliance for Democracy (PAD) started in his second term as prime minister and his tenure ended with a coup d'état in 2006. The junta installed a military government which lasted a year.
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In 2007, a civilian government led by the Thaksin-allied People's Power Party (PPP) was elected. Another protest led by PAD ended with the dissolution of PPP, and the Democrat Party led a coalition government in its place. The pro-Thaksin United Front for Democracy Against Dictatorship (UDD) protested both in 2009 and in 2010.
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After the general election of 2011, the populist Pheu Thai Party won a majority and Yingluck Shinawatra, Thaksin's younger sister, became prime minister. The People's Democratic Reform Committee organised another anti-Shinawatra protest[c] after the ruling party proposed an amnesty bill which would benefit Thaksin.[50] Yingluck dissolved parliament and a general election was scheduled, but was invalidated by the Constitution Court. The crisis ended with another coup d'état in 2014, the second coup in a decade.[d] The National Council for Peace and Order, a military junta led by General Prayut Chan-o-cha, has led the country since. Civil and political rights were restricted, and the country saw a surge in lèse-majesté cases. Political opponents and dissenters were sent to "attitude adjustment" camps.[51] Bhumibol, the longest-reigning Thai king, died in 2016, and his son Vajiralongkorn ascended to the throne. The referendum and adoption of Thailand's current constitution happened under the junta's rule.[e] In 2019, the junta agreed to schedule a general election in March.[51] Prayut continued his premiership with the support of Palang Pracharath Party-coalition in the House and junta-appointed Senate, amid allegations of election fraud.[53]
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Prior to 1932, Thai kings were absolute monarchs. During Sukhothai Kingdom, the king was seen as a Dharmaraja or 'king who rules in accordance with Dharma'. The system of government was a network of tributaries ruled by local lords. Modern absolute monarchy and statehood was established by Chulalongkorn when he transformed the decentralized protectorate system into a unitary state. On 24 June 1932, Khana Ratsadon (People's Party) carried out a bloodless revolution which marked the beginning of constitutional monarchy.
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Thailand has had 20 constitutions and charters since 1932, including the latest and current 2017 Constitution. Throughout this time, the form of government has ranged from military dictatorship to electoral democracy.[54][55] Thailand has had the fourth-most coups in the world.[56] "Uniformed or ex-military men have led Thailand for 55 of the 83 years" between 1932 and 2009.[57] Most recently, the National Council for Peace and Order ruled the country between 2014 and 2019.
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The politics of Thailand is conducted within the framework of a constitutional monarchy, whereby a hereditary monarch is head of state. The current King of Thailand is Vajiralongkorn (or Rama X), who has reigned since October 2016. The powers of the king are limited by the constitution and he is primarily a symbolic figurehead. The monarch is head of the armed forces and is required to be Buddhist as well as the Defender of the Faith. He has the power to appoint his heirs, the power to grant pardons, and the royal assent. The king is aided in his duties by the Privy Council of Thailand. However, the monarch still occasionally intervenes in Thai politics, as all constitutions pave the way for customary royal rulings. The monarchy is widely revered and lèse majesté is a severe crime in Thailand.
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Government is separated into three branches:
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Military and bureaucratic aristocrats fully controlled political parties between 1946 and 1980s.[59]:16 Most parties in Thailand are short-lived.[60]:246 Between 1992 and 2006, Thailand had a two-party system.[60]:245 Since 2000, two political parties dominated Thai general elections: one was the Pheu Thai Party (which was a successor of People's Power Party and the Thai Rak Thai Party), and the other was the Democrat Party. The political parties which support Thaksin Shinawatra won the most representatives every general election since 2001. Later constitutions created a multi-party system where a single party cannot gain a majority in the house.
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The 2007 constitution was partially abrogated by the military dictatorship that came to power in May 2014.[61]
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Thailand's kings are protected by lèse-majesté laws which allow critics to be jailed for three to fifteen years.[62] After the 2014 Thai coup d'état, Thailand had the highest number of lèse-majesté prisoners in the nation's history.[63][64] In 2017, the military court in Thailand sentenced a man to 35 years in prison for violating the country's lèse-majesté law.[64] Thailand has been rated not free on the Freedom House Index since 2014.[65] Thai activist and magazine editor Somyot Prueksakasemsuk, who was sentenced to eleven years' imprisonment for lèse-majesté in 2013,[66] is a designated prisoner of conscience by Amnesty International.[67]
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Totalling 513,120 square kilometres (198,120 sq mi), Thailand is the 50th-largest country by total area. It is slightly smaller than Yemen and slightly larger than Spain.[1]
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Thailand comprises several distinct geographic regions, partly corresponding to the provincial groups. The north of the country is the mountainous area of the Thai highlands, with the highest point being Doi Inthanon in the Thanon Thong Chai Range at 2,565 metres (8,415 ft) above sea level. The northeast, Isan, consists of the Khorat Plateau, bordered to the east by the Mekong River. The centre of the country is dominated by the predominantly flat Chao Phraya river valley, which runs into the Gulf of Thailand.
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Southern Thailand consists of the narrow Kra Isthmus that widens into the Malay Peninsula. Politically, there are six geographical regions which differ from the others in population, basic resources, natural features, and level of social and economic development. The diversity of the regions is the most pronounced attribute of Thailand's physical setting.
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The Chao Phraya and the Mekong River are the indispensable water courses of rural Thailand. Industrial scale production of crops use both rivers and their tributaries. The Gulf of Thailand covers 320,000 square kilometres (124,000 sq mi) and is fed by the Chao Phraya, Mae Klong, Bang Pakong, and Tapi Rivers. It contributes to the tourism sector owing to its clear shallow waters along the coasts in the southern region and the Kra Isthmus. The eastern shore of the Gulf of Thailand is an industrial centre of Thailand with the kingdom's premier deepwater port in Sattahip and its busiest commercial port, Laem Chabang.
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The Andaman Sea is a precious natural resource as it hosts popular and luxurious resorts. Phuket, Krabi, Ranong, Phang Nga and Trang, and their islands, all lay along the coasts of the Andaman Sea and, despite the 2004 tsunami, they remain a tourist magnet.
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Thailand's climate is influenced by monsoon winds that have a seasonal character (the southwest and northeast monsoon).[68]:2 Most of the country is classified as Köppen's tropical savanna climate.[69] The majority of the south as well as the eastern tip of the east have a tropical monsoon climate. Parts of the south also have a tropical rainforest climate.
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Thailand is divided into three seasons.[68]:2 The first is the rainy or southwest monsoon season (mid–May to mid–October), which is caused by southwestern wind from Indian Ocean.[68]:2 Rainfall is also contributed by Intertropical Convergence Zone (ITCZ) and tropical cyclones.[68]:2 August and September being the wettest period of the year.[68]:2 The country receives a mean annual rainfall of 1,200 to 1,600 mm (47 to 63 in).[68]:4 Winter or the northeast monsoon starts from mid–October until mid–February.[68]:2 Most of Thailand experiences dry weather with mild temperatures.[68]:2,4 The exception is southern Thailand where it receives abundant rainfall, particularly during October to November.[68]:2 Summer or the pre–monsoon season runs from mid–February until mid–May.[68]:3 Due to its inland nature and latitude, the north, northeast, central and eastern parts of Thailand experience a long period of warm weather, where emperatures can reach up to 40 °C (104 °F) during March to May,[68]:3 in contrast to close to or below 0 °C (32 °F) in some areas in winter.[68]:3 Southern Thailand is characterised by mild weather year-round with less diurnal and seasonal variations in temperatures due to maritime influences.[68]:3
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Thailand is among the world's ten countries that are most exposed to climate change; in particular, it is highly vulnerable to rising sea levels and extreme weather events.[70][71]
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Thailand has a mediocre but improving performance in the global Environmental Performance Index (EPI) with an overall ranking of 91 out of 180 countries in 2016. The environmental areas where Thailand performs worst (i.e., highest ranking) are air quality (167), environmental effects of the agricultural industry (106), and the climate and energy sector (93), the later mainly because of a high CO2 emission per KWh produced. Thailand performs best (i.e., lowest ranking) in water resource management (66), with some major improvements expected for the future, and sanitation (68).[73][74]
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The population of elephants, the country's national symbol, has fallen from fallen from 100,000 in 1850 to an estimated 2,000.[72] Poachers have long hunted elephants for ivory and hides, and now increasingly for meat.[75] Young elephants are often captured for use in tourist attractions or as work animals, which there were claims of mistreatment.[76] Although their use has declined since the government banned logging in 1989.
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Poaching of protected species remains a major problem. Tigers, leopards, and other large cats are hunted for their pelts. Many are farmed or hunted for their meat, supposedly has medicinal properties. Although such trade is illegal, the well-known Bangkok market Chatuchak is still known for the sale of endangered species.[77] The practice of keeping wild animals as pets affects species such as Asiatic black bear, Malayan sun bear, white-handed lar, pileated gibbon, and binturong.[78]
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Thailand is a unitary state; the administrative services of the executive branch are divided into three levels by National Government Organisation Act, BE 2534 (1991): central, provincial and local. Thailand is composed of 76 provinces (จังหวัด, changwat), which are first-level administrative divisions.[79] There are also two specially governed districts: the capital Bangkok and Pattaya. Bangkok is at provincial level and thus often counted as a province. Each province is divided into districts (อำเภอ, amphoe) and the districts are further divided into sub-districts (ตำบล, tambons). The name of each province's capital city (เมือง, mueang) is the same as that of the province. For example, the capital of Chiang Mai Province (Changwat Chiang Mai) is Mueang Chiang Mai or Chiang Mai. All provincial governors and district chiefs, which are administrators of provinces and districts respectively, are appointed by the central government.[80] Thailand's provinces are sometimes grouped into four to six regions, depending on the source.
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The foreign relations of Thailand are handled by the Minister of Foreign Affairs.
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Thailand participates fully in international and regional organisations. It is a major non-NATO ally and Priority Watch List Special 301 Report of the United States. The country remains an active member of ASEAN Association of Southeast Asian Nations. Thailand has developed increasingly close ties with other ASEAN members: Indonesia, Malaysia, the Philippines, Singapore, Brunei, Laos, Cambodia, Myanmar, and Vietnam, whose foreign and economic ministers hold annual meetings. Regional co-operation is progressing in economic, trade, banking, political, and cultural matters. In 2003, Thailand served as APEC (Asia Pacific Economic Cooperation) host. Dr. Supachai Panitchpakdi, the former Deputy Prime Minister of Thailand, currently serves as Secretary-General of the United Nations Conference on Trade and Development (UNCTAD). In 2005 Thailand attended the inaugural East Asia Summit.
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In recent years, Thailand has taken an increasingly active role on the international stage. When East Timor gained independence from Indonesia, Thailand, for the first time in its history, contributed troops to the international peacekeeping effort. Its troops remain there today as part of a UN peacekeeping force. As part of its effort to increase international ties, Thailand has reached out to such regional organisations as the Organization of American States (OAS) and the Organisation for Security and Cooperation in Europe (OSCE). Thailand has contributed troops to reconstruction efforts in Afghanistan and Iraq.
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Thaksin initiated negotiations for several free trade agreements with China, Australia, Bahrain, India, and the US. The latter especially was criticised, with claims that uncompetitive Thai industries could be wiped out.[81]
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Thaksin also announced that Thailand would forsake foreign aid, and work with donor countries to assist in the development of neighbours in the Greater Mekong Sub-region.[82] Thaksin sought to position Thailand as a regional leader, initiating various development projects in poorer neighbouring countries like Laos. More controversially, he established close, friendly ties with the Burmese dictatorship.[83]
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Thailand joined the US-led invasion of Iraq, sending a 423-strong humanitarian contingent.[84] It withdrew its troops on 10 September 2004. Two Thai soldiers died in Iraq in an insurgent attack.
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Abhisit appointed Peoples Alliance for Democracy leader Kasit Piromya as foreign minister. In April 2009, fighting broke out between Thai and Cambodian troops on territory immediately adjacent to the 900-year-old ruins of Cambodia's Preah Vihear Hindu temple near the border. The Cambodian government claimed its army had killed at least four Thais and captured 10 more, although the Thai government denied that any Thai soldiers were killed or injured. Two Cambodian and three Thai soldiers were killed. Both armies blamed the other for firing first and denied entering the other's territory.[85][86]
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The Royal Thai Armed Forces (กองทัพไทย; RTGS: Kong Thap Thai) constitute the military of the Kingdom of Thailand. It consists of the Royal Thai Army (กองทัพบกไทย), the Royal Thai Navy (กองทัพเรือไทย), and the Royal Thai Air Force (กองทัพอากาศไทย). It also incorporates various paramilitary forces.
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The Thai Armed Forces have a combined manpower of 306,000 active duty personnel and another 245,000 active reserve personnel.[87] The head of the Thai Armed Forces (จอมทัพไทย, Chom Thap Thai) is the king,[88] although this position is only nominal. The armed forces are managed by the Ministry of Defence of Thailand, which is headed by the Minister of Defence (a member of the cabinet of Thailand) and commanded by the Royal Thai Armed Forces Headquarters, which in turn is headed by the Chief of Defence Forces of Thailand.[89] Thai annual defense budget almost tripled from 78 billion baht in 2005 to 207 billion baht in 2016, accounting for approximately 1.5% of 2019 Thai GDP.[90][91] Thailand ranked 16th worldwide in the Military Strength Index based on the Credit Suisse report in September 2015.
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The military is also tasked with humanitarian missions, such as escorting Rohingya to Malaysia or Indonesia,[92] ensuring security and welfare for refugees during Indochina refugee crisis.[93]
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According to the constitution, serving in the armed forces is a duty of all Thai citizens.[94] Thailand still use active draft system for males over the age of 21. They are subjected to varying lengths of active service depending on the duration of reserve training as Territorial Defence Student and their level of education. Those who have completed three years or more of reserve training will be exempted entirely. The practice has long been criticized, as some media question its efficacy and value.[95][96] It is alleged that conscripts end up as servants to senior officers[97] or clerks in military cooperative shops.[98][99] In a report issued in March 2020, Amnesty International charged that Thai military conscripts face institutionalised abuse systematically hushed up by military authorities.[100]
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Critics observed that Thai military's main objective is to deal with internal rather than external threats.[101] Internal Security Operations Command is called the political arm of the Thai military, which has overlapping social and political functions with civilian bureaucracy. It also has anti-democracy mission.[101] The military is also notorious for numerous corruption incidents, such as accusation of human trafficking,[102] and nepotism in promotion of high-ranking officers.[103] The military is deeply entrenched in politics. Most recently, the appointed senators include more than 100 active and retired military.[104]
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In 2018 the literacy rate was 93.8%. The youth literacy rate was 98.1% in 2015.[105] Education is provided by a well-organised school system of kindergartens, primary, lower secondary and upper secondary schools, numerous vocational colleges, and universities. The private sector of education is well developed and significantly contributes to the overall provision of education which the government would not be able to meet with public establishments. Education is compulsory up to and including age 14, with the government providing free education through to age 17. Thailand is the 3rd most popular study destination in Asean. The number of international degree students in Thailand increased by fully 979% between 1999 and 2012, from 1,882 to 20,309 students. The most of international students come from Asian neighbor countries[106] from China, Myanmar, Cambodia and Vietnam.[107] The number of higher education institutions in Thailand has grown strongly over the past decades from just a handful of universities in the 1970s to 156 officially. The two top-ranking universities in Thailand are Chulalongkorn University and Mahidol University.[108] Thai universities research output still relatively low by international ranking comparison, Recent initiatives, such as the National Research University from 9 universities around the country[109] and Graduate research intensive university: VISTEC, designed to strengthen Thailand's national research universities, however, appear to be gaining traction. Thailand's research output, as measured by journal publications, increased by 20% between 2011 and 2016.[110]
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Teaching relies heavily on rote learning rather than on student-centred methodology. The establishment of reliable and coherent curricula for its primary and secondary schools is subject to such rapid changes that schools and their teachers are not always sure what they are supposed to be teaching, and authors and publishers of textbooks are unable to write and print new editions quickly enough to keep up with the volatility. Issues concerning university entrance has been in constant upheaval for a number of years. Nevertheless, Thai education has seen its greatest progress in the years since 2001. Most of the present generation of students are computer literate. Thailand was ranked 74th out of 100 countries globally for English proficiency.[111] Thailand has the second highest number of English-medium private international schools in Southeast Asian Nations, according to the International School Consultancy Group 181 schools around the country in 2017 compared to just 10 international schools for expatriate children in 1992.[106]
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Students in ethnic minority areas score consistently lower in standardised national and international tests.[112][113][114]
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This is likely due to unequal allocation of educational resources, weak teacher training, poverty, and low Thai language skill, the language of the tests.[112]
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[115]
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[116]
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Extensive nationwide IQ tests were administered to 72,780 Thai students from December 2010 to January 2011. The average IQ was found to be 98.59, which is higher than previous studies have found. IQ levels were found to be inconsistent throughout the country, with the lowest average of 88.07 found in the southern region of Narathiwat Province and the highest average of 108.91 reported in Nonthaburi Province. The Ministry of Public Health blames the discrepancies on iodine deficiency, and as of 2011[update] steps were being taken to require that iodine be added to table salt, a practice common in many Western countries.[117]
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In 2013, the Ministry of Information and Communication Technology announced that 27,231 schools would receive classroom-level access to high-speed internet.[118]
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In modern times, Thai scientists have made many significant contributions in various fields of study. For example, In chemistry, Krisana Kraisintu as known as the "Gypsy pharmacist".[119] She developed one of the first generic ARV fixed-dose combinations and dedicated her life to making medicines more affordable and accessible. Her efforts have saved countless lives in Africa,GPO-VIR has now been chosen by World Health Organization as the first regimen treatment for HIV/AIDS patients in poor countries.[120] In Thailand, this drug (GPO-VIR) is used in the national HIV/AIDS treatment programme, making it free of charge for 100,000 patients.[121] while Pongrama Ramasoota, He discoveries production of therapeutic human monoclonal antibodies against dengue virus and the world's first Dengue fever medication, include DNA vaccine development for dengue and Canine parvovirus.[122]
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Thailand has also made significant advances technology in the development of Medical Robotics. Medical robots have been used and promoted in Thailand in many areas, including surgery, diagnosis, rehabilitation and services.[123] and their use has been increasing. such as, an elderly care robot made by Thai manufacturer that Japanese nursing homes are widely using.[124] In surgery, back in 2019, The Medical Services Department has unveiled Thailand's robot created to help surgeons in brain surgery on patients afflicted with epilepsy.[125] back in 2017, Ramathibodi Hospital, a leading government hospital in Bangkok and a reputable medical school, successfully performed the first robot-assisted brain surgery in Asia.[126] For rehabilitation and therapy robots, were developed to help patients with arm and leg injuries perform practiced movements aided by the robots is the first prize winner of the i-MEDBOT Innovation Contest 2018 held by Thailand Center of Excellence for Life Sciences (TCELS).[127]
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According to the UNESCO Institute for Statistics, Thailand devoted 1% of its GDP to science research and development in 2017.[128] Between 2014 and 2016, Research and development workforce in Thailand increased from 84,216 people to 112,386 people.[129] The Thai government is developing new
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growth hubs by starting with the
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Eastern Economic Corridor of Innovation (EECi) to accelerating human resource and research development.[130] The National Science and Technology Development Agency is an agency of the government of Thailand which supports research in science and technology and its application in the Thai economy.[131]
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Thailand is an emerging economy and is considered a newly industrialised country. Thailand had a 2017 GDP of US$1.236 trillion (on a purchasing power parity basis).[137] Thailand is the 2nd largest economy in Southeast Asia after Indonesia. Thailand ranks midway in the wealth spread in Southeast Asia as it is the 4th richest nation according to GDP per capita, after Singapore, Brunei, and Malaysia.
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Thailand functions as an anchor economy for the neighbouring developing economies of Laos, Myanmar, and Cambodia. In the third quarter of 2014, the unemployment rate in Thailand stood at 0.84% according to Thailand's National Economic and Social Development Board (NESDB).[138]
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Thailand experienced the world's highest economic growth rate from 1985 to 1996 – averaging 12.4% annually. In 1997 increased pressure on the baht, a year in which the economy contracted by 1.9%, led to a crisis that uncovered financial sector weaknesses and forced the Chavalit Yongchaiyudh administration to float the currency. Prime Minister Chavalit Yongchaiyudh was forced to resign after his cabinet came under fire for its slow response to the economic crisis. The baht was pegged at 25 to the US dollar from 1978 to 1997. The baht reached its lowest point of 56 to the US dollar in January 1998 and the economy contracted by 10.8% that year, triggering the Asian financial crisis.
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Thailand's economy started to recover in 1999, expanding 4.2–4.4% in 2000, thanks largely to strong exports. Growth (2.2%) was dampened by the softening of the global economy in 2001, but picked up in the subsequent years owing to strong growth in Asia, a relatively weak baht encouraging exports, and increased domestic spending as a result of several mega projects and incentives of Prime Minister Thaksin Shinawatra, known as Thaksinomics. Growth in 2002, 2003, and 2004 was 5–7% annually.
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Growth in 2005, 2006, and 2007 hovered around 4–5%. Due both to the weakening of the US dollar and an increasingly strong Thai currency, by March 2008 the dollar was hovering around the 33 baht mark. While Thaksinomics has received criticism, official economic data reveals that between 2001 and 2011, Isan's GDP per capita more than doubled to US$1,475, while, over the same period, GDP in the Bangkok area increased from US$7,900 to nearly US$13,000.[139]
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With the instability surrounding major 2010 protests, the GDP growth of Thailand settled at around 4–5%, from highs of 5–7% under the previous civilian administration. Political uncertainty was identified as the primary cause of a decline in investor and consumer confidence. The IMF predicted that the Thai economy would rebound strongly from the low 0.1% GDP growth in 2011, to 5.5% in 2012 and then 7.5% in 2013, due to the monetary policy of the Bank of Thailand, as well as a package of fiscal stimulus measures introduced by the former Yingluck Shinawatra government.[140]
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Following the Thai military coup of 22 May 2014. In 2017, Concluded with information on the Thai economy's grew an inflation-adjusted 3.9%, up from 3.3% in 2016, marking its fastest expansion since 2012.[141]
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Thais have median wealth per one adult person of $1,469 in 2016,[142]:98 increasing from $605 in 2010.[142]:34 In 2016, Thailand was ranked 87th in Human Development Index, and 70th in the inequality-adjusted HDI.[143]
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In 2017, Thailand's median household income was ฿26,946 per month.[144]:1 Top quintile households had a 45.0% share of all income, while bottom quintile households had 7.1%.[144]:4 There were 26.9 million persons who had the bottom 40% of income earning less than ฿5,344 per person per month.[145]:5 During 2013–2014 Thai political crisis, a survey found that anti-government PDRC mostly (32%) had a monthly income of more than ฿50,000, while pro-government UDD mostly (27%) had between ฿10,000 and ฿20,000.[146]:7
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In 2014, Credit Suisse reported that Thailand was the world's third most unequal country, behind Russia and India.[147] Top 10% richest held 79% of the country's asset.[147] Top 1% richest held 58% worth of the economy.[147] Thai 50 richest families had a total net worth accounting to 30% of GDP.[147]
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In 2016, 5.81 million people lived in poverty, or 11.6 million people (17.2% of population) if "near poor" is included.[145]:1 Proportion of the poor relative to total population in each region was 12.96% in the Northeast, 12.35% in the South, and 9.83% in the North.[145]:2 In 2017, there were 14 million people who applied for social welfare (yearly income of less than ฿100,000 was required).[147] At the end of 2017, Thailand's total household debt was ฿11.76 trillion.[134]:5 In 2010, 3% of all household were bankrupt.[136]:5 In 2016, there were estimated 30,000 homeless persons in the country.[148]
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The economy of Thailand is heavily export-dependent, with exports accounting for more than two-thirds of gross domestic product (GDP). Thailand exports over US$105 billion worth of goods and services annually.[1] Major exports include cars, computers, electrical appliances, rice, textiles and footwear, fishery products, rubber, and jewellery.[1]
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Substantial industries include electric appliances, components, computer components, and vehicles. Thailand's recovery from the 1997–1998 Asian financial crisis depended mainly on exports, among various other factors. As of 2012[update], the Thai automotive industry was the largest in Southeast Asia and the 9th largest in the world.[149][150][151] The Thailand industry has an annual output of near 1.5 million vehicles, mostly commercial vehicles.[151]
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Most of the vehicles built in Thailand are developed and licensed by foreign producers, mainly Japanese and American. The Thai car industry takes advantage of the ASEAN Free Trade Area (AFTA) to find a market for many of its products. Eight manufacturers, five Japanese, two US, and Tata of India, produce pick-up trucks in Thailand.[152] As of 2012, Thailand was the second largest consumer of pick-up trucks in the world, after the US.[153] In 2014, pick-ups accounted for 42% of all new vehicle sales in Thailand.[152]
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Tourism makes up about 6% of the country's economy. Thailand was the most visited country in Southeast Asia in 2013, according to the World Tourism Organisation. Estimates of tourism receipts directly contributing to the Thai GDP of 12 trillion baht range from 9 percent (1 trillion baht) (2013) to 16 percent.[154] When including the indirect effects of tourism, it is said to account for 20.2 percent (2.4 trillion baht) of Thailand's GDP.[155]:1
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Asian tourists primarily visit Thailand for Bangkok and the historical, natural, and cultural sights in its vicinity. Western tourists not only visit Bangkok and surroundings, but in addition many travel to the southern beaches and islands. The north is the chief destination for trekking and adventure travel with its diverse ethnic minority groups and forested mountains. The region hosting the fewest tourists is Isan. To accommodate foreign visitors, a separate tourism police with offices were set up in the major tourist areas and an emergency telephone number.[156]
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Thailand ranks 5th biggest medical tourism destination of inbound medical tourism spending, according to World Travel and Tourism Council, attracting over 2.5 million visitors in 2018.[157] The country is also Asia's number one.[158] The country is popular for the growing practice of sex reassignment surgery (SRS) and cosmetic surgery. In 2010–2012, more than 90% of the visitors travelled to Thailand for SRS.[159]
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Prostitution in Thailand and sex tourism also form a de facto part of the economy. Campaigns promote Thailand as exotic to attract tourists.[160] One estimate published in 2003 placed the trade at US$4.3 billion per year or about 3% of the Thai economy.[161] It is believed that at least 10% of tourist dollars are spent on the sex trade.[162]
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Forty-nine per cent of Thailand's labour force is employed in agriculture.[163] This is down from 70% in 1980.[163] Rice is the most important crop in the country and Thailand had long been the world's leading exporter of rice, until recently falling behind both India and Vietnam.[164] Thailand has the highest percentage of arable land, 27.25%, of any nation in the Greater Mekong Subregion.[165] About 55% of the arable land area is used for rice production.[166]
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Agriculture has been experiencing a transition from labour-intensive and transitional methods to a more industrialised and competitive sector.[163] Between 1962 and 1983, the agricultural sector grew by 4.1% per year on average and continued to grow at 2.2% between 1983 and 2007.[163] The relative contribution of agriculture to GDP has declined while exports of goods and services have increased.
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Furthermore, access to biocapacity in Thailand is lower than world average. In 2016, Thailand had 1.2 global hectares[167] of biocapacity per person within its territory, a little less than world average of 1.6 global hectares per person.[168] In contrast, in 2016, they used 2.5 global hectares of biocapacity – their ecological footprint of consumption. This means they use about twice as much biocapacity as Thailand contains. As a result, Thailand is running a biocapacity deficit.[167]
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The State Railway of Thailand (SRT) operates all of Thailand's national rail lines. Bangkok Railway Station (Hua Lamphong Station) is the main terminus of all routes. Phahonyothin and ICD Lat Krabang are the main freight terminals. As of 2017[update] SRT had 4,507 km (2,801 mi) of track, all of it meter gauge except the Airport Link. Nearly all is single-track (4,097 km), although some important sections around Bangkok are double (303 km or 188 mi) or triple-tracked (107 km or 66 mi) and there are plans to extend this.[169] Rail transport in Bangkok includes long-distance services, and some daily commuter trains running from and to the outskirts of the city during the rush hour, but passenger numbers have remained low. There are also three rapid transit rail systems in the capital.
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Thailand has 390,000 km (242,335 miles) of highways.[170] According to the BBC Thailand has 462,133 roads and many multi-lane highways. As of 2017[update] Thailand has 37 million registered vehicles, 20 million of them motorbikes. A number of undivided two-lane highways have been converted into divided four-lane highways. A Bangkok – Chon Buri motorway (Route 7) now links to the new airport and Eastern Seaboard. There are 4,125 public vans operating on 114 routes from Bangkok alone.[171] Other forms of road transport includes tuk-tuks, taxis—as of November 2018, Thailand has 80,647 registered taxis nationwide[172]—vans (minibus), motorbike taxis and songthaews.
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As of 2012, Thailand had 103 airports with 63 paved runways, in addition to 6 heliports. The busiest airport in the county is Bangkok's Suvarnabhumi Airport.
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75% of Thailand's electrical generation is powered by natural gas in 2014.[173] Coal-fired power plants produce an additional 20% of electricity, with the remainder coming from biomass, hydro, and biogas.[173]
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Thailand produces roughly one-third of the oil it consumes. It is the second largest importer of oil in SE Asia. Thailand is a large producer of natural gas, with reserves of at least 10 trillion cubic feet. After Indonesia, it is the largest coal producer in SE Asia, but must import additional coal to meet domestic demand.
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Thailand has a diverse and robust informal labour sector—in 2012, it was estimated that informal workers comprised 62.6% of the Thai workforce. The Ministry of Labour defines informal workers to be individuals who work in informal economies and do not have employee status under a given country's Labour Protection Act (LPA). The informal sector in Thailand has grown significantly over the past 60 years over the course of Thailand's gradual transition from an agriculture-based economy to becoming more industrialised and service-oriented.[174] Between 1993 and 1995, ten percent of the Thai labour force moved from the agricultural sector to urban and industrial jobs, especially in the manufacturing sector. It is estimated that between 1988 and 1995, the number of factory workers in the country doubled from two to four million, as Thailand's GDP tripled.[175] While the Asian Financial Crisis that followed in 1997 hit the Thai economy hard, the industrial sector continued to expand under widespread deregulation, as Thailand was mandated to adopt a range of structural adjustment reforms upon receiving funding from the IMF and World Bank. These reforms implemented an agenda of increased privatisation and trade liberalisation in the country, and decreased federal subsidisation of public goods and utilities, agricultural price supports, and regulations on fair wages and labour conditions.[176] These changes put further pressure on the agricultural sector, and prompted continued migration from the rural countryside to the growing cities. Many migrant farmers found work in Thailand's growing manufacturing industry, and took jobs in sweatshops and factories with few labour regulations and often exploitative conditions.[177]
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Those that could not find formal factory work, including illegal migrants and the families of rural Thai migrants that followed their relatives to the urban centres, turned to the informal sector to provide the extra support needed for survival—under the widespread regulation imposed by the structural adjustment programs, one family member working in a factory or sweatshop made very little. Scholars argue that the economic consequences and social costs of Thailand's labour reforms in the wake of the 1997 Asian Financial Crisis fell on individuals and families rather than the state. This can be described as the "externalisation of market risk", meaning that as the country's labour market became increasingly deregulated, the burden and responsibility of providing an adequate livelihood shifted from employers and the state to the workers themselves, whose families had to find jobs in the informal sector to make up for the losses and subsidise the wages being made by their relatives in the formal sector. The weight of these economic changes hit migrants and the urban poor especially hard, and the informal sector expanded rapidly as a result.[176]
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Today, informal labour in Thailand is typically broken down into three main groups: subcontracted/self employed/home-based workers, service workers (including those that are employed in restaurants, as street vendors, masseuses, taxi drivers, and as domestic workers), and agricultural workers. Not included in these categories are those that work in entertainment, nightlife, and the sex industry. Individuals employed in these facets of the informal labour sector face additional vulnerabilities, including recruitment into circles of sexual exploitation and human trafficking.[174]
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In general, education levels are low in the informal sector. A 2012 study found that 64% of informal workers had not completed education beyond primary school. Many informal workers are also migrants, only some of which have legal status in the country. Education and citizenship are two main barriers to entry for those looking to work in formal industries, and enjoy the labour protections and social security benefits that come along with formal employment. Because the informal labour sector is not recognised under the Labour Protection Act (LPA), informal workers are much more vulnerable labour to exploitation and unsafe working conditions than those employed in more formal and federally recognised industries. While some Thai labour laws provide minimal protections to domestic and agricultural workers, they are often weak and difficult to enforce. Furthermore, Thai social security policies fail to protect against the risks many informal workers face, including workplace accidents and compensation as well as unemployment and retirement insurance. Many informal workers are not legally contracted for their employment, and many do not make a living wage.[174] As a result, labour trafficking is common in the region, affecting children and adults, men and women, and migrants and Thai citizens alike.
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Thailand had a population of 66,558,935 as of 2019.[4] Thailand's population is largely rural, concentrated in the rice-growing areas of the central, northeastern and northern regions. About 45.7% of Thailand's population lived in urban areas as of 2010[update], concentrated mostly in and around the Bangkok Metropolitan Area.
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Thailand's government-sponsored family planning program resulted in a dramatic decline in population growth from 3.1% in 1960 to around 0.4% today. In 1970, an average of 5.7 people lived in a Thai household. At the time of the 2010 census, the average Thai household size was 3.2 people.
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Thai nationals make up the majority of Thailand's population, 95.9% in 2010. The remaining 4.1% of the population are Burmese (2.0%), others 1.3%, and unspecified 0.9%.[1]
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According to the Royal Thai Government's 2011 Country Report to the UN Committee responsible for the International Convention for the Elimination of All Forms of Racial Discrimination, available from the Department of Rights and Liberties Promotion of the Thai Ministry of Justice,:3 62 ethnic communities are officially recognised in Thailand. Twenty million Central Thai (together with approximately 650,000 Khorat Thai) make up approximately 20,650,000 (34.1 percent) of the nation's population of 60,544,937[180] at the time of completion of the Mahidol University Ethnolinguistic Maps of Thailand data (1997).[181]
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The 2011 Thailand Country Report provides population numbers for mountain peoples ('hill tribes') and ethnic communities in the Northeast and is explicit about its main reliance on the Mahidol University Ethnolinguistic Maps of Thailand data.[181] Thus, though over 3.288 million people in the Northeast alone could not be categorised, the population and percentages of other ethnic communities circa 1997 are known for all of Thailand and constitute minimum populations. In descending order, the largest (equal to or greater than 400,000) are a) 15,080,000 Lao (24.9 percent) consisting of the Thai Lao (14 million) and other smaller Lao groups, namely the Thai Loei (400–500,000), Lao Lom (350,000), Lao Wiang/Klang (200,000), Lao Khrang (90,000), Lao Ngaew (30,000), and Lao Ti (10,000; b) six million Khon Muang (9.9 percent, also called Northern Thais); c) 4.5 million Pak Tai (7.5 percent, also called Southern Thais); d) 1.4 million Khmer Leu (2.3 percent, also called Northern Khmer); e) 900,000 Malay (1.5%); f) 500,000 Nyaw (0.8 percent); g) 470,000 Phu Thai (0.8 percent); h) 400,000 Kuy/Kuay (also known as Suay) (0.7 percent), and i) 350,000 Karen (0.6 percent).:7–13 Thai Chinese, those of significant Chinese heritage, are 14% of the population, while Thais with partial Chinese ancestry comprise up to 40% of the population.[182] Thai Malays represent 3% of the population, with the remainder consisting of Mons, Khmers and various "hill tribes". The country's official language is Thai and the primary religion is Theravada Buddhism, which is practised by around 95% of the population.
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Increasing numbers of migrants from neighbouring Myanmar, Laos, and Cambodia, as well as from Nepal and India, have pushed the total number of non-national residents to around 3.5 million as of 2009[update], up from an estimated 2 million in 2008, and about 1.3 million in 2000.[183] Some 41,000 Britons and 20,000 Australians live in Thailand.[184][185]
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The official language of Thailand is Thai, a Kra–Dai language closely related to Lao, Shan in Myanmar, and numerous smaller languages spoken in an arc from Hainan and Yunnan south to the Chinese border. It is the principal language of education and government and spoken throughout the country. The standard is based on the dialect of the central Thai people, and it is written in the Thai alphabet, an abugida script that evolved from the Khmer alphabet.
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Sixty-two languages were recognised by the Royal Thai Government in the 2011 Country Report to the UN Committee responsible for the International Convention for the Elimination of All Forms of Racial Discrimination, which employed an ethnolinguistic approach and is available from the Department of Rights and Liberties Promotion of the Thai Ministry of Justice.:3 Southern Thai is spoken in the southern provinces, and Northern Thai is spoken in the provinces that were formerly part of the independent kingdom of Lan Na. For the purposes of the national census, which does not recognise all 62 languages recognised by the Royal Thai Government in the 2011 Country Report, four dialects of Thai exist; these partly coincide with regional designations.
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The largest of Thailand's minority languages is the Lao dialect of Isan spoken in the northeastern provinces. Although sometimes considered a Thai dialect, it is a Lao dialect, and the region where it is traditionally spoken was historically part of the Lao kingdom of Lan Xang.[citation needed] In the far south, Kelantan-Pattani Malay is the primary language of Malay Muslims. Varieties of Chinese are also spoken by the large Thai Chinese population, with the Teochew dialect best-represented.
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Numerous tribal languages are also spoken, including many Austroasiatic languages such as Mon, Khmer, Viet, Mlabri and Orang Asli; Austronesian languages such as Cham and Moken; Sino-Tibetan languages like Lawa, Akha, and Karen; and other Tai languages such as Tai Yo, Phu Thai, and Saek. Hmong is a member of the Hmong–Mien languages, which is now regarded as a language family of its own.
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English is a mandatory school subject, but the number of fluent speakers remains low, especially outside cities.
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Thailand's prevalent religion is Theravada Buddhism, which is an integral part of Thai identity and culture. Active participation in Buddhism is among the highest in the world. According to the 2000 census, 94.6% and 93.58% in 2010 of the country's population self-identified as Buddhists of the Theravada tradition. Muslims constitute the second largest religious group in Thailand, comprising 4.29% of the population in 2015.[187]
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Islam is concentrated mostly in the country's southernmost provinces: Pattani, Yala, Satun, Narathiwat, and part of Songkhla Chumphon, which are predominantly Malay, most of whom are Sunni Muslims. Christians represented 1.17% (2015) of the population in 2015, with the remaining population consisting of Hindus and Sikhs, who live mostly in the country's cities. There is also a small but historically significant Jewish community in Thailand dating back to the 17th century.
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The constitution does not name official state religion, and provides for freedom of religion. Even the authority formally does not register new religious groups that have not been accepted and limit the number of missionaries, unregistered religious organisations as well as missionaries who are allowed to operate freely. There have been no widespread reports of societal abuses or discrimination based on religious belief or practice.[188]
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Thailand ranks world's 6th, and Asia's 1st in the 2019 Global Health Security Index of global health security capabilities in 195 countries,[189] making it the only developing country on the world's top ten. Thailand had 62 hospitals accredited by Joint Commission International.[190] In 2002, Bumrungrad became the first hospital in Asia to meet the standard.
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Health and medical care is overseen by the Ministry of Public Health (MOPH), along with several other non-ministerial government agencies, with total national expenditures on health amounting to 4.3 percent of GDP in 2009. Non-communicable diseases form the major burden of morbidity and mortality, while infectious diseases including malaria and tuberculosis, as well as traffic accidents, are also important public health issues. The current Minister for Public Health is Anutin Charnvirakul.
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In December 2018 the interim parliament voted to legalise the use of cannabis for medical reasons. Recreational use remained unlawful. The National Legislative Assembly had 166 votes in favour of the amendment to the Narcotics Bill, while there were no nay votes and 13 abstentions. The vote makes Thailand the first Southeast Asian country to allow the use of medical cannabis.[191]
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Thai culture and traditions incorporate a great deal of influence from India, China, Cambodia, and the rest of Southeast Asia. Thailand's national religion, Theravada Buddhism, is central to modern Thai identity. Thai Buddhism has evolved over time to include many regional beliefs originating from Hinduism, animism, as well as ancestor worship. The official calendar in Thailand is based on the Eastern version of the Buddhist Era (BE). Thai identity today is a social construct of Phibun regime in 1940s.
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Several ethnic groups mediated change between their traditional local culture, national Thai, and global cultural influences. Overseas Chinese also form a significant part of Thai society, particularly in and around Bangkok. Their successful integration into Thai society has allowed them to hold positions of economic and political power. Thai Chinese businesses prosper as part of the larger bamboo network.[192]
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Wai is traditional Thai greeting, and is generally offered first by person who is younger or lower in social status and position. Respects for elderly (by age, position, monks, or certain professions) is Thai mores. As with other Asian cultures, respect towards ancestors is an essential part of Thai spiritual practice. Thais have strong sense of social hierarchy, reflecting in many classes of honorifics. Seniority is paramount in Thai culture. Elders have by tradition ruled in family decisions or ceremonies. Older siblings have duties to younger ones.
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Thais have a strong sense of hospitality and generosity.
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Taboos in Thailand include touching someone's head or pointing with the feet, as the head is considered the most sacred and the foot the lowest part of the body.
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The origins of Thai art were very much influenced by Buddhist art and by scenes from the Indian epics. Traditional Thai sculpture almost exclusively depicts images of the Buddha, being very similar with the other styles from Southeast Asia. Traditional Thai paintings usually consist of book illustrations, and painted ornamentation of buildings such as palaces and temples. Thai art was influenced by indigenous civilisations of the Mon and other civilisations. By the Sukothai and Ayutthaya period, thai had developed into its own unique style and was later further influenced by the other Asian styles, mostly by Sri Lankan and Chinese. Thai sculpture and painting, and the royal courts provided patronage, erecting temples and other religious shrines as acts of merit or to commemorate important events.[193]
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Traditional Thai paintings showed subjects in two dimensions without perspective. The size of each element in the picture reflected its degree of importance. The primary technique of composition is that of apportioning areas: the main elements are isolated from each other by space transformers. This eliminated the intermediate ground, which would otherwise imply perspective. Perspective was introduced only as a result of Western influence in the mid-19th century. Monk artist Khrua In Khong is well known as the first artist to introduce linear perspective to Thai traditional art.[194]
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The most frequent narrative subjects for paintings were or are: the Jataka stories, episodes from the life of the Buddha, the Buddhist heavens and hells, themes derived from the Thai versions of the Ramayana and Mahabharata, not to mention scenes of daily life. Some of the scenes are influenced by Thai folklore instead of following strict Buddhist iconography.[193]
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Architecture is the preeminent medium of the country's cultural legacy and reflects both the challenges of living in Thailand's sometimes extreme climate as well as, historically, the importance of architecture to the Thai people's sense of community and religious beliefs. Influenced by the architectural traditions of many of Thailand's neighbours, it has also developed significant regional variation within its vernacular and religious buildings.
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The Ayutthaya Kingdom movement, which went from approximately 1350 to 1767, was one of the most fruitful and creative periods in Thai architecture The identity of architecture in Ayutthaya period is designed to display might and riches so it has great size and appearance. The temples in Ayutthaya seldom built eaves stretching from the masterhead. The dominant feature of this style is sunlight shining into buildings. During the latter part of the Ayutthaya period, architecture was regarded as a peak achievement that responded to the requirements of people and expressed the gracefulness of Thainess.[195]
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Buddhist temples in Thailand are known as "wats", from the Pāḷi vāṭa, meaning an enclosure. A temple has an enclosing wall that divides it from the secular world. Wat architecture has seen many changes in Thailand in the course of history. Although there are many differences in layout and style, they all adhere to the same principles.[196]
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Thai literature has had a long history. Even before the establishment of the Sukhothai Kingdom there existed oral and written works.
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During the Sukhothai, Most literary works were written in simple prose with certain alliteration schemes. Major works include King Ram Khamhaeng Inscription. King Ram Khamhaeng's Stone Inscription is considered the first Thai literary work in Thai script. It gives an account of the life of King Ramkhamhaeng the Great, the way of life of Thai people in general, laws, religion, economic and political stability. Trai Phum Phra Ruang, was written in 1345 by King Maha Thammaracha I, the fifth king of Sukhothai. It expounds Buddhist philosophy based on a profound and extensive study with reference to over 30 sacred texts. The work could be considered the nation's first piece of research dissertation. It was written in beautiful prose rich in allusions and imagery. It is a treatise on Buddhist cosmology, ethics, biology and belief system.[197]
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During the Ayutthaya, The period produced a variety of forms on diverse subjects. New poetic forms were created, with different rhyme schemes and metres. It is common to find a combination of different poetic forms in one poetic work. Lilit Yuan Phai is a narrative poem describing the war between King Borommatrailokkanat of Ayutthaya and Prince Tilokkarat of Lan Na. One of the most beautiful literary works is Kap He Ruea composed by Prince Thammathibet comparing the scenic beauty to that of his beloved lady on a boat journey in the nirat tradition. Traditionally, the verse is sung during the colourful royal barge procession. It has been the model for subsequent poets to emulate. The same prince also composed the greatly admired Kap Ho Khlong on the Visit to Than Thongdaeng and Kap Ho Khlong Nirat Phrabat.[198]
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Despite its short period of 15 years, Thon Buri Period produced Ramakian, a verse drama to which King Taksin the Great contributed his poetic talent. The revival of literature at this time is remarkable since the country had not quite recovered from the aftermath of war. Some poets who later became a major force in the early Rattanakosin Period had already begun writing at this time.
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During the 18th century Rattanakosin Period. After sporadic fighting at the beginning of the period, the country gradually returned to normal. It is only natural that many of the early Rattanakosin works should deal with war and military strategy. Some examples are Nirat Rop Phama Thi Tha Din Daeng, Phleng Yao Rop Phama Thi Nakhon Si Thammarat.In the performing arts, perhaps the most important dramatic achievement is the complete work of Ramakian by King Rama I. In addition, There were also verse recitals with musical accompaniment, such as Mahori telling the story of Kaki, Sepha relating the story of Khun Chang Khun Phaen. Other recitals include Sri Thanonchai. The most important Thai poet in this period was Sunthorn Phu (สุนทรภู่) (1786–1855), widely known as "the bard of Rattanakosin" (Thai: กวีเอกแห่งกรุงรัตนโกสินทร์). Sunthorn Phu is best known for his epic poem Phra Aphai Mani (Thai: พระอภัยมณี), which he started in 1822 (while in jail) and finished in 1844. Phra Aphai Mani is a versified fantasy-adventure novel, a genre of Siamese literature known as nithan kham klon (Thai: นิทานคำกลอน).[198]
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Aside from folk and regional dances (southern Thailand's Menora (dance) and Ramwong, for example), the two major forms of Thai classical dance drama are Khon and Lakhon nai. In the beginning, both were exclusively court entertainments and it was not until much later that a popular style of dance theatre, likay, evolved as a diversion for common folk who had no access to royal performances.[199]
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Folk dance forms include dance theater forms like likay, numerous regional dances (ram), the ritual dance ram muay, and homage to the teacher, wai khru. Both ram muay and wai khru take place before all traditional muay Thai matches. The wai is also an annual ceremony performed by Thai classical dance groups to honor their artistic ancestors.
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Thai classical music is synonymous with those stylized court ensembles and repertoires that emerged in their present form within the royal centers of Central Thailand some 800 years ago. These ensembles, while being influenced by older practices are today uniquely Thai expressions. While the three primary classical ensembles, the Piphat, Khrueang sai and Mahori differ in significant ways, they all share a basic instrumentation and theoretical approach. Each employs small ching hand cymbals and krap wooden sticks to mark the primary beat reference. Thai classical music has had a wide influence on the musical traditions of neighboring countries. The traditional music of Myanmar was strongly influenced by the Thai music repertoire, called Yodaya (ယိုးဒယား), which was brought over from the Ayutthaya Kingdom. As Siam expanded its political and cultural influence to Laos and Cambodia during the early Rattanakosin period, its music was quickly absorbed by the Cambodia and Lao courts.
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Thai films are exported and exhibited in Southeast Asia.[200] Thai cinema has developed its own unique identity and now being internationally recognized for their culture-driven.[201] Films such as Ong-Bak: Muay Thai Warrior (2003) and 'Tom-Yum-Goong (2005), starred Tony Jaa, feature distinctive aspects of Thai martial arts "Muay Thai".
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Thai horror has always had a significant cult following, unique take on tales from beyond the grave. More recently, horror films such as Shutter (2004), was one of the best-known Thai horror movies and recognized worldwide.[202] Other examples include The Unseeable (2006), Alone (2007), Body (2007), Coming Soon (2008), 4bia (2008), Phobia 2 (2009), Ladda Land (2011), Pee Mak (2013), and The Promise (2017).
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Thai heist thriller film Bad Genius (2017), was one of the most internationally successful Thai film, It broke Thai film earning records in several Asian countries,[203][204] Bad Genius won in 12 categories at the 27th Suphannahong National Film Awards, and also won the Jury Award at the 16th New York Asian Film Festival with a worldwide collection of more than $42 million.[205]
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Thailand television dramas, known as Lakorn, Lakorn have become popular in Thailand and its neighbors.[206] Many dramas tend to have a romantic focus, such as Khluen Chiwit, U-Prince, Ugly Duckling, The Crown Princess and teen dramas television series, such as 2gether: The Series, The Gifted, Girl From Nowhere, Hormones: The Series.
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The Entertainment industries (film and television) are estimated to have directly contributed $2.1 billion in gross domestic product (GDP) to the Thai economy in 2011. They also directly supported 86,600 jobs.[207] Amongst several Dance-pop artists who have made internationally successful can be mentioned "Lisa" Lalisa Manoban[208] and Tata Young.
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Thai cuisine is one of the most popular in the world.[209][210] Thai food blends five fundamental tastes: sweet, spicy, sour, bitter, and salty. The herbs and spices most used in Thai cooking themselves have medicinal qualities such as garlic, lemongrass, Kaffir lime, galangal, turmeric, coriander, coconut milk.[211] Each region of Thailand has its specialities: kaeng khiao wan (green curry) in the central region, som tam (green papaya salad) in the northeast, khao soi in the north, Massaman curry in the south.
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In 2017, seven Thai dishes appeared on a list of the "World's 50 Best Foods"— an online poll of worldwide by CNN Travel. Thailand had more dishes on the list than any other country. They were: tom yam goong (4th), pad Thai (5th), som tam (6th), Massaman curry (10th), green curry (19th), Thai fried rice (24th) and mu nam tok (36th).[212]
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The staple food in Thailand is rice, particularly jasmine rice (also known as hom Mali) which forms part of almost every meal. Thailand is a leading exporter of rice, and Thais consume over 100 kg of milled rice per person per year.[213]
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Thailand generally uses the metric system, but traditional units of measurement for land area are used, and imperial units of measurement are occasionally used for building materials, such as wood and plumbing fixtures. Years are numbered as B.E. (Buddhist Era) in educational settings, civil service, government, contracts, and newspaper datelines. However, in banking, and increasingly in industry and commerce, standard Western year (Christian or Common Era) counting is the standard practice.[214]
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Muay Thai (Thai: มวยไทย, RTGS: Muai Thai, [muaj tʰaj], lit. "Thai boxing") is a combat sport of Thailand that uses stand-up striking along with various clinching techniques. Muay Thai became widespread internationally in the late-20th to 21st century, when Westernized practitioners from Thailand began competing in kickboxing and mixed rules matches as well as matches under muay Thai rules around the world, Famous practitioners such as Buakaw Banchamek, Samart Payakaroon, Dieselnoi Chor Thanasukarn and Apidej Sit-Hirun. Buakaw Banchamek has probably brought more international interest in Muay Thai than any other Muay Thai fighters ever had.[215]
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Association football has overtaken muay Thai as the most widely followed sport in contemporary Thai society. Thailand national football team has played the AFC Asian Cup six times and reached the semifinals in 1972. The country has hosted the Asian Cup twice, in 1972 and in 2007. The 2007 edition was co-hosted together with Indonesia, Malaysia and Vietnam. It is not uncommon to see Thais cheering their favourite English Premier League teams on television and walking around in replica kit. Another widely enjoyed pastime, and once a competitive sport, is kite flying.
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Volleyball is rapidly growing as one of the most popular sports. The women's team has often participated in the World Championship, World Cup, and World Grand Prix Asian Championship. They have won the Asian Championship twice and Asian Cup once. By the success of the women's team, the men team has been growing as well.
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Takraw (Thai: ตะกร้อ) is a sport native to Thailand, in which the players hit a rattan ball and are only allowed to use their feet, knees, chest, and head to touch the ball. Sepak takraw is a form of this sport which is similar to volleyball. The players must volley a ball over a net and force it to hit the ground on the opponent's side. It is also a popular sport in other countries in Southeast Asia. A rather similar game but played only with the feet is buka ball.
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Snooker has enjoyed increasing popularity in Thailand in recent years, with interest in the game being stimulated by the success of Thai snooker player James Wattana in the 1990s.[216] Other notable players produced by the country include Ratchayothin Yotharuck, Noppon Saengkham and Dechawat Poomjaeng.[217]
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Rugby is also a growing sport in Thailand with the Thailand national rugby union team rising to be ranked 61st in the world.[218] Thailand became the first country in the world to host an international 80 welterweight rugby tournament in 2005.[219] The national domestic Thailand Rugby Union (TRU) competition includes several universities and services teams such as Chulalongkorn University, Mahasarakham University, Kasetsart University, Prince of Songkla University, Thammasat University, Rangsit University, the Thai Police, the Thai Army, the Thai Navy and the Royal Thai Air Force. Local sports clubs which also compete in the TRU include the British Club of Bangkok, the Southerners Sports Club (Bangkok) and the Royal Bangkok Sports Club.
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Thailand has been called the golf capital of Asia[220] as it is a popular destination for golf. The country attracts a large number of golfers from Japan, Korea, Singapore, South Africa, and Western countries who come to play golf in Thailand every year.[221] The growing popularity of golf, especially among the middle classes and immigrants, is evident as there are more than 200 world-class golf courses nationwide,[222] and some of them are chosen to host PGA and LPGA tournaments, such as Amata Spring Country Club, Alpine Golf and Sports Club, Thai Country Club, and Black Mountain Golf Club.
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Basketball is a growing sport in Thailand, especially on the professional sports club level. The Chang Thailand Slammers won the 2011 ASEAN Basketball League Championship.[223] The Thailand national basketball team had its most successful year at the 1966 Asian Games where it won the silver medal.[224]
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Other sports in Thailand are slowly growing as the country develops its sporting infrastructure. The success in sports like weightlifting and taekwondo at the last two summer Olympic Games has demonstrated that boxing is no longer the only medal option for Thailand.
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The well-known Lumpinee Boxing Stadium originally sited at Rama IV Road near Lumphini Park hosted its final Muay Thai boxing matches on 8 February 2014 after the venue first opened in December 1956. Managed by the Royal Thai Army, the stadium was officially selected for the purpose of muay Thai bouts following a competition that was staged on 15 March 1956. From 11 February 2014, the stadium will relocate to Ram Intra Road, due to the new venue's capacity to accommodate audiences of up to 3,500. Foreigners typically pay between 1,000–2,000 baht to view a match, with prices depending on the location of the seating.[225]
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Thammasat Stadium is a multi-purpose stadium in Bangkok. It is currently used mostly for football matches. The stadium holds 25,000. It is on Thammasat University's Rangsit campus. It was built for the 1998 Asian Games by construction firm Christiani and Nielsen, the same company that constructed the Democracy Monument in Bangkok.
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Rajamangala National Stadium is the biggest sporting arena in Thailand. It currently has a capacity of 65,000. It is in Bang Kapi, Bangkok. The stadium was built in 1998 for the 1998 Asian Games and is the home stadium of the Thailand national football team.
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Thales of Miletus (/ˈθeɪliːz/ THAY-leez; Greek: Θαλῆς (ὁ Μιλήσιος), Thalēs; c. 624/623 – c. 548/545 BC) was a Greek mathematician, astronomer and pre-Socratic philosopher from Miletus in Ionia, Asia Minor. He was one of the Seven Sages of Greece. Many, most notably Aristotle, regarded him as the first philosopher in the Greek tradition,[1][2] and he is otherwise historically recognized as the first individual in Western civilization known to have entertained and engaged in scientific philosophy.[3][4]
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Thales is recognized for breaking from the use of mythology to explain the world and the universe, and instead explaining natural objects and phenomena by naturalistic theories and hypotheses, in a precursor to modern science. Almost all the other pre-Socratic philosophers followed him in explaining nature as deriving from a unity of everything based on the existence of a single ultimate substance, instead of using mythological explanations. Aristotle regarded him as the founder of the Ionian School and reported Thales' hypothesis that the originating principle of nature and the nature of matter was a single material substance: water.[5]
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In mathematics, Thales used geometry to calculate the heights of pyramids and the distance of ships from the shore. He is the first known individual to use deductive reasoning applied to geometry, by deriving four corollaries to Thales' theorem. He is the first known individual to whom a mathematical discovery has been attributed.[6]
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The dates of Thales' life are not exactly known, but are roughly established by a few datable events mentioned in the sources. According to Herodotus, Thales predicted the solar eclipse of May 28, 585 BC.[7] Diogenes Laërtius quotes the chronicle of Apollodorus of Athens as saying that Thales died at the age of 78 during the 58th Olympiad (548–545 BC) and attributes his death to heat stroke while watching the games.[8]
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Thales was probably born in the city of Miletus around the mid-620s BC. The ancient writer Apollodorus of Athens[9] writing during the 2nd century BC,[4] thought Thales was born about the year 625 BC.[9] Herodotus, writing in the fifth century BC, described Thales as "a Phoenician by remote descent".[10] Tim Whitmarsh noted that Thales regarded water as the primal matter, and because thal is the Phoenician word for moisture, his name may have derived from this circumstance."[11]
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According to the later historian Diogenes Laërtius, in his third century AD Lives of the Philosophers, references Herodotus, Duris, and Democritus, who all agree "that Thales was the son of Examyas and Cleobulina, and belonged to the Thelidae who are Phoenicians."[12][13] Their names are indigenous Carian and Greek, respectively.[10] Diogenes then states that "Most writers, however, represent him as a native of Miletus and of a distinguished family."[12][13] However, his supposed mother Cleobulina has also been described as his companion.[14] Diogenes then delivers more conflicting reports: one that Thales married and either fathered a son (Cybisthus or Cybisthon) or adopted his nephew of the same name; the second that he never married, telling his mother as a young man that it was too early to marry, and as an older man that it was too late. Plutarch had earlier told this version: Solon visited Thales and asked him why he remained single; Thales answered that he did not like the idea of having to worry about children. Nevertheless, several years later, anxious for family, he adopted his nephew Cybisthus.[15]
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It has been claimed that he was roughly the professional equivalent of a contemporary option trader.[16]
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It is assumed that Thales at one point in his life visited Egypt, where he learned about geometry.[17] Diogenes Laërtius wrote that Thales identifies the Milesians as Athenian colonists.[18]
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Thales (who died around 30 years before the time of Pythagoras and 300 years before Euclid, Eudoxus of Cnidus, and Eudemus of Rhodes) is often hailed as "the first Greek mathematician".[19] While some historians, such as Colin R. Fletcher, point out that there could have been a predecessor to Thales who would have been named in Eudemus' lost book History of Geometry, it is admitted that without the work "the question becomes mere speculation."[19] Fletcher holds that as there is no viable predecessor to the title of first Greek mathematician, the only question is whether Thales qualifies as a practitioner in that field; he holds that "Thales had at his command the techniques of observation, experimentation, superposition and deduction…he has proved himself mathematician."[19]
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Aristotle wrote in Metaphysics, "Thales, the founder of this school of philosophy, says the permanent entity is water (which is why he also propounded that the earth floats on water). Presumably he derived this assumption from seeing that the nutriment of everything is moist, and that heat itself is generated from moisture and depends upon it for its existence (and that from which a thing is generated is always its first principle). He derived his assumption from this; and also from the fact that the seeds of everything have a moist nature, whereas water is the first principle of the nature of moist things."[5]
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Thales involved himself in many activities, including engineering.[20] Some say that he left no writings. Others say that he wrote On the Solstice and On the Equinox. The Nautical Star-guide has been attributed to him, but this was disputed in ancient times.[21] No writing attributed to him has survived. Diogenes Laërtius quotes two letters from Thales: one to Pherecydes of Syros, offering to review his book on religion, and one to Solon, offering to keep him company on his sojourn from Athens.[clarification needed]
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A story, with different versions, recounts how Thales achieved riches from an olive harvest by prediction of the weather. In one version, he bought all the olive presses in Miletus after predicting the weather and a good harvest for a particular year. Another version of the story has Aristotle explain that Thales had reserved presses in advance, at a discount, and could rent them out at a high price when demand peaked, following his prediction of a particularly good harvest. This first version of the story would constitute the first historically known creation and use of futures, whereas the second version would be the first historically known creation and use of options.[22]
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25 |
+
Aristotle explains that Thales' objective in doing this was not to enrich himself but to prove to his fellow Milesians that philosophy could be useful, contrary to what they thought,[23] or alternatively, Thales had made his foray into enterprise because of a personal challenge put to him by an individual who had asked why, if Thales was an intelligent famous philosopher, he had yet to attain wealth.
|
26 |
+
|
27 |
+
Diogenes Laërtius tells us that Thales gained fame as a counselor when he advised the Milesians not to engage in a symmachia, a "fighting together", with the Lydians. This has sometimes been interpreted as an alliance.[24][failed verification] Another story by Herodotus is that Croesus sent his army to the Persian territory. He was stopped by the river Halys, then unbridged. Thales then got the army across the river by digging a diversion upstream so as to reduce the flow, making it possible to cross the river.[25] While Herodotus reported that most of his fellow Greeks believe that Thales did divert the river Halys to assist King Croesus' military endeavors, he himself finds it doubtful.[26]
|
28 |
+
|
29 |
+
Croesus was defeated before the city of Sardis by Cyrus, who subsequently spared Miletus because it had taken no action. Cyrus was so impressed by Croesus’ wisdom and his connection with the sages that he spared him and took his advice on various matters.[citation needed] The Ionian cities should be demoi, or "districts".
|
30 |
+
|
31 |
+
He counselled them to establish a single seat of government, and pointed out Teos as the fittest place for it; "for that," he said, "was the centre of Ionia. Their other cities might still continue to enjoy their own laws, just as if they were independent states."[27]
|
32 |
+
|
33 |
+
Miletus, however, received favorable terms from Cyrus. The others remained in an Ionian League of twelve cities (excluding Miletus), and were subjugated by the Persians.[citation needed]
|
34 |
+
|
35 |
+
According to Herodotus, Thales predicted the solar eclipse of May 28, 585 BC.[7] Thales also described the position of Ursa Minor, and he thought the constellation might be useful as a guide for navigation at sea. He calculated the duration of the year and the timings of the equinoxes and solstices. He is additionally attributed with the first observation of the Hyades and with calculating the position of the Pleiades.[28] Plutarch indicates that in his day (c. AD 100) there was an extant work, the Astronomy, composed in verse and attributed to Thales.[29]
|
36 |
+
|
37 |
+
Herodotus writes that in the sixth year of the war, the Lydians under King Alyattes and the Medes under Cyaxares were engaged in an indecisive battle when suddenly day turned into night, leading to both parties halting the fighting and negotiating a peace agreement. Herodotus also mentions that the loss of daylight had been predicted by Thales. He does not, however, mention the location of the battle.[30]
|
38 |
+
|
39 |
+
Afterwards, on the refusal of Alyattes to give up his suppliants when Cyaxares sent to demand them of him, war broke out between the Lydians and the Medes, and continued for five years, with various success. In the course of it the Medes gained many victories over the Lydians, and the Lydians also gained many victories over the Medes. Among their other battles there was one night engagement. As, however, the balance had not inclined in favour of either nation, another combat took place in the sixth year, in the course of which, just as the battle was growing warm, day was on a sudden changed into night. This event had been foretold by Thales, the Milesian, who forewarned the Ionians of it, fixing for it the very year in which it actually took place. The Medes and Lydians, when they observed the change, ceased fighting, and were alike anxious to have terms of peace agreed on.[27]
|
40 |
+
|
41 |
+
However, based on the list of Medean kings and the duration of their reign reported elsewhere by Herodotus, Cyaxares died 10 years before the eclipse.[31][32]
|
42 |
+
|
43 |
+
Diogenes Laërtius[33] tells us that the Seven Sages were created in the archonship of Damasius at Athens about 582 BC and that Thales was the first sage. The same story, however, asserts that Thales emigrated to Miletus. There is also a report that he did not become a student of nature until after his political career. Much as we would like to have a date on the seven sages, we must reject these stories and the tempting date if we are to believe that Thales was a native of Miletus, predicted the eclipse, and was with Croesus in the campaign against Cyrus.
|
44 |
+
|
45 |
+
Thales received instruction from an Egyptian priest.[citation needed] It was fairly certain that he came from a wealthy, established family, in a class which customarily provided higher education for their children.[citation needed] Moreover, the ordinary citizen, unless he was a seafaring man or a merchant, could not afford the grand tour in Egypt, and did not consort with noble lawmakers such as Solon.[citation needed]
|
46 |
+
|
47 |
+
In Diogenes Laërtius' Lives of Eminent Philosophers Chapter 1.39, Laërtius relates several stories of an expensive object that is to go to the most wise. In one version (that Laërtius credits to Callimachus in his Iambics) Bathycles of Arcadia states in his will that an expensive bowl "'should be given to him who had done most good by his wisdom.' So it was given to Thales, went the round of all the sages, and came back to Thales again. And he sent it to Apollo at Didyma, with this dedication...'Thales the Milesian, son of Examyas [dedicates this] to Delphinian Apollo after twice winning the prize from all the Greeks.'"[34]
|
48 |
+
|
49 |
+
Early Greeks, and other civilizations before them, often invoked idiosyncratic explanations of natural phenomena with reference to the will of anthropomorphic gods and heroes. Instead, Thales aimed to explain natural phenomena via rational hypotheses that referenced natural processes themselves. For example, rather than assuming that earthquakes were the result of supernatural whims, Thales explained them by hypothesizing that the Earth floats on water and that earthquakes occur when the Earth is rocked by waves.[35][36]
|
50 |
+
|
51 |
+
Thales was a hylozoist (one who thinks that matter is alive,[37] i.e. containing soul(s)). Aristotle wrote (De Anima 411 a7-8) of Thales: ...Thales thought all things are full of gods. Aristotle posits the origin of Thales thought on matter generally containing souls, to Thales thinking initially on the fact of, because magnets move iron, the presence of movement of matter indicated this matter contained life.[38][39]
|
52 |
+
|
53 |
+
Thales, according to Aristotle, asked what was the nature (Greek arche) of the object so that it would behave in its characteristic way. Physis (φύσις) comes from phyein (φύειν), "to grow", related to our word "be".[40][41] (G)natura is the way a thing is "born",[42] again with the stamp of what it is in itself.
|
54 |
+
|
55 |
+
Aristotle characterizes most of the philosophers "at first" (πρῶτον) as thinking that the "principles in the form of matter were the only principles of all things", where "principle" is arche, "matter" is hyle ("wood" or "matter", "material") and "form" is eidos.[43]
|
56 |
+
|
57 |
+
Arche is translated as "principle", but the two words do not have precisely the same meaning. A principle of something is merely prior (related to pro-) to it either chronologically or logically. An arche (from ἄρχειν, "to rule") dominates an object in some way. If the arche is taken to be an origin, then specific causality is implied; that is, B is supposed to be characteristically B just because it comes from A, which dominates it.
|
58 |
+
|
59 |
+
The archai that Aristotle had in mind in his well-known passage on the first Greek scientists are not necessarily chronologically prior to their objects, but are constituents of it. For example, in pluralism objects are composed of earth, air, fire and water, but those elements do not disappear with the production of the object. They remain as archai within it, as do the atoms of the atomists.
|
60 |
+
|
61 |
+
What Aristotle is really saying is that the first philosophers were trying to define the substance(s) of which all material objects are composed. As a matter of fact, that is exactly what modern scientists are attempting to accomplish in nuclear physics, which is a second reason why Thales is described as the first western scientist,[citation needed] but some contemporary scholars reject this interpretation.[44]
|
62 |
+
|
63 |
+
Thales was known for his innovative use of geometry. His understanding was theoretical as well as practical. For example, he said:
|
64 |
+
|
65 |
+
Megiston topos: apanta gar chorei (Μέγιστον τόπος· ἄπαντα γὰρ χωρεῖ.)
|
66 |
+
|
67 |
+
The greatest is space, for it holds all things.[45]
|
68 |
+
|
69 |
+
Topos is in Newtonian-style space, since the verb, chorei, has the connotation of yielding before things, or spreading out to make room for them, which is extension. Within this extension, things have a position. Points, lines, planes and solids related by distances and angles follow from this presumption.
|
70 |
+
|
71 |
+
Thales understood similar triangles and right triangles, and what is more, used that knowledge in practical ways. The story is told in Diogenes Laërtius (loc. cit.) that he measured the height of the pyramids by their shadows at the moment when his own shadow was equal to his height. A right triangle with two equal legs is a 45-degree right triangle, all of which are similar. The length of the pyramid's shadow measured from the center of the pyramid at that moment must have been equal to its height.
|
72 |
+
|
73 |
+
This story indicates that he was familiar with the Egyptian seked, or seqed, the ratio of the run to the rise of a slope (cotangent).[citation needed] The seked is at the base of problems 56, 57, 58, 59 and 60 of the Rhind papyrus — an ancient Egyptian mathematical document.
|
74 |
+
|
75 |
+
More practically Thales used the same method to measure the distances of ships at sea, said Eudemus as reported by Proclus ("in Euclidem"). According to Kirk & Raven (reference cited below), all you need for this feat is three straight sticks pinned at one end and knowledge of your altitude. One stick goes vertically into the ground. A second is made level. With the third you sight the ship and calculate the seked from the height of the stick and its distance from the point of insertion to the line of sight (Proclus, In Euclidem, 352).
|
76 |
+
|
77 |
+
There are two theorems of Thales in elementary geometry, one known as Thales' theorem having to do with a triangle inscribed in a circle and having the circle's diameter as one leg, the other theorem being also called the intercept theorem. In addition Eudemus attributed to him the discovery that a circle is bisected by its diameter, that the base angles of an isosceles triangle are equal and that vertical angles are equal. According to a historical Note,[46] when Thales visited Egypt,[17] he observed that whenever the Egyptians drew two intersecting lines, they would measure the vertical angles to make sure that they were equal. Thales concluded that one could prove that all vertical angles are equal if one accepted some general notions such as: all straight angles are equal, equals added to equals are equal, and equals subtracted from equals are equal.
|
78 |
+
|
79 |
+
The evidence for the primacy of Thales comes to us from a book by Proclus who wrote a thousand years after Thales but is believed to have had a copy of Eudemus' book. Proclus wrote "Thales was the first to go to Egypt and bring back to Greece this study."[19] He goes on to tell us that in addition to applying the knowledge he gained in Egypt "He himself discovered many propositions and disclosed the underlying principles of many others to his successors, in some case his method being more general, in others more empirical."[19]
|
80 |
+
|
81 |
+
Other quotes from Proclus list more of Thales' mathematical achievements:
|
82 |
+
|
83 |
+
They say that Thales was the first to demonstrate that the circle is bisected by the diameter, the cause of the bisection being the unimpeded passage of the straight line through the centre.[19]
|
84 |
+
|
85 |
+
[Thales] is said to have been the first to have known and to have enunciated [the theorem] that the angles at the base of any isosceles triangle are equal, though in the more archaic manner he described the equal angles as similar.[19]
|
86 |
+
|
87 |
+
This theorem, that when two straight lines cut one another, the vertical and opposite angles are equal, was first discovered, as Eudemus says, by Thales, though the scientific demonstration was improved by the writer of Elements.[19]
|
88 |
+
|
89 |
+
Eudemus in his History of Geometry attributes this theorem [the equality of triangles having two angles and one side equal] to Thales. For he says that the method by which Thales showed how to find the distance of ships at sea necessarily involves this method.[19]
|
90 |
+
|
91 |
+
Pamphila says that, having learnt geometry from the Egyptians, he [Thales] was the first to inscribe in a circle a right-angled triangle, whereupon he sacrificed an ox.[19]
|
92 |
+
|
93 |
+
In addition to Proclus, Hieronymus of Rhodes also cites Thales as the first Greek mathematician. Hieronymus held that Thales was able to measure the height of the pyramids by using a theorem of geometry now known as the intercept theorem, (after gathering data by using his walking-stick and comparing its shadow to those cast by the pyramids). We receive variations of Hieronymus' story through Diogenes Laërtius,[47] Pliny the Elder, and Plutarch.[19][48] According to Hieronymus, historically quoted by Diogenes Laërtius, Thales found the height of pyramids by comparison between the lengths of the shadows cast by a person and by the pyramids.[49]
|
94 |
+
|
95 |
+
Due to the variations among testimonies, such as the "story of the sacrifice of an ox on the occasion of the discovery that the angle on a diameter of a circle is a right angle" in the version told by Diogenes Laërtius being accredited to Pythagoras rather than Thales, some historians (such as D. R. Dicks) question whether such anecdotes have any historical worth whatsoever.[26]
|
96 |
+
|
97 |
+
Thales' most famous philosophical position was his cosmological thesis, which comes down to us through a passage from Aristotle's Metaphysics.[50] In the work Aristotle unequivocally reported Thales' hypothesis about the nature of all matter – that the originating principle of nature was a single material substance: water. Aristotle then proceeded to proffer a number of conjectures based on his own observations to lend some credence to why Thales may have advanced this idea (though Aristotle did not hold it himself).
|
98 |
+
|
99 |
+
Aristotle laid out his own thinking about matter and form which may shed some light on the ideas of Thales, in Metaphysics 983 b6 8–11, 17–21. (The passage contains words that were later adopted by science with quite different meanings.)
|
100 |
+
|
101 |
+
That from which is everything that exists and from which it first becomes and into which it is rendered at last, its substance remaining under it, but transforming in qualities, that they say is the element and principle of things that are. …For it is necessary that there be some nature (φύσις), either one or more than one, from which become the other things of the object being saved... Thales the founder of this type of philosophy says that it is water.
|
102 |
+
|
103 |
+
In this quote we see Aristotle's depiction of the problem of change and the definition of substance. He asked if an object changes, is it the same or different? In either case how can there be a change from one to the other? The answer is that the substance "is saved", but acquires or loses different qualities (πάθη, the things you "experience").
|
104 |
+
|
105 |
+
Aristotle conjectured that Thales reached his conclusion by contemplating that the "nourishment of all things is moist and that even the hot is created from the wet and lives by it." While Aristotle's conjecture on why Thales held water as the originating principle of matter is his own thinking, his statement that Thales held it as water is generally accepted as genuinely originating with Thales and he is seen as an incipient matter-and-formist.[citation needed]
|
106 |
+
|
107 |
+
Thales thought the Earth must be a flat disk which is floating in an expanse of water.[51]
|
108 |
+
|
109 |
+
Heraclitus Homericus states that Thales drew his conclusion from seeing moist substance turn into air, slime and earth. It seems likely that Thales viewed the Earth as solidifying from the water on which it floated and the oceans that surround it.
|
110 |
+
|
111 |
+
Writing centuries later, Diogenes Laërtius also states that Thales taught "Water constituted (ὑπεστήσατο, 'stood under') the principle of all things."[52]
|
112 |
+
|
113 |
+
Aristotle considered Thales’ position to be roughly the equivalent to the later ideas of Anaximenes, who held that everything was composed of air.[53]
|
114 |
+
The 1870 book Dictionary of Greek and Roman Biography and Mythology noted:[2]
|
115 |
+
|
116 |
+
In his dogma that water is the origin of things, that is, that it is that out of which every thing arises, and into which every thing resolves itself, Thales may have followed Orphic cosmogonies, while, unlike them, he sought to establish the truth of the assertion. Hence, Aristotle, immediately after he has called him the originator of philosophy brings forward the reasons which Thales was believed to have adduced in confirmation of that assertion; for that no written development of it, or indeed any book by Thales, was extant, is proved by the expressions which Aristotle uses when he brings forward the doctrines and proofs of the Milesian. (p. 1016)
|
117 |
+
|
118 |
+
According to Aristotle, Thales thought lodestones had souls, because iron is attracted to them (by the force of magnetism).[54]
|
119 |
+
|
120 |
+
Aristotle defined the soul as the principle of life, that which imbues the matter and makes it live, giving it the animation, or power to act. The idea did not originate with him, as the Greeks in general believed in the distinction between mind and matter, which was ultimately to lead to a distinction not only between body and soul but also between matter and energy.[citation needed] If things were alive, they must have souls. This belief was no innovation, as the ordinary ancient populations of the Mediterranean did believe that natural actions were caused by divinities. Accordingly, Aristotle and other ancient writers state that Thales believed that "all things were full of gods."[55][56] In their zeal to make him the first in everything some said he was the first to hold the belief, which must have been widely known to be false.[citation needed] However, Thales was looking for something more general, a universal substance of mind.[citation needed] That also was in the polytheism of the times. Zeus was the very personification of supreme mind, dominating all the subordinate manifestations. From Thales on, however, philosophers had a tendency to depersonify or objectify mind, as though it were the substance of animation per se and not actually a god like the other gods. The end result was a total removal of mind from substance, opening the door to a non-divine principle of action.[citation needed]
|
121 |
+
|
122 |
+
Classical thought, however, had proceeded only a little way along that path. Instead of referring to the person, Zeus, they talked about the great mind:
|
123 |
+
|
124 |
+
"Thales", says Cicero,[57] "assures that water is the principle of all things; and that God is that Mind which shaped and created all things from water."
|
125 |
+
|
126 |
+
The universal mind appears as a Roman belief in Virgil as well:
|
127 |
+
|
128 |
+
In the beginning, SPIRIT within (spiritus intus) strengthens Heaven and Earth,
|
129 |
+
The watery fields, and the lucid globe of Luna, and then –
|
130 |
+
Titan stars; and mind (mens) infused through the limbs
|
131 |
+
Agitates the whole mass, and mixes itself with GREAT MATTER (magno corpore)[58]
|
132 |
+
|
133 |
+
According to Henry Fielding (1775), Diogenes Laërtius (1.35) affirmed that Thales posed "the independent pre-existence of God from all eternity, stating "that God was the oldest of all beings, for he existed without a previous cause even in the way of generation; that the world was the most beautiful of all things; for it was created by God."[59]
|
134 |
+
|
135 |
+
Due to the scarcity of sources concerning Thales and the discrepancies between the accounts given in the sources that have survived, there is a scholarly debate over possible influences on Thales and the Greek mathematicians that came after him.
|
136 |
+
Historian Roger L. Cooke points out that Proclus does not make any mention of Mesopotamian influence on Thales or Greek geometry, but "is shown clearly in Greek astronomy, in the use of sexagesimal system of measuring angles and in Ptolemy's explicit use of Mesopotamian astronomical observations."[60] Cooke notes that it may possibly also appear in the second book of Euclid's Elements, "which contains geometric constructions equivalent to certain algebraic relations that are frequently encountered in the cuneiform tablets." Cooke notes "This relation, however, is controversial."[60]
|
137 |
+
|
138 |
+
Historian B.L. Van der Waerden is among those advocating the idea of Mesopotamian influence, writing "It follows that we have to abandon the traditional belief that the oldest Greek mathematicians discovered geometry entirely by themselves…a belief that was tenable only as long as nothing was known about Babylonian mathematics. This in no way diminishes the stature of Thales; on the contrary, his genius receives only now the honour that is due to it, the honour of having developed a logical structure for geometry, of having introduced proof into geometry."[19]
|
139 |
+
|
140 |
+
Some historians, such as D. R. Dicks takes issue with the idea that we can determine from the questionable sources we have, just how influenced Thales was by Babylonian sources. He points out that while Thales is held to have been able to calculate an eclipse using a cycle called the "Saros" held to have been "borrowed from the Babylonians", "The Babylonians, however, did not use cycles to predict solar eclipses, but computed them from observations of the latitude of the moon made shortly before the expected syzygy."[26] Dicks cites historian O. Neugebauer who relates that "No Babylonian theory for predicting solar eclipse existed at 600 B.C., as one can see from the very unsatisfactory situation 400 years later; nor did the Babylonians ever develop any theory which took the influence of geographical latitude into account." Dicks examines the cycle referred to as 'Saros' – which Thales is held to have used and which is believed to stem from the Babylonians. He points out that Ptolemy makes use of this and another cycle in his book Mathematical Syntaxis but attributes it to Greek astronomers earlier than Hipparchus and not to Babylonians.[26] Dicks notes Herodotus does relate that Thales made use of a cycle to predict the eclipse, but maintains that "if so, the fulfillment of the 'prediction' was a stroke of pure luck not science".[26] He goes further joining with other historians (F. Martini, J.L. E. Dreyer, O. Neugebauer) in rejecting the historicity of the eclipse story altogether.[26] Dicks links the story of Thales discovering the cause for a solar eclipse with Herodotus' claim that Thales discovered the cycle of the sun with relation to the solstices, and concludes "he could not possibly have possessed this knowledge which neither the Egyptians nor the Babylonians nor his immediate successors possessed."[26] Josephus is the only ancient historian that claims Thales visited Babylonia.
|
141 |
+
|
142 |
+
Herodotus wrote that the Greeks learnt the practice of dividing the day into 12 parts, about the polos, and the gnomon from the Babylonians. (The exact meaning of his use of the word polos is unknown, current theories include: "the heavenly dome", "the tip of the axis of the celestial sphere", or a spherical concave sundial.) Yet even Herodotus' claims on Babylonian influence are contested by some modern historians, such as L. Zhmud, who points out that the division of the day into twelve parts (and by analogy the year) was known to the Egyptians already in the second millennium, the gnomon was known to both Egyptians and Babylonians, and the idea of the "heavenly sphere" was not used outside of Greece at this time.[61]
|
143 |
+
|
144 |
+
Less controversial than the position that Thales learnt Babylonian mathematics is the claim he was influenced by Egyptians. Pointedly historian S. N. Bychkov holds that the idea that the base angles of an isosceles triangle are equal likely came from Egypt. This is because, when building a roof for a home – having a cross section be exactly an isosceles triangle isn't crucial (as it's the ridge of the roof that must fit precisely), in contrast a symmetric square pyramid cannot have errors in the base angles of the faces or they will not fit together tightly.[60]
|
145 |
+
Historian D.R. Dicks agrees that compared to the Greeks in the era of Thales, there was a more advanced state of mathematics among the Babylonians and especially the Egyptians – "both cultures knew the correct formulae for determining the areas and volumes of simple geometrical figures such as triangles, rectangles, trapezoids, etc.; the Egyptians could also calculate correctly the volume of the frustum of a pyramid with a square base (the Babylonians used an incorrect formula for this), and used a formula for the area of a circle...which gives a value for π of 3.1605—a good approximation."[26] Dicks also agrees that this would have had an effect on Thales (whom the most ancient sources agree was interested in mathematics and astronomy) but he holds that tales of Thales' travels in these lands are pure myth.
|
146 |
+
|
147 |
+
The ancient civilization and massive monuments of Egypt had "a profound and ineradicable impression on the Greeks". They attributed to Egyptians "an immemorial knowledge of certain subjects" (including geometry) and would claim Egyptian origin for some of their own ideas to try and lend them "a respectable antiquity" (such as the "Hermetic" literature of the Alexandrian period).[26]
|
148 |
+
|
149 |
+
Dicks holds that since Thales was a prominent figure in Greek history by the time of Eudemus but "nothing certain was known except that he lived in Miletus".[26] A tradition developed that as "Milesians were in a position to be able to travel widely" Thales must have gone to Egypt.[26] As Herodotus says Egypt was the birthplace of geometry he must have learnt that while there. Since he had to have been there, surely one of the theories on Nile Flooding laid out by Herodotus must have come from Thales. Likewise as he must have been in Egypt he had to have done something with the Pyramids – thus the tale of measuring them. Similar apocryphal stories exist of Pythagoras and Plato traveling to Egypt with no corroborating evidence.
|
150 |
+
|
151 |
+
As the Egyptian and Babylonian geometry at the time was "essentially arithmetical", they used actual numbers and "the procedure is then described with explicit instructions as to what to do with these numbers" there was no mention of how the rules of procedure were made, and nothing toward a logically arranged corpus of generalized geometrical knowledge with analytical 'proofs' such as we find in the words of Euclid, Archimedes, and Apollonius."[26] So even had Thales traveled there he could not have learnt anything about the theorems he is held to have picked up there (especially because there is no evidence that any Greeks of this age could use Egyptian hieroglyphics).[26]
|
152 |
+
|
153 |
+
Likewise until around the second century BC and the time of Hipparchus (c. 190–120 BC) the Babylonian general division of the circle into 360 degrees and their sexagesimal system was unknown.[26] Herodotus says almost nothing about Babylonian literature and science, and very little about their history. Some historians, like P. Schnabel, hold that the Greeks only learned more about Babylonian culture from Berossus, a Babylonian priest who is said to have set up a school in Cos around 270 BC (but to what extent this had in the field of geometry is contested).
|
154 |
+
|
155 |
+
Dicks points out that the primitive state of Greek mathematics and astronomical ideas exhibited by the peculiar notions of Thales' successors (such as Anaximander, Anaximenes, Xenophanes, and Heraclitus), which historian J. L. Heiberg calls "a mixture of brilliant intuition and childlike analogies",[62] argues against the assertions from writers in late antiquity that Thales discovered and taught advanced concepts in these fields.
|
156 |
+
|
157 |
+
John Burnet (1892) noted[63]
|
158 |
+
|
159 |
+
Lastly, we have one admitted instance of a philosophic guild, that of the Pythagoreans. And it will be found that the hypothesis, if it is to be called by that name, of a regular organisation of scientific activity will alone explain all the facts. The development of doctrine in the hands of Thales, Anaximander, and Anaximenes, for instance, can only be understood as the elaboration of a single idea in a school with a continuous tradition.
|
160 |
+
|
161 |
+
In the long sojourn of philosophy, there has existed hardly a philosopher or historian of philosophy who did not mention Thales and try to characterize him in some way. He is generally recognized as having brought something new to human thought. Mathematics, astronomy, and medicine already existed. Thales added something to these different collections of knowledge to produce a universality, which, as far as writing tells us, was not in tradition before, but resulted in a new field.
|
162 |
+
|
163 |
+
Ever since, interested persons have been asking what that new something is. Answers fall into (at least) two categories, the theory and the method. Once an answer has been arrived at, the next logical step is to ask how Thales compares to other philosophers, which leads to his classification (rightly or wrongly).
|
164 |
+
|
165 |
+
The most natural epithets of Thales are "materialist" and "naturalist", which are based on ousia and physis. The Catholic Encyclopedia notes that Aristotle called him a physiologist, with the meaning "student of nature."[64] On the other hand, he would have qualified as an early physicist, as did Aristotle. They studied corpora, "bodies", the medieval descendants of substances.
|
166 |
+
|
167 |
+
Most agree that Thales' stamp on thought is the unity of substance, hence Bertrand Russell:[65]
|
168 |
+
|
169 |
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The view that all matter is one is quite a reputable scientific hypothesis.
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...But it is still a handsome feat to have discovered that a substance remains the same in different states of aggregation.
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Russell was only reflecting an established tradition; for example: Nietzsche, in his Philosophy in the Tragic Age of the Greeks, wrote:[66]
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Greek philosophy seems to begin with an absurd notion, with the proposition that water is the primal origin and the womb of all things. Is it really necessary for us to take serious notice of this proposition? It is, and for three reasons. First, because it tells us something about the primal origin of all things; second, because it does so in language devoid of image or fable, and finally, because contained in it, if only embryonically, is the thought, "all things are one."
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This sort of materialism, however, should not be confused with deterministic materialism. Thales was only trying to explain the unity observed in the free play of the qualities. The arrival of uncertainty in the modern world made possible a return to Thales; for example, John Elof Boodin writes ("God and Creation"):
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We cannot read the universe from the past...
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Boodin defines an "emergent" materialism, in which the objects of sense emerge uncertainly from the substrate. Thales is the innovator of this sort of materialism.
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Later scholastic thinkers would maintain that in his choice of water Thales was influenced by Babylonian or Chaldean religion, that held that a god had begun creation by acting upon the pre-existing water. Historian Abraham Feldman holds this does not stand up under closer examination. In Babylonian religion the water is lifeless and sterile until a god acts upon it, but for Thales water itself was divine and creative. He maintained that "All things are full of gods", and to understand the nature of things was to discover the secrets of the deities, and through this knowledge open the possibility that one could be greater than the grandest Olympian.[67]
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Feldman points out that while other thinkers recognized the wetness of the world "none of them was inspired to conclude that everything was ultimately aquatic."[67] He further points out that Thales was "a wealthy citizen of the fabulously rich Oriental port of Miletus...a dealer in the staples of antiquity, wine and oil...He certainly handled the shell-fish of the Phoenicians that secreted the dye of imperial purple."[67] Feldman recalls the stories of Thales measuring the distance of boats in the harbor, creating mechanical improvements for ship navigation, giving an explanation for the flooding of the Nile (vital to Egyptian agriculture and Greek trade), and changing the course of the river Halys so an army could ford it. Rather than seeing water as a barrier Thales contemplated the Ionian yearly religious gathering for athletic ritual (held on the promontory of Mycale and believed to be ordained by the ancestral kindred of Poseidon, the god of the sea). He called for the Ionian mercantile states participating in this ritual to convert it into a democratic federation under the protection of Poseidon that would hold off the forces of pastoral Persia. Feldman concludes that Thales saw "that water was a revolutionary leveler and the elemental factor determining the subsistence and business of the world"[67] and "the common channel of states."[67]
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Feldman considers Thales' environment and holds that Thales would have seen tears, sweat, and blood as granting value to a person's work and the means how life giving commodities travelled (whether on bodies of water or through the sweat of slaves and pack-animals). He would have seen that minerals could be processed from water such as life-sustaining salt and gold taken from rivers. He would’ve seen fish and other food stuffs gathered from it. Feldman points out that Thales held that the lodestone was alive as it drew metals to itself. He holds that Thales "living ever in sight of his beloved sea" would see water seem to draw all "traffic in wine and oil, milk and honey, juices and dyes" to itself, leading him to "a vision of the universe melting into a single substance that was valueless in itself and still the source of wealth."[67] Feldman concludes that for Thales "...water united all things. The social significance of water in the time of Thales induced him to discern through hardware and dry-goods, through soil and sperm, blood, sweat and tears, one fundamental fluid stuff...water, the most commonplace and powerful material known to him."[67] This combined with his contemporary's idea of "spontaneous generation" allow us to see how Thales could hold that water could be divine and creative.
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Feldman points to the lasting association of the theory that "all whatness is wetness" with Thales himself, pointing out that Diogenes Laërtius speaks of a poem, probably a satire, where Thales is snatched to heaven by the sun.[67]
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In the West, Thales represents a new kind of inquiring community as well. Edmund Husserl[68] attempts to capture the new movement as follows. Philosophical man is a "new cultural configuration" based in stepping back from "pregiven tradition" and taking up a rational "inquiry into what is true in itself;" that is, an ideal of truth. It begins with isolated individuals such as Thales, but they are supported and cooperated with as time goes on. Finally the ideal transforms the norms of society, leaping across national borders.
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The term "Pre-Socratic" derives ultimately from the philosopher Aristotle, who distinguished the early philosophers as concerning themselves with substance.
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Diogenes Laërtius on the other hand took a strictly geographic and ethnic approach. Philosophers were either Ionian or Italian. He used "Ionian" in a broader sense, including also the Athenian academics, who were not Pre-Socratics. From a philosophic point of view, any grouping at all would have been just as effective. There is no basis for an Ionian or Italian unity. Some scholars, however, concede to Diogenes' scheme as far as referring to an "Ionian" school. There was no such school in any sense.
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The most popular approach refers to a Milesian school, which is more justifiable socially and philosophically. They sought for the substance of phenomena and may have studied with each other. Some ancient writers qualify them as Milesioi, "of Miletus."
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Thales had a profound influence on other Greek thinkers and therefore on Western history. Some believe Anaximander was a pupil of Thales. Early sources report that one of Anaximander's more famous pupils, Pythagoras, visited Thales as a young man, and that Thales advised him to travel to Egypt to further his philosophical and mathematical studies.
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Many philosophers followed Thales' lead in searching for explanations in nature rather than in the supernatural; others returned to supernatural explanations, but couched them in the language of philosophy rather than of myth or of religion.
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Looking specifically at Thales' influence during the pre-Socratic era, it is clear that he stood out as one of the first thinkers who thought more in the way of logos than mythos. The difference between these two more profound ways of seeing the world is that mythos is concentrated around the stories of holy origin, while logos is concentrated around the argumentation. When the mythical man wants to explain the world the way he sees it, he explains it based on gods and powers. Mythical thought does not differentiate between things and persons[citation needed] and furthermore it does not differentiate between nature and culture[citation needed]. The way a logos thinker would present a world view is radically different from the way of the mythical thinker. In its concrete form, logos is a way of thinking not only about individualism[clarification needed], but also the abstract[clarification needed]. Furthermore, it focuses on sensible and continuous argumentation. This lays the foundation of philosophy and its way of explaining the world in terms of abstract argumentation, and not in the way of gods and mythical stories[citation needed].
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Because of Thales' elevated status in Greek culture an intense interest and admiration followed his reputation. Due to this following, the oral stories about his life were open to amplification and historical fabrication, even before they were written down generations later. Most modern dissension comes from trying to interpret what we know, in particular, distinguishing legend from fact.
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Historian D.R. Dicks and other historians divide the ancient sources about Thales into those before 320 BC and those after that year (some such as Proclus writing in the 5th century C.E. and Simplicius of Cilicia in the 6th century C.E. writing nearly a millennium after his era).[26] The first category includes Herodotus, Plato, Aristotle, Aristophanes, and Theophrastus among others. The second category includes Plautus, Aetius, Eusebius, Plutarch, Josephus, Iamblichus, Diogenes Laërtius, Theon of Smyrna, Apuleius, Clement of Alexandria, Pliny the Elder, and John Tzetzes among others.
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The earliest sources on Thales (living before 320 BC) are often the same for the other Milesian philosophers (Anaximander, and Anaximenes). These sources were either roughly contemporaneous (such as Herodotus) or lived within a few hundred years of his passing. Moreover, they were writing from an oral tradition that was widespread and well known in the Greece of their day.
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The latter sources on Thales are several "ascriptions of commentators and compilers who lived anything from 700 to 1,000 years after his death"[26] which include "anecdotes of varying degrees of plausibility"[26] and in the opinion of some historians (such as D. R. Dicks) of "no historical worth whatsoever".[26] Dicks points out that there is no agreement "among the 'authorities' even on the most fundamental facts of his life—e.g. whether he was a Milesian or a Phoenician, whether he left any writings or not, whether he was married or single-much less on the actual ideas and achievements with which he is credited."[26]
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Contrasting the work of the more ancient writers with those of the later, Dicks points out that in the works of the early writers Thales and the other men who would be hailed as "the Seven Sages of Greece" had a different reputation than that which would be assigned to them by later authors. Closer to their own era, Thales, Solon, Bias of Priene, Pittacus of Mytilene and others were hailed as "essentially practical men who played leading roles in the affairs of their respective states, and were far better known to the earlier Greeks as lawgivers and statesmen than as profound thinkers and philosophers."[26] For example, Plato praises him (coupled with Anacharsis) for being the originator of the potter's wheel and the anchor.
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Only in the writings of the second group of writers (working after 320 BC) do "we obtain the picture of Thales as the pioneer in Greek scientific thinking, particularly in regard to mathematics and astronomy which he is supposed to have learnt about in Babylonia and Egypt."[26] Rather than "the earlier tradition [where] he is a favourite example of the intelligent man who possesses some technical 'know how'...the later doxographers [such as Dicaearchus in the latter half of the fourth century BC] foist on to him any number of discoveries and achievements, in order to build him up as a figure of superhuman wisdom."[26]
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Dicks points out a further problem arises in the surviving information on Thales, for rather than using ancient sources closer to the era of Thales, the authors in later antiquity ("epitomators, excerptors, and compilers"[26]) actually "preferred to use one or more intermediaries, so that what we actually read in them comes to us not even at second, but at third or fourth or fifth hand. ...Obviously this use of intermediate sources, copied and recopied from century to century, with each writer adding additional pieces of information of greater or less plausibility from his own knowledge, provided a fertile field for errors in transmission, wrong ascriptions, and fictitious attributions".[26] Dicks points out that "certain doctrines that later commentators invented for Thales...were then accepted into the biographical tradition" being copied by subsequent writers who were then cited by those coming after them "and thus, because they may be repeated by different authors relying on different sources, may produce an illusory impression of genuineness."[26]
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Doubts even exist when considering the philosophical positions held to originate in Thales "in reality these stem directly from Aristotle's own interpretations which then became incorporated in the doxographical tradition as erroneous ascriptions to Thales".[26] (The same treatment was given by Aristotle to Anaxagoras.)
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Most philosophic analyses of the philosophy of Thales come from Aristotle, a professional philosopher, tutor of Alexander the Great, who wrote 200 years after Thales' death. Aristotle, judging from his surviving books, does not seem to have access to any works by Thales, although he probably had access to works of other authors about Thales, such as Herodotus, Hecataeus, Plato etc., as well as others whose work is now extinct. It was Aristotle's express goal to present Thales' work not because it was significant in itself, but as a prelude to his own work in natural philosophy.[69] Geoffrey Kirk and John Raven, English compilers of the fragments of the Pre-Socratics, assert that Aristotle's "judgments are often distorted by his view of earlier philosophy as a stumbling progress toward the truth that Aristotle himself revealed in his physical doctrines."[70] There was also an extensive oral tradition. Both the oral and the written were commonly read or known by all educated men in the region.
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Aristotle's philosophy had a distinct stamp: it professed the theory of matter and form, which modern scholastics have dubbed hylomorphism. Though once very widespread, it was not generally adopted by rationalist and modern science, as it mainly is useful in metaphysical analyses, but does not lend itself to the detail that is of interest to modern science. It is not clear that the theory of matter and form existed as early as Thales, and if it did, whether Thales espoused it.
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While some historians, like B. Snell, maintain that Aristotle was relying on a pre-Platonic written record by Hippias rather than oral tradition, this is a controversial position. Representing the scholarly consensus Dicks states that "the tradition about him even as early as the fifth century B.C., was evidently based entirely on hearsay....It would seem that already by Aristotle's time the early Ionians were largely names only to which popular tradition attached various ideas or achievements with greater or less plausibility".[26] He points out that works confirmed to have existed in the sixth century BC by Anaximander and Xenophanes had already disappeared by the fourth century BC, so the chances of Pre-Socratic material surviving to the age of Aristotle is almost nil (even less likely for Aristotle's pupils Theophrastus and Eudemus and less likely still for those following after them).
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The main secondary source concerning the details of Thales' life and career is Diogenes Laërtius, "Lives of Eminent Philosophers".[71] This is primarily a biographical work, as the name indicates. Compared to Aristotle, Diogenes is not much of a philosopher. He is the one who, in the Prologue to that work, is responsible for the division of the early philosophers into "Ionian" and "Italian", but he places the Academics in the Ionian school and otherwise evidences considerable disarray and contradiction, especially in the long section on forerunners of the "Ionian School". Diogenes quotes two letters attributed to Thales, but Diogenes wrote some eight centuries after Thales' death and that his sources often contained "unreliable or even fabricated information",[72] hence the concern for separating fact from legend in accounts of Thales.
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It is due to this use of hearsay and a lack of citing original sources that leads some historians, like Dicks and Werner Jaeger, to look at the late origin of the traditional picture of Pre-Socratic philosophy and view the whole idea as a construct from a later age, "the whole picture that has come down to us of the history of early philosophy was fashioned during the two or three generations from Plato to the immediate pupils of Aristotle".[73]
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Thales of Miletus (/ˈθeɪliːz/ THAY-leez; Greek: Θαλῆς (ὁ Μιλήσιος), Thalēs; c. 624/623 – c. 548/545 BC) was a Greek mathematician, astronomer and pre-Socratic philosopher from Miletus in Ionia, Asia Minor. He was one of the Seven Sages of Greece. Many, most notably Aristotle, regarded him as the first philosopher in the Greek tradition,[1][2] and he is otherwise historically recognized as the first individual in Western civilization known to have entertained and engaged in scientific philosophy.[3][4]
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Thales is recognized for breaking from the use of mythology to explain the world and the universe, and instead explaining natural objects and phenomena by naturalistic theories and hypotheses, in a precursor to modern science. Almost all the other pre-Socratic philosophers followed him in explaining nature as deriving from a unity of everything based on the existence of a single ultimate substance, instead of using mythological explanations. Aristotle regarded him as the founder of the Ionian School and reported Thales' hypothesis that the originating principle of nature and the nature of matter was a single material substance: water.[5]
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In mathematics, Thales used geometry to calculate the heights of pyramids and the distance of ships from the shore. He is the first known individual to use deductive reasoning applied to geometry, by deriving four corollaries to Thales' theorem. He is the first known individual to whom a mathematical discovery has been attributed.[6]
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The dates of Thales' life are not exactly known, but are roughly established by a few datable events mentioned in the sources. According to Herodotus, Thales predicted the solar eclipse of May 28, 585 BC.[7] Diogenes Laërtius quotes the chronicle of Apollodorus of Athens as saying that Thales died at the age of 78 during the 58th Olympiad (548–545 BC) and attributes his death to heat stroke while watching the games.[8]
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Thales was probably born in the city of Miletus around the mid-620s BC. The ancient writer Apollodorus of Athens[9] writing during the 2nd century BC,[4] thought Thales was born about the year 625 BC.[9] Herodotus, writing in the fifth century BC, described Thales as "a Phoenician by remote descent".[10] Tim Whitmarsh noted that Thales regarded water as the primal matter, and because thal is the Phoenician word for moisture, his name may have derived from this circumstance."[11]
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According to the later historian Diogenes Laërtius, in his third century AD Lives of the Philosophers, references Herodotus, Duris, and Democritus, who all agree "that Thales was the son of Examyas and Cleobulina, and belonged to the Thelidae who are Phoenicians."[12][13] Their names are indigenous Carian and Greek, respectively.[10] Diogenes then states that "Most writers, however, represent him as a native of Miletus and of a distinguished family."[12][13] However, his supposed mother Cleobulina has also been described as his companion.[14] Diogenes then delivers more conflicting reports: one that Thales married and either fathered a son (Cybisthus or Cybisthon) or adopted his nephew of the same name; the second that he never married, telling his mother as a young man that it was too early to marry, and as an older man that it was too late. Plutarch had earlier told this version: Solon visited Thales and asked him why he remained single; Thales answered that he did not like the idea of having to worry about children. Nevertheless, several years later, anxious for family, he adopted his nephew Cybisthus.[15]
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It has been claimed that he was roughly the professional equivalent of a contemporary option trader.[16]
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It is assumed that Thales at one point in his life visited Egypt, where he learned about geometry.[17] Diogenes Laërtius wrote that Thales identifies the Milesians as Athenian colonists.[18]
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Thales (who died around 30 years before the time of Pythagoras and 300 years before Euclid, Eudoxus of Cnidus, and Eudemus of Rhodes) is often hailed as "the first Greek mathematician".[19] While some historians, such as Colin R. Fletcher, point out that there could have been a predecessor to Thales who would have been named in Eudemus' lost book History of Geometry, it is admitted that without the work "the question becomes mere speculation."[19] Fletcher holds that as there is no viable predecessor to the title of first Greek mathematician, the only question is whether Thales qualifies as a practitioner in that field; he holds that "Thales had at his command the techniques of observation, experimentation, superposition and deduction…he has proved himself mathematician."[19]
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Aristotle wrote in Metaphysics, "Thales, the founder of this school of philosophy, says the permanent entity is water (which is why he also propounded that the earth floats on water). Presumably he derived this assumption from seeing that the nutriment of everything is moist, and that heat itself is generated from moisture and depends upon it for its existence (and that from which a thing is generated is always its first principle). He derived his assumption from this; and also from the fact that the seeds of everything have a moist nature, whereas water is the first principle of the nature of moist things."[5]
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Thales involved himself in many activities, including engineering.[20] Some say that he left no writings. Others say that he wrote On the Solstice and On the Equinox. The Nautical Star-guide has been attributed to him, but this was disputed in ancient times.[21] No writing attributed to him has survived. Diogenes Laërtius quotes two letters from Thales: one to Pherecydes of Syros, offering to review his book on religion, and one to Solon, offering to keep him company on his sojourn from Athens.[clarification needed]
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A story, with different versions, recounts how Thales achieved riches from an olive harvest by prediction of the weather. In one version, he bought all the olive presses in Miletus after predicting the weather and a good harvest for a particular year. Another version of the story has Aristotle explain that Thales had reserved presses in advance, at a discount, and could rent them out at a high price when demand peaked, following his prediction of a particularly good harvest. This first version of the story would constitute the first historically known creation and use of futures, whereas the second version would be the first historically known creation and use of options.[22]
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Aristotle explains that Thales' objective in doing this was not to enrich himself but to prove to his fellow Milesians that philosophy could be useful, contrary to what they thought,[23] or alternatively, Thales had made his foray into enterprise because of a personal challenge put to him by an individual who had asked why, if Thales was an intelligent famous philosopher, he had yet to attain wealth.
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Diogenes Laërtius tells us that Thales gained fame as a counselor when he advised the Milesians not to engage in a symmachia, a "fighting together", with the Lydians. This has sometimes been interpreted as an alliance.[24][failed verification] Another story by Herodotus is that Croesus sent his army to the Persian territory. He was stopped by the river Halys, then unbridged. Thales then got the army across the river by digging a diversion upstream so as to reduce the flow, making it possible to cross the river.[25] While Herodotus reported that most of his fellow Greeks believe that Thales did divert the river Halys to assist King Croesus' military endeavors, he himself finds it doubtful.[26]
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Croesus was defeated before the city of Sardis by Cyrus, who subsequently spared Miletus because it had taken no action. Cyrus was so impressed by Croesus’ wisdom and his connection with the sages that he spared him and took his advice on various matters.[citation needed] The Ionian cities should be demoi, or "districts".
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He counselled them to establish a single seat of government, and pointed out Teos as the fittest place for it; "for that," he said, "was the centre of Ionia. Their other cities might still continue to enjoy their own laws, just as if they were independent states."[27]
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Miletus, however, received favorable terms from Cyrus. The others remained in an Ionian League of twelve cities (excluding Miletus), and were subjugated by the Persians.[citation needed]
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According to Herodotus, Thales predicted the solar eclipse of May 28, 585 BC.[7] Thales also described the position of Ursa Minor, and he thought the constellation might be useful as a guide for navigation at sea. He calculated the duration of the year and the timings of the equinoxes and solstices. He is additionally attributed with the first observation of the Hyades and with calculating the position of the Pleiades.[28] Plutarch indicates that in his day (c. AD 100) there was an extant work, the Astronomy, composed in verse and attributed to Thales.[29]
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Herodotus writes that in the sixth year of the war, the Lydians under King Alyattes and the Medes under Cyaxares were engaged in an indecisive battle when suddenly day turned into night, leading to both parties halting the fighting and negotiating a peace agreement. Herodotus also mentions that the loss of daylight had been predicted by Thales. He does not, however, mention the location of the battle.[30]
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+
Afterwards, on the refusal of Alyattes to give up his suppliants when Cyaxares sent to demand them of him, war broke out between the Lydians and the Medes, and continued for five years, with various success. In the course of it the Medes gained many victories over the Lydians, and the Lydians also gained many victories over the Medes. Among their other battles there was one night engagement. As, however, the balance had not inclined in favour of either nation, another combat took place in the sixth year, in the course of which, just as the battle was growing warm, day was on a sudden changed into night. This event had been foretold by Thales, the Milesian, who forewarned the Ionians of it, fixing for it the very year in which it actually took place. The Medes and Lydians, when they observed the change, ceased fighting, and were alike anxious to have terms of peace agreed on.[27]
|
40 |
+
|
41 |
+
However, based on the list of Medean kings and the duration of their reign reported elsewhere by Herodotus, Cyaxares died 10 years before the eclipse.[31][32]
|
42 |
+
|
43 |
+
Diogenes Laërtius[33] tells us that the Seven Sages were created in the archonship of Damasius at Athens about 582 BC and that Thales was the first sage. The same story, however, asserts that Thales emigrated to Miletus. There is also a report that he did not become a student of nature until after his political career. Much as we would like to have a date on the seven sages, we must reject these stories and the tempting date if we are to believe that Thales was a native of Miletus, predicted the eclipse, and was with Croesus in the campaign against Cyrus.
|
44 |
+
|
45 |
+
Thales received instruction from an Egyptian priest.[citation needed] It was fairly certain that he came from a wealthy, established family, in a class which customarily provided higher education for their children.[citation needed] Moreover, the ordinary citizen, unless he was a seafaring man or a merchant, could not afford the grand tour in Egypt, and did not consort with noble lawmakers such as Solon.[citation needed]
|
46 |
+
|
47 |
+
In Diogenes Laërtius' Lives of Eminent Philosophers Chapter 1.39, Laërtius relates several stories of an expensive object that is to go to the most wise. In one version (that Laërtius credits to Callimachus in his Iambics) Bathycles of Arcadia states in his will that an expensive bowl "'should be given to him who had done most good by his wisdom.' So it was given to Thales, went the round of all the sages, and came back to Thales again. And he sent it to Apollo at Didyma, with this dedication...'Thales the Milesian, son of Examyas [dedicates this] to Delphinian Apollo after twice winning the prize from all the Greeks.'"[34]
|
48 |
+
|
49 |
+
Early Greeks, and other civilizations before them, often invoked idiosyncratic explanations of natural phenomena with reference to the will of anthropomorphic gods and heroes. Instead, Thales aimed to explain natural phenomena via rational hypotheses that referenced natural processes themselves. For example, rather than assuming that earthquakes were the result of supernatural whims, Thales explained them by hypothesizing that the Earth floats on water and that earthquakes occur when the Earth is rocked by waves.[35][36]
|
50 |
+
|
51 |
+
Thales was a hylozoist (one who thinks that matter is alive,[37] i.e. containing soul(s)). Aristotle wrote (De Anima 411 a7-8) of Thales: ...Thales thought all things are full of gods. Aristotle posits the origin of Thales thought on matter generally containing souls, to Thales thinking initially on the fact of, because magnets move iron, the presence of movement of matter indicated this matter contained life.[38][39]
|
52 |
+
|
53 |
+
Thales, according to Aristotle, asked what was the nature (Greek arche) of the object so that it would behave in its characteristic way. Physis (φύσις) comes from phyein (φύειν), "to grow", related to our word "be".[40][41] (G)natura is the way a thing is "born",[42] again with the stamp of what it is in itself.
|
54 |
+
|
55 |
+
Aristotle characterizes most of the philosophers "at first" (πρῶτον) as thinking that the "principles in the form of matter were the only principles of all things", where "principle" is arche, "matter" is hyle ("wood" or "matter", "material") and "form" is eidos.[43]
|
56 |
+
|
57 |
+
Arche is translated as "principle", but the two words do not have precisely the same meaning. A principle of something is merely prior (related to pro-) to it either chronologically or logically. An arche (from ἄρχειν, "to rule") dominates an object in some way. If the arche is taken to be an origin, then specific causality is implied; that is, B is supposed to be characteristically B just because it comes from A, which dominates it.
|
58 |
+
|
59 |
+
The archai that Aristotle had in mind in his well-known passage on the first Greek scientists are not necessarily chronologically prior to their objects, but are constituents of it. For example, in pluralism objects are composed of earth, air, fire and water, but those elements do not disappear with the production of the object. They remain as archai within it, as do the atoms of the atomists.
|
60 |
+
|
61 |
+
What Aristotle is really saying is that the first philosophers were trying to define the substance(s) of which all material objects are composed. As a matter of fact, that is exactly what modern scientists are attempting to accomplish in nuclear physics, which is a second reason why Thales is described as the first western scientist,[citation needed] but some contemporary scholars reject this interpretation.[44]
|
62 |
+
|
63 |
+
Thales was known for his innovative use of geometry. His understanding was theoretical as well as practical. For example, he said:
|
64 |
+
|
65 |
+
Megiston topos: apanta gar chorei (Μέγιστον τόπος· ἄπαντα γὰρ χωρεῖ.)
|
66 |
+
|
67 |
+
The greatest is space, for it holds all things.[45]
|
68 |
+
|
69 |
+
Topos is in Newtonian-style space, since the verb, chorei, has the connotation of yielding before things, or spreading out to make room for them, which is extension. Within this extension, things have a position. Points, lines, planes and solids related by distances and angles follow from this presumption.
|
70 |
+
|
71 |
+
Thales understood similar triangles and right triangles, and what is more, used that knowledge in practical ways. The story is told in Diogenes Laërtius (loc. cit.) that he measured the height of the pyramids by their shadows at the moment when his own shadow was equal to his height. A right triangle with two equal legs is a 45-degree right triangle, all of which are similar. The length of the pyramid's shadow measured from the center of the pyramid at that moment must have been equal to its height.
|
72 |
+
|
73 |
+
This story indicates that he was familiar with the Egyptian seked, or seqed, the ratio of the run to the rise of a slope (cotangent).[citation needed] The seked is at the base of problems 56, 57, 58, 59 and 60 of the Rhind papyrus — an ancient Egyptian mathematical document.
|
74 |
+
|
75 |
+
More practically Thales used the same method to measure the distances of ships at sea, said Eudemus as reported by Proclus ("in Euclidem"). According to Kirk & Raven (reference cited below), all you need for this feat is three straight sticks pinned at one end and knowledge of your altitude. One stick goes vertically into the ground. A second is made level. With the third you sight the ship and calculate the seked from the height of the stick and its distance from the point of insertion to the line of sight (Proclus, In Euclidem, 352).
|
76 |
+
|
77 |
+
There are two theorems of Thales in elementary geometry, one known as Thales' theorem having to do with a triangle inscribed in a circle and having the circle's diameter as one leg, the other theorem being also called the intercept theorem. In addition Eudemus attributed to him the discovery that a circle is bisected by its diameter, that the base angles of an isosceles triangle are equal and that vertical angles are equal. According to a historical Note,[46] when Thales visited Egypt,[17] he observed that whenever the Egyptians drew two intersecting lines, they would measure the vertical angles to make sure that they were equal. Thales concluded that one could prove that all vertical angles are equal if one accepted some general notions such as: all straight angles are equal, equals added to equals are equal, and equals subtracted from equals are equal.
|
78 |
+
|
79 |
+
The evidence for the primacy of Thales comes to us from a book by Proclus who wrote a thousand years after Thales but is believed to have had a copy of Eudemus' book. Proclus wrote "Thales was the first to go to Egypt and bring back to Greece this study."[19] He goes on to tell us that in addition to applying the knowledge he gained in Egypt "He himself discovered many propositions and disclosed the underlying principles of many others to his successors, in some case his method being more general, in others more empirical."[19]
|
80 |
+
|
81 |
+
Other quotes from Proclus list more of Thales' mathematical achievements:
|
82 |
+
|
83 |
+
They say that Thales was the first to demonstrate that the circle is bisected by the diameter, the cause of the bisection being the unimpeded passage of the straight line through the centre.[19]
|
84 |
+
|
85 |
+
[Thales] is said to have been the first to have known and to have enunciated [the theorem] that the angles at the base of any isosceles triangle are equal, though in the more archaic manner he described the equal angles as similar.[19]
|
86 |
+
|
87 |
+
This theorem, that when two straight lines cut one another, the vertical and opposite angles are equal, was first discovered, as Eudemus says, by Thales, though the scientific demonstration was improved by the writer of Elements.[19]
|
88 |
+
|
89 |
+
Eudemus in his History of Geometry attributes this theorem [the equality of triangles having two angles and one side equal] to Thales. For he says that the method by which Thales showed how to find the distance of ships at sea necessarily involves this method.[19]
|
90 |
+
|
91 |
+
Pamphila says that, having learnt geometry from the Egyptians, he [Thales] was the first to inscribe in a circle a right-angled triangle, whereupon he sacrificed an ox.[19]
|
92 |
+
|
93 |
+
In addition to Proclus, Hieronymus of Rhodes also cites Thales as the first Greek mathematician. Hieronymus held that Thales was able to measure the height of the pyramids by using a theorem of geometry now known as the intercept theorem, (after gathering data by using his walking-stick and comparing its shadow to those cast by the pyramids). We receive variations of Hieronymus' story through Diogenes Laërtius,[47] Pliny the Elder, and Plutarch.[19][48] According to Hieronymus, historically quoted by Diogenes Laërtius, Thales found the height of pyramids by comparison between the lengths of the shadows cast by a person and by the pyramids.[49]
|
94 |
+
|
95 |
+
Due to the variations among testimonies, such as the "story of the sacrifice of an ox on the occasion of the discovery that the angle on a diameter of a circle is a right angle" in the version told by Diogenes Laërtius being accredited to Pythagoras rather than Thales, some historians (such as D. R. Dicks) question whether such anecdotes have any historical worth whatsoever.[26]
|
96 |
+
|
97 |
+
Thales' most famous philosophical position was his cosmological thesis, which comes down to us through a passage from Aristotle's Metaphysics.[50] In the work Aristotle unequivocally reported Thales' hypothesis about the nature of all matter – that the originating principle of nature was a single material substance: water. Aristotle then proceeded to proffer a number of conjectures based on his own observations to lend some credence to why Thales may have advanced this idea (though Aristotle did not hold it himself).
|
98 |
+
|
99 |
+
Aristotle laid out his own thinking about matter and form which may shed some light on the ideas of Thales, in Metaphysics 983 b6 8–11, 17–21. (The passage contains words that were later adopted by science with quite different meanings.)
|
100 |
+
|
101 |
+
That from which is everything that exists and from which it first becomes and into which it is rendered at last, its substance remaining under it, but transforming in qualities, that they say is the element and principle of things that are. …For it is necessary that there be some nature (φύσις), either one or more than one, from which become the other things of the object being saved... Thales the founder of this type of philosophy says that it is water.
|
102 |
+
|
103 |
+
In this quote we see Aristotle's depiction of the problem of change and the definition of substance. He asked if an object changes, is it the same or different? In either case how can there be a change from one to the other? The answer is that the substance "is saved", but acquires or loses different qualities (πάθη, the things you "experience").
|
104 |
+
|
105 |
+
Aristotle conjectured that Thales reached his conclusion by contemplating that the "nourishment of all things is moist and that even the hot is created from the wet and lives by it." While Aristotle's conjecture on why Thales held water as the originating principle of matter is his own thinking, his statement that Thales held it as water is generally accepted as genuinely originating with Thales and he is seen as an incipient matter-and-formist.[citation needed]
|
106 |
+
|
107 |
+
Thales thought the Earth must be a flat disk which is floating in an expanse of water.[51]
|
108 |
+
|
109 |
+
Heraclitus Homericus states that Thales drew his conclusion from seeing moist substance turn into air, slime and earth. It seems likely that Thales viewed the Earth as solidifying from the water on which it floated and the oceans that surround it.
|
110 |
+
|
111 |
+
Writing centuries later, Diogenes Laërtius also states that Thales taught "Water constituted (ὑπεστήσατο, 'stood under') the principle of all things."[52]
|
112 |
+
|
113 |
+
Aristotle considered Thales’ position to be roughly the equivalent to the later ideas of Anaximenes, who held that everything was composed of air.[53]
|
114 |
+
The 1870 book Dictionary of Greek and Roman Biography and Mythology noted:[2]
|
115 |
+
|
116 |
+
In his dogma that water is the origin of things, that is, that it is that out of which every thing arises, and into which every thing resolves itself, Thales may have followed Orphic cosmogonies, while, unlike them, he sought to establish the truth of the assertion. Hence, Aristotle, immediately after he has called him the originator of philosophy brings forward the reasons which Thales was believed to have adduced in confirmation of that assertion; for that no written development of it, or indeed any book by Thales, was extant, is proved by the expressions which Aristotle uses when he brings forward the doctrines and proofs of the Milesian. (p. 1016)
|
117 |
+
|
118 |
+
According to Aristotle, Thales thought lodestones had souls, because iron is attracted to them (by the force of magnetism).[54]
|
119 |
+
|
120 |
+
Aristotle defined the soul as the principle of life, that which imbues the matter and makes it live, giving it the animation, or power to act. The idea did not originate with him, as the Greeks in general believed in the distinction between mind and matter, which was ultimately to lead to a distinction not only between body and soul but also between matter and energy.[citation needed] If things were alive, they must have souls. This belief was no innovation, as the ordinary ancient populations of the Mediterranean did believe that natural actions were caused by divinities. Accordingly, Aristotle and other ancient writers state that Thales believed that "all things were full of gods."[55][56] In their zeal to make him the first in everything some said he was the first to hold the belief, which must have been widely known to be false.[citation needed] However, Thales was looking for something more general, a universal substance of mind.[citation needed] That also was in the polytheism of the times. Zeus was the very personification of supreme mind, dominating all the subordinate manifestations. From Thales on, however, philosophers had a tendency to depersonify or objectify mind, as though it were the substance of animation per se and not actually a god like the other gods. The end result was a total removal of mind from substance, opening the door to a non-divine principle of action.[citation needed]
|
121 |
+
|
122 |
+
Classical thought, however, had proceeded only a little way along that path. Instead of referring to the person, Zeus, they talked about the great mind:
|
123 |
+
|
124 |
+
"Thales", says Cicero,[57] "assures that water is the principle of all things; and that God is that Mind which shaped and created all things from water."
|
125 |
+
|
126 |
+
The universal mind appears as a Roman belief in Virgil as well:
|
127 |
+
|
128 |
+
In the beginning, SPIRIT within (spiritus intus) strengthens Heaven and Earth,
|
129 |
+
The watery fields, and the lucid globe of Luna, and then –
|
130 |
+
Titan stars; and mind (mens) infused through the limbs
|
131 |
+
Agitates the whole mass, and mixes itself with GREAT MATTER (magno corpore)[58]
|
132 |
+
|
133 |
+
According to Henry Fielding (1775), Diogenes Laërtius (1.35) affirmed that Thales posed "the independent pre-existence of God from all eternity, stating "that God was the oldest of all beings, for he existed without a previous cause even in the way of generation; that the world was the most beautiful of all things; for it was created by God."[59]
|
134 |
+
|
135 |
+
Due to the scarcity of sources concerning Thales and the discrepancies between the accounts given in the sources that have survived, there is a scholarly debate over possible influences on Thales and the Greek mathematicians that came after him.
|
136 |
+
Historian Roger L. Cooke points out that Proclus does not make any mention of Mesopotamian influence on Thales or Greek geometry, but "is shown clearly in Greek astronomy, in the use of sexagesimal system of measuring angles and in Ptolemy's explicit use of Mesopotamian astronomical observations."[60] Cooke notes that it may possibly also appear in the second book of Euclid's Elements, "which contains geometric constructions equivalent to certain algebraic relations that are frequently encountered in the cuneiform tablets." Cooke notes "This relation, however, is controversial."[60]
|
137 |
+
|
138 |
+
Historian B.L. Van der Waerden is among those advocating the idea of Mesopotamian influence, writing "It follows that we have to abandon the traditional belief that the oldest Greek mathematicians discovered geometry entirely by themselves…a belief that was tenable only as long as nothing was known about Babylonian mathematics. This in no way diminishes the stature of Thales; on the contrary, his genius receives only now the honour that is due to it, the honour of having developed a logical structure for geometry, of having introduced proof into geometry."[19]
|
139 |
+
|
140 |
+
Some historians, such as D. R. Dicks takes issue with the idea that we can determine from the questionable sources we have, just how influenced Thales was by Babylonian sources. He points out that while Thales is held to have been able to calculate an eclipse using a cycle called the "Saros" held to have been "borrowed from the Babylonians", "The Babylonians, however, did not use cycles to predict solar eclipses, but computed them from observations of the latitude of the moon made shortly before the expected syzygy."[26] Dicks cites historian O. Neugebauer who relates that "No Babylonian theory for predicting solar eclipse existed at 600 B.C., as one can see from the very unsatisfactory situation 400 years later; nor did the Babylonians ever develop any theory which took the influence of geographical latitude into account." Dicks examines the cycle referred to as 'Saros' – which Thales is held to have used and which is believed to stem from the Babylonians. He points out that Ptolemy makes use of this and another cycle in his book Mathematical Syntaxis but attributes it to Greek astronomers earlier than Hipparchus and not to Babylonians.[26] Dicks notes Herodotus does relate that Thales made use of a cycle to predict the eclipse, but maintains that "if so, the fulfillment of the 'prediction' was a stroke of pure luck not science".[26] He goes further joining with other historians (F. Martini, J.L. E. Dreyer, O. Neugebauer) in rejecting the historicity of the eclipse story altogether.[26] Dicks links the story of Thales discovering the cause for a solar eclipse with Herodotus' claim that Thales discovered the cycle of the sun with relation to the solstices, and concludes "he could not possibly have possessed this knowledge which neither the Egyptians nor the Babylonians nor his immediate successors possessed."[26] Josephus is the only ancient historian that claims Thales visited Babylonia.
|
141 |
+
|
142 |
+
Herodotus wrote that the Greeks learnt the practice of dividing the day into 12 parts, about the polos, and the gnomon from the Babylonians. (The exact meaning of his use of the word polos is unknown, current theories include: "the heavenly dome", "the tip of the axis of the celestial sphere", or a spherical concave sundial.) Yet even Herodotus' claims on Babylonian influence are contested by some modern historians, such as L. Zhmud, who points out that the division of the day into twelve parts (and by analogy the year) was known to the Egyptians already in the second millennium, the gnomon was known to both Egyptians and Babylonians, and the idea of the "heavenly sphere" was not used outside of Greece at this time.[61]
|
143 |
+
|
144 |
+
Less controversial than the position that Thales learnt Babylonian mathematics is the claim he was influenced by Egyptians. Pointedly historian S. N. Bychkov holds that the idea that the base angles of an isosceles triangle are equal likely came from Egypt. This is because, when building a roof for a home – having a cross section be exactly an isosceles triangle isn't crucial (as it's the ridge of the roof that must fit precisely), in contrast a symmetric square pyramid cannot have errors in the base angles of the faces or they will not fit together tightly.[60]
|
145 |
+
Historian D.R. Dicks agrees that compared to the Greeks in the era of Thales, there was a more advanced state of mathematics among the Babylonians and especially the Egyptians – "both cultures knew the correct formulae for determining the areas and volumes of simple geometrical figures such as triangles, rectangles, trapezoids, etc.; the Egyptians could also calculate correctly the volume of the frustum of a pyramid with a square base (the Babylonians used an incorrect formula for this), and used a formula for the area of a circle...which gives a value for π of 3.1605—a good approximation."[26] Dicks also agrees that this would have had an effect on Thales (whom the most ancient sources agree was interested in mathematics and astronomy) but he holds that tales of Thales' travels in these lands are pure myth.
|
146 |
+
|
147 |
+
The ancient civilization and massive monuments of Egypt had "a profound and ineradicable impression on the Greeks". They attributed to Egyptians "an immemorial knowledge of certain subjects" (including geometry) and would claim Egyptian origin for some of their own ideas to try and lend them "a respectable antiquity" (such as the "Hermetic" literature of the Alexandrian period).[26]
|
148 |
+
|
149 |
+
Dicks holds that since Thales was a prominent figure in Greek history by the time of Eudemus but "nothing certain was known except that he lived in Miletus".[26] A tradition developed that as "Milesians were in a position to be able to travel widely" Thales must have gone to Egypt.[26] As Herodotus says Egypt was the birthplace of geometry he must have learnt that while there. Since he had to have been there, surely one of the theories on Nile Flooding laid out by Herodotus must have come from Thales. Likewise as he must have been in Egypt he had to have done something with the Pyramids – thus the tale of measuring them. Similar apocryphal stories exist of Pythagoras and Plato traveling to Egypt with no corroborating evidence.
|
150 |
+
|
151 |
+
As the Egyptian and Babylonian geometry at the time was "essentially arithmetical", they used actual numbers and "the procedure is then described with explicit instructions as to what to do with these numbers" there was no mention of how the rules of procedure were made, and nothing toward a logically arranged corpus of generalized geometrical knowledge with analytical 'proofs' such as we find in the words of Euclid, Archimedes, and Apollonius."[26] So even had Thales traveled there he could not have learnt anything about the theorems he is held to have picked up there (especially because there is no evidence that any Greeks of this age could use Egyptian hieroglyphics).[26]
|
152 |
+
|
153 |
+
Likewise until around the second century BC and the time of Hipparchus (c. 190–120 BC) the Babylonian general division of the circle into 360 degrees and their sexagesimal system was unknown.[26] Herodotus says almost nothing about Babylonian literature and science, and very little about their history. Some historians, like P. Schnabel, hold that the Greeks only learned more about Babylonian culture from Berossus, a Babylonian priest who is said to have set up a school in Cos around 270 BC (but to what extent this had in the field of geometry is contested).
|
154 |
+
|
155 |
+
Dicks points out that the primitive state of Greek mathematics and astronomical ideas exhibited by the peculiar notions of Thales' successors (such as Anaximander, Anaximenes, Xenophanes, and Heraclitus), which historian J. L. Heiberg calls "a mixture of brilliant intuition and childlike analogies",[62] argues against the assertions from writers in late antiquity that Thales discovered and taught advanced concepts in these fields.
|
156 |
+
|
157 |
+
John Burnet (1892) noted[63]
|
158 |
+
|
159 |
+
Lastly, we have one admitted instance of a philosophic guild, that of the Pythagoreans. And it will be found that the hypothesis, if it is to be called by that name, of a regular organisation of scientific activity will alone explain all the facts. The development of doctrine in the hands of Thales, Anaximander, and Anaximenes, for instance, can only be understood as the elaboration of a single idea in a school with a continuous tradition.
|
160 |
+
|
161 |
+
In the long sojourn of philosophy, there has existed hardly a philosopher or historian of philosophy who did not mention Thales and try to characterize him in some way. He is generally recognized as having brought something new to human thought. Mathematics, astronomy, and medicine already existed. Thales added something to these different collections of knowledge to produce a universality, which, as far as writing tells us, was not in tradition before, but resulted in a new field.
|
162 |
+
|
163 |
+
Ever since, interested persons have been asking what that new something is. Answers fall into (at least) two categories, the theory and the method. Once an answer has been arrived at, the next logical step is to ask how Thales compares to other philosophers, which leads to his classification (rightly or wrongly).
|
164 |
+
|
165 |
+
The most natural epithets of Thales are "materialist" and "naturalist", which are based on ousia and physis. The Catholic Encyclopedia notes that Aristotle called him a physiologist, with the meaning "student of nature."[64] On the other hand, he would have qualified as an early physicist, as did Aristotle. They studied corpora, "bodies", the medieval descendants of substances.
|
166 |
+
|
167 |
+
Most agree that Thales' stamp on thought is the unity of substance, hence Bertrand Russell:[65]
|
168 |
+
|
169 |
+
The view that all matter is one is quite a reputable scientific hypothesis.
|
170 |
+
...But it is still a handsome feat to have discovered that a substance remains the same in different states of aggregation.
|
171 |
+
|
172 |
+
Russell was only reflecting an established tradition; for example: Nietzsche, in his Philosophy in the Tragic Age of the Greeks, wrote:[66]
|
173 |
+
|
174 |
+
Greek philosophy seems to begin with an absurd notion, with the proposition that water is the primal origin and the womb of all things. Is it really necessary for us to take serious notice of this proposition? It is, and for three reasons. First, because it tells us something about the primal origin of all things; second, because it does so in language devoid of image or fable, and finally, because contained in it, if only embryonically, is the thought, "all things are one."
|
175 |
+
|
176 |
+
This sort of materialism, however, should not be confused with deterministic materialism. Thales was only trying to explain the unity observed in the free play of the qualities. The arrival of uncertainty in the modern world made possible a return to Thales; for example, John Elof Boodin writes ("God and Creation"):
|
177 |
+
|
178 |
+
We cannot read the universe from the past...
|
179 |
+
|
180 |
+
Boodin defines an "emergent" materialism, in which the objects of sense emerge uncertainly from the substrate. Thales is the innovator of this sort of materialism.
|
181 |
+
|
182 |
+
Later scholastic thinkers would maintain that in his choice of water Thales was influenced by Babylonian or Chaldean religion, that held that a god had begun creation by acting upon the pre-existing water. Historian Abraham Feldman holds this does not stand up under closer examination. In Babylonian religion the water is lifeless and sterile until a god acts upon it, but for Thales water itself was divine and creative. He maintained that "All things are full of gods", and to understand the nature of things was to discover the secrets of the deities, and through this knowledge open the possibility that one could be greater than the grandest Olympian.[67]
|
183 |
+
|
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Feldman points out that while other thinkers recognized the wetness of the world "none of them was inspired to conclude that everything was ultimately aquatic."[67] He further points out that Thales was "a wealthy citizen of the fabulously rich Oriental port of Miletus...a dealer in the staples of antiquity, wine and oil...He certainly handled the shell-fish of the Phoenicians that secreted the dye of imperial purple."[67] Feldman recalls the stories of Thales measuring the distance of boats in the harbor, creating mechanical improvements for ship navigation, giving an explanation for the flooding of the Nile (vital to Egyptian agriculture and Greek trade), and changing the course of the river Halys so an army could ford it. Rather than seeing water as a barrier Thales contemplated the Ionian yearly religious gathering for athletic ritual (held on the promontory of Mycale and believed to be ordained by the ancestral kindred of Poseidon, the god of the sea). He called for the Ionian mercantile states participating in this ritual to convert it into a democratic federation under the protection of Poseidon that would hold off the forces of pastoral Persia. Feldman concludes that Thales saw "that water was a revolutionary leveler and the elemental factor determining the subsistence and business of the world"[67] and "the common channel of states."[67]
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Feldman considers Thales' environment and holds that Thales would have seen tears, sweat, and blood as granting value to a person's work and the means how life giving commodities travelled (whether on bodies of water or through the sweat of slaves and pack-animals). He would have seen that minerals could be processed from water such as life-sustaining salt and gold taken from rivers. He would’ve seen fish and other food stuffs gathered from it. Feldman points out that Thales held that the lodestone was alive as it drew metals to itself. He holds that Thales "living ever in sight of his beloved sea" would see water seem to draw all "traffic in wine and oil, milk and honey, juices and dyes" to itself, leading him to "a vision of the universe melting into a single substance that was valueless in itself and still the source of wealth."[67] Feldman concludes that for Thales "...water united all things. The social significance of water in the time of Thales induced him to discern through hardware and dry-goods, through soil and sperm, blood, sweat and tears, one fundamental fluid stuff...water, the most commonplace and powerful material known to him."[67] This combined with his contemporary's idea of "spontaneous generation" allow us to see how Thales could hold that water could be divine and creative.
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Feldman points to the lasting association of the theory that "all whatness is wetness" with Thales himself, pointing out that Diogenes Laërtius speaks of a poem, probably a satire, where Thales is snatched to heaven by the sun.[67]
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In the West, Thales represents a new kind of inquiring community as well. Edmund Husserl[68] attempts to capture the new movement as follows. Philosophical man is a "new cultural configuration" based in stepping back from "pregiven tradition" and taking up a rational "inquiry into what is true in itself;" that is, an ideal of truth. It begins with isolated individuals such as Thales, but they are supported and cooperated with as time goes on. Finally the ideal transforms the norms of society, leaping across national borders.
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The term "Pre-Socratic" derives ultimately from the philosopher Aristotle, who distinguished the early philosophers as concerning themselves with substance.
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Diogenes Laërtius on the other hand took a strictly geographic and ethnic approach. Philosophers were either Ionian or Italian. He used "Ionian" in a broader sense, including also the Athenian academics, who were not Pre-Socratics. From a philosophic point of view, any grouping at all would have been just as effective. There is no basis for an Ionian or Italian unity. Some scholars, however, concede to Diogenes' scheme as far as referring to an "Ionian" school. There was no such school in any sense.
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The most popular approach refers to a Milesian school, which is more justifiable socially and philosophically. They sought for the substance of phenomena and may have studied with each other. Some ancient writers qualify them as Milesioi, "of Miletus."
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Thales had a profound influence on other Greek thinkers and therefore on Western history. Some believe Anaximander was a pupil of Thales. Early sources report that one of Anaximander's more famous pupils, Pythagoras, visited Thales as a young man, and that Thales advised him to travel to Egypt to further his philosophical and mathematical studies.
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Many philosophers followed Thales' lead in searching for explanations in nature rather than in the supernatural; others returned to supernatural explanations, but couched them in the language of philosophy rather than of myth or of religion.
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Looking specifically at Thales' influence during the pre-Socratic era, it is clear that he stood out as one of the first thinkers who thought more in the way of logos than mythos. The difference between these two more profound ways of seeing the world is that mythos is concentrated around the stories of holy origin, while logos is concentrated around the argumentation. When the mythical man wants to explain the world the way he sees it, he explains it based on gods and powers. Mythical thought does not differentiate between things and persons[citation needed] and furthermore it does not differentiate between nature and culture[citation needed]. The way a logos thinker would present a world view is radically different from the way of the mythical thinker. In its concrete form, logos is a way of thinking not only about individualism[clarification needed], but also the abstract[clarification needed]. Furthermore, it focuses on sensible and continuous argumentation. This lays the foundation of philosophy and its way of explaining the world in terms of abstract argumentation, and not in the way of gods and mythical stories[citation needed].
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Because of Thales' elevated status in Greek culture an intense interest and admiration followed his reputation. Due to this following, the oral stories about his life were open to amplification and historical fabrication, even before they were written down generations later. Most modern dissension comes from trying to interpret what we know, in particular, distinguishing legend from fact.
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Historian D.R. Dicks and other historians divide the ancient sources about Thales into those before 320 BC and those after that year (some such as Proclus writing in the 5th century C.E. and Simplicius of Cilicia in the 6th century C.E. writing nearly a millennium after his era).[26] The first category includes Herodotus, Plato, Aristotle, Aristophanes, and Theophrastus among others. The second category includes Plautus, Aetius, Eusebius, Plutarch, Josephus, Iamblichus, Diogenes Laërtius, Theon of Smyrna, Apuleius, Clement of Alexandria, Pliny the Elder, and John Tzetzes among others.
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The earliest sources on Thales (living before 320 BC) are often the same for the other Milesian philosophers (Anaximander, and Anaximenes). These sources were either roughly contemporaneous (such as Herodotus) or lived within a few hundred years of his passing. Moreover, they were writing from an oral tradition that was widespread and well known in the Greece of their day.
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The latter sources on Thales are several "ascriptions of commentators and compilers who lived anything from 700 to 1,000 years after his death"[26] which include "anecdotes of varying degrees of plausibility"[26] and in the opinion of some historians (such as D. R. Dicks) of "no historical worth whatsoever".[26] Dicks points out that there is no agreement "among the 'authorities' even on the most fundamental facts of his life—e.g. whether he was a Milesian or a Phoenician, whether he left any writings or not, whether he was married or single-much less on the actual ideas and achievements with which he is credited."[26]
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Contrasting the work of the more ancient writers with those of the later, Dicks points out that in the works of the early writers Thales and the other men who would be hailed as "the Seven Sages of Greece" had a different reputation than that which would be assigned to them by later authors. Closer to their own era, Thales, Solon, Bias of Priene, Pittacus of Mytilene and others were hailed as "essentially practical men who played leading roles in the affairs of their respective states, and were far better known to the earlier Greeks as lawgivers and statesmen than as profound thinkers and philosophers."[26] For example, Plato praises him (coupled with Anacharsis) for being the originator of the potter's wheel and the anchor.
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Only in the writings of the second group of writers (working after 320 BC) do "we obtain the picture of Thales as the pioneer in Greek scientific thinking, particularly in regard to mathematics and astronomy which he is supposed to have learnt about in Babylonia and Egypt."[26] Rather than "the earlier tradition [where] he is a favourite example of the intelligent man who possesses some technical 'know how'...the later doxographers [such as Dicaearchus in the latter half of the fourth century BC] foist on to him any number of discoveries and achievements, in order to build him up as a figure of superhuman wisdom."[26]
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Dicks points out a further problem arises in the surviving information on Thales, for rather than using ancient sources closer to the era of Thales, the authors in later antiquity ("epitomators, excerptors, and compilers"[26]) actually "preferred to use one or more intermediaries, so that what we actually read in them comes to us not even at second, but at third or fourth or fifth hand. ...Obviously this use of intermediate sources, copied and recopied from century to century, with each writer adding additional pieces of information of greater or less plausibility from his own knowledge, provided a fertile field for errors in transmission, wrong ascriptions, and fictitious attributions".[26] Dicks points out that "certain doctrines that later commentators invented for Thales...were then accepted into the biographical tradition" being copied by subsequent writers who were then cited by those coming after them "and thus, because they may be repeated by different authors relying on different sources, may produce an illusory impression of genuineness."[26]
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Doubts even exist when considering the philosophical positions held to originate in Thales "in reality these stem directly from Aristotle's own interpretations which then became incorporated in the doxographical tradition as erroneous ascriptions to Thales".[26] (The same treatment was given by Aristotle to Anaxagoras.)
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Most philosophic analyses of the philosophy of Thales come from Aristotle, a professional philosopher, tutor of Alexander the Great, who wrote 200 years after Thales' death. Aristotle, judging from his surviving books, does not seem to have access to any works by Thales, although he probably had access to works of other authors about Thales, such as Herodotus, Hecataeus, Plato etc., as well as others whose work is now extinct. It was Aristotle's express goal to present Thales' work not because it was significant in itself, but as a prelude to his own work in natural philosophy.[69] Geoffrey Kirk and John Raven, English compilers of the fragments of the Pre-Socratics, assert that Aristotle's "judgments are often distorted by his view of earlier philosophy as a stumbling progress toward the truth that Aristotle himself revealed in his physical doctrines."[70] There was also an extensive oral tradition. Both the oral and the written were commonly read or known by all educated men in the region.
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Aristotle's philosophy had a distinct stamp: it professed the theory of matter and form, which modern scholastics have dubbed hylomorphism. Though once very widespread, it was not generally adopted by rationalist and modern science, as it mainly is useful in metaphysical analyses, but does not lend itself to the detail that is of interest to modern science. It is not clear that the theory of matter and form existed as early as Thales, and if it did, whether Thales espoused it.
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While some historians, like B. Snell, maintain that Aristotle was relying on a pre-Platonic written record by Hippias rather than oral tradition, this is a controversial position. Representing the scholarly consensus Dicks states that "the tradition about him even as early as the fifth century B.C., was evidently based entirely on hearsay....It would seem that already by Aristotle's time the early Ionians were largely names only to which popular tradition attached various ideas or achievements with greater or less plausibility".[26] He points out that works confirmed to have existed in the sixth century BC by Anaximander and Xenophanes had already disappeared by the fourth century BC, so the chances of Pre-Socratic material surviving to the age of Aristotle is almost nil (even less likely for Aristotle's pupils Theophrastus and Eudemus and less likely still for those following after them).
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The main secondary source concerning the details of Thales' life and career is Diogenes Laërtius, "Lives of Eminent Philosophers".[71] This is primarily a biographical work, as the name indicates. Compared to Aristotle, Diogenes is not much of a philosopher. He is the one who, in the Prologue to that work, is responsible for the division of the early philosophers into "Ionian" and "Italian", but he places the Academics in the Ionian school and otherwise evidences considerable disarray and contradiction, especially in the long section on forerunners of the "Ionian School". Diogenes quotes two letters attributed to Thales, but Diogenes wrote some eight centuries after Thales' death and that his sources often contained "unreliable or even fabricated information",[72] hence the concern for separating fact from legend in accounts of Thales.
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It is due to this use of hearsay and a lack of citing original sources that leads some historians, like Dicks and Werner Jaeger, to look at the late origin of the traditional picture of Pre-Socratic philosophy and view the whole idea as a construct from a later age, "the whole picture that has come down to us of the history of early philosophy was fashioned during the two or three generations from Plato to the immediate pupils of Aristotle".[73]
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Theatre or theater[a] is a collaborative form of performing art that uses live performers, typically actors or actresses, to present the experience of a real or imagined event before a live audience in a specific place, often a stage. The performers may communicate this experience to the audience through combinations of gesture, speech, song, music, and dance. Elements of art, such as painted scenery and stagecraft such as lighting are used to enhance the physicality, presence and immediacy of the experience.[1] The specific place of the performance is also named by the word "theatre" as derived from the Ancient Greek θέατρον (théatron, "a place for viewing"), itself from θεάομαι (theáomai, "to see", "to watch", "to observe").
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Modern Western theatre comes, in large measure, from the theatre of ancient Greece, from which it borrows technical terminology, classification into genres, and many of its themes, stock characters, and plot elements. Theatre artist Patrice Pavis defines theatricality, theatrical language, stage writing and the specificity of theatre as synonymous expressions that differentiate theatre from the other performing arts, literature and the arts in general.[2][b]
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Modern theatre includes performances of plays and musical theatre. The art forms of ballet and opera are also theatre and use many conventions such as acting, costumes and staging. They were influential to the development of musical theatre; see those articles for more information.
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The city-state of Athens is where western theatre originated.[3][4][5][c] It was part of a broader culture of theatricality and performance in classical Greece that included festivals, religious rituals, politics, law, athletics and gymnastics, music, poetry, weddings, funerals, and symposia.[6][5][7][8][d]
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Participation in the city-state's many festivals—and mandatory attendance at the City Dionysia as an audience member (or even as a participant in the theatrical productions) in particular—was an important part of citizenship.[10] Civic participation also involved the evaluation of the rhetoric of orators evidenced in performances in the law-court or political assembly, both of which were understood as analogous to the theatre and increasingly came to absorb its dramatic vocabulary.[11][12] The Greeks also developed the concepts of dramatic criticism and theatre architecture.[13][14][15] Actors were either amateur or at best semi-professional.[16] The theatre of ancient Greece consisted of three types of drama: tragedy, comedy, and the satyr play.[17]
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The origins of theatre in ancient Greece, according to Aristotle (384–322 BCE), the first theoretician of theatre, are to be found in the festivals that honoured Dionysus. The performances were given in semi-circular auditoria cut into hillsides, capable of seating 10,000–20,000 people. The stage consisted of a dancing floor (orchestra), dressing room and scene-building area (skene). Since the words were the most important part, good acoustics and clear delivery were paramount. The actors (always men) wore masks appropriate to the characters they represented, and each might play several parts.[18]
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Athenian tragedy—the oldest surviving form of tragedy—is a type of dance-drama that formed an important part of the theatrical culture of the city-state.[3][4][5][19][20][e] Having emerged sometime during the 6th century BCE, it flowered during the 5th century BCE (from the end of which it began to spread throughout the Greek world), and continued to be popular until the beginning of the Hellenistic period.[22][23][4][f]
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No tragedies from the 6th century BCE and only 32 of the more than a thousand that were performed in during the 5th century BCE have survived.[25][26][g] We have complete texts extant by Aeschylus, Sophocles, and Euripides.[27][h] The origins of tragedy remain obscure, though by the 5th century BCE it was institutionalised in competitions (agon) held as part of festivities celebrating Dionysus (the god of wine and fertility).[28][29] As contestants in the City Dionysia's competition (the most prestigious of the festivals to stage drama) playwrights were required to present a tetralogy of plays (though the individual works were not necessarily connected by story or theme), which usually consisted of three tragedies and one satyr play.[30][31][i] The performance of tragedies at the City Dionysia may have begun as early as 534 BCE; official records (didaskaliai) begin from 501 BCE, when the satyr play was introduced.[32][30][j]
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Most Athenian tragedies dramatise events from Greek mythology, though The Persians—which stages the Persian response to news of their military defeat at the Battle of Salamis in 480 BCE—is the notable exception in the surviving drama.[30][k] When Aeschylus won first prize for it at the City Dionysia in 472 BCE, he had been writing tragedies for more than 25 years, yet its tragic treatment of recent history is the earliest example of drama to survive.[30][34] More than 130 years later, the philosopher Aristotle analysed 5th-century Athenian tragedy in the oldest surviving work of dramatic theory—his Poetics (c. 335 BCE).
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Athenian comedy is conventionally divided into three periods, "Old Comedy", "Middle Comedy", and "New Comedy". Old Comedy survives today largely in the form of the eleven surviving plays of Aristophanes, while Middle Comedy is largely lost (preserved only in relatively short fragments in authors such as Athenaeus of Naucratis). New Comedy is known primarily from the substantial papyrus fragments of Menander. Aristotle defined comedy as a representation of laughable people that involves some kind of blunder or ugliness that does not cause pain or disaster.[l]
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In addition to the categories of comedy and tragedy at the City Dionysia, the festival also included the Satyr Play. Finding its origins in rural, agricultural rituals dedicated to Dionysus, the satyr play eventually found its way to Athens in its most well-known form. Satyr's themselves were tied to the god Dionysus as his loyal woodland companions, often engaging in drunken revelry and mischief at his side. The satyr play itself was classified as tragicomedy, erring on the side of the more modern burlesque traditions of the early twentieth century. The plotlines of the plays were typically concerned with the dealings of the pantheon of Gods and their involvement in human affairs, backed by the chorus of Satyrs. However, according to Webster, satyr actors did not always perform typical satyr actions and would break from the acting traditions assigned to the character type of a mythical forest creature.[35]
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Western theatre developed and expanded considerably under the Romans. The Roman historian Livy wrote that the Romans first experienced theatre in the 4th century BCE, with a performance by Etruscan actors.[36] Beacham argues that they had been familiar with "pre-theatrical practices" for some time before that recorded contact.[37] The theatre of ancient Rome was a thriving and diverse art form, ranging from festival performances of street theatre, nude dancing, and acrobatics, to the staging of Plautus's broadly appealing situation comedies, to the high-style, verbally elaborate tragedies of Seneca. Although Rome had a native tradition of performance, the Hellenization of Roman culture in the 3rd century BCE had a profound and energizing effect on Roman theatre and encouraged the development of Latin literature of the highest quality for the stage. The only surviving Roman tragedies, indeed the only plays of any kind from the Roman Empire, are ten dramas attributed to Lucius Annaeus Seneca (4 BCE–65 CE), the Corduba-born Stoic philosopher and tutor of Nero.[38]
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The earliest-surviving fragments of Sanskrit drama date from the 1st century CE.[39][40] The wealth of archeological evidence from earlier periods offers no indication of the existence of a tradition of theatre.[41] The ancient Vedas (hymns from between 1500 and 1000 BCE that are among the earliest examples of literature in the world) contain no hint of it (although a small number are composed in a form of dialogue) and the rituals of the Vedic period do not appear to have developed into theatre.[41] The Mahābhāṣya by Patañjali contains the earliest reference to what may have been the seeds of Sanskrit drama.[42] This treatise on grammar from 140 BCE provides a feasible date for the beginnings of theatre in India.[42]
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The major source of evidence for Sanskrit theatre is A Treatise on Theatre (Nātyaśāstra), a compendium whose date of composition is uncertain (estimates range from 200 BCE to 200 CE) and whose authorship is attributed to Bharata Muni. The Treatise is the most complete work of dramaturgy in the ancient world. It addresses acting, dance, music, dramatic construction, architecture, costuming, make-up, props, the organisation of companies, the audience, competitions, and offers a mythological account of the origin of theatre.[42] In doing so, it provides indications about the nature of actual theatrical practices. Sanskrit theatre was performed on sacred ground by priests who had been trained in the necessary skills (dance, music, and recitation) in a [hereditary process]. Its aim was both to educate and to entertain.
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Under the patronage of royal courts, performers belonged to professional companies that were directed by a stage manager (sutradhara), who may also have acted.[39][42] This task was thought of as being analogous to that of a puppeteer—the literal meaning of "sutradhara" is "holder of the strings or threads".[42] The performers were trained rigorously in vocal and physical technique.[43] There were no prohibitions against female performers; companies were all-male, all-female, and of mixed gender. Certain sentiments were considered inappropriate for men to enact, however, and were thought better suited to women. Some performers played characters their own age, while others played ages different from their own (whether younger or older). Of all the elements of theatre, the Treatise gives most attention to acting (abhinaya), which consists of two styles: realistic (lokadharmi) and conventional (natyadharmi), though the major focus is on the latter.[43][m]
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Its drama is regarded as the highest achievement of Sanskrit literature.[39] It utilised stock characters, such as the hero (nayaka), heroine (nayika), or clown (vidusaka). Actors may have specialised in a particular type. Kālidāsa in the 1st century BCE, is arguably considered to be ancient India's greatest Sanskrit dramatist. Three famous romantic plays written by Kālidāsa are the Mālavikāgnimitram (Mālavikā and Agnimitra), Vikramuurvashiiya (Pertaining to Vikrama and Urvashi), and Abhijñānaśākuntala (The Recognition of Shakuntala). The last was inspired by a story in the Mahabharata and is the most famous. It was the first to be translated into English and German. Śakuntalā (in English translation) influenced Goethe's Faust (1808–1832).[39]
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The next great Indian dramatist was Bhavabhuti (c. 7th century CE). He is said to have written the following three plays: Malati-Madhava, Mahaviracharita and Uttar Ramacharita. Among these three, the last two cover between them the entire epic of Ramayana. The powerful Indian emperor Harsha (606–648) is credited with having written three plays: the comedy Ratnavali, Priyadarsika, and the Buddhist drama Nagananda.
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The Tang dynasty is sometimes known as "The Age of 1000 Entertainments". During this era, Ming Huang formed an acting school known as The Pear Garden to produce a form of drama that was primarily musical. That is why actors are commonly called "Children of the Pear Garden." During the dynasty of Empress Ling, shadow puppetry first emerged as a recognized form of theatre in China. There were two distinct forms of shadow puppetry, Pekingese (northern) and Cantonese (southern). The two styles were differentiated by the method of making the puppets and the positioning of the rods on the puppets, as opposed to the type of play performed by the puppets. Both styles generally performed plays depicting great adventure and fantasy, rarely was this very stylized form of theatre used for political propaganda.
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Cantonese shadow puppets were the larger of the two. They were built using thick leather which created more substantial shadows. Symbolic color was also very prevalent; a black face represented honesty, a red one bravery. The rods used to control Cantonese puppets were attached perpendicular to the puppets' heads. Thus, they were not seen by the audience when the shadow was created. Pekingese puppets were more delicate and smaller. They were created out of thin, translucent leather (usually taken from the belly of a donkey). They were painted with vibrant paints, thus they cast a very colorful shadow. The thin rods which controlled their movements were attached to a leather collar at the neck of the puppet. The rods ran parallel to the bodies of the puppet then turned at a ninety degree angle to connect to the neck. While these rods were visible when the shadow was cast, they laid outside the shadow of the puppet; thus they did not interfere with the appearance of the figure. The rods attached at the necks to facilitate the use of multiple heads with one body. When the heads were not being used, they were stored in a muslin book or fabric lined box. The heads were always removed at night. This was in keeping with the old superstition that if left intact, the puppets would come to life at night. Some puppeteers went so far as to store the heads in one book and the bodies in another, to further reduce the possibility of reanimating puppets. Shadow puppetry is said to have reached its highest point of artistic development in the eleventh century before becoming a tool of the government.
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In the Song dynasty, there were many popular plays involving acrobatics and music. These developed in the Yuan dynasty into a more sophisticated form known as zaju, with a four- or five-act structure. Yuan drama spread across China and diversified into numerous regional forms, one of the best known of which is Peking Opera which is still popular today.
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Xiangsheng is a certain traditional Chinese comedic performance in the forms of monologue or dialogue.
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Theatre took on many alternate forms in the West between the 15th and 19th centuries, including commedia dell'arte and melodrama. The general trend was away from the poetic drama of the Greeks and the Renaissance and toward a more naturalistic prose style of dialogue, especially following the Industrial Revolution.[44]
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Theatre took a big pause during 1642 and 1660 in England because of the Puritan Interregnum. Viewing theatre as something sinful, the Puritans ordered the closure of London theatres in 1642. This stagnant period ended once Charles II came back to the throne in 1660 in the Restoration. Theatre (among other arts) exploded, with influence from French culture, since Charles had been exiled in France in the years previous to his reign.
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One of the big changes was the new theatre house. Instead of the type of the Elizabethan era, such as the Globe Theatre, round with no place for the actors to really prep for the next act and with no "theatre manners", the theatre house became transformed into a place of refinement, with a stage in front and stadium seating facing it. Since seating was no longer all the way around the stage, it became prioritized—some seats were obviously better than others. The king would have the best seat in the house: the very middle of the theatre, which got the widest view of the stage as well as the best way to see the point of view and vanishing point that the stage was constructed around. Philippe Jacques de Loutherbourg was one of the most influential set designers of the time because of his use of floor space and scenery.
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Because of the turmoil before this time, there was still some controversy about what should and should not be put on the stage. Jeremy Collier, a preacher, was one of the heads in this movement through his piece A Short View of the Immorality and Profaneness of the English Stage. The beliefs in this paper were mainly held by non-theatre goers and the remainder of the Puritans and very religious of the time. The main question was if seeing something immoral on stage affects behavior in the lives of those who watch it, a controversy that is still playing out today.[46]
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The seventeenth century had also introduced women to the stage, which was considered inappropriate earlier. These women were regarded as celebrities (also a newer concept, thanks to ideas on individualism that arose in the wake of Renaissance Humanism), but on the other hand, it was still very new and revolutionary that they were on the stage, and some said they were unladylike, and looked down on them. Charles II did not like young men playing the parts of young women, so he asked that women play their own parts.[47] Because women were allowed on the stage, playwrights had more leeway with plot twists, like women dressing as men, and having narrow escapes from morally sticky situations as forms of comedy.
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Comedies were full of the young and very much in vogue, with the storyline following their love lives: commonly a young roguish hero professing his love to the chaste and free minded heroine near the end of the play, much like Sheridan's The School for Scandal. Many of the comedies were fashioned after the French tradition, mainly Molière, again hailing back to the French influence brought back by the King and the Royals after their exile. Molière was one of the top comedic playwrights of the time, revolutionizing the way comedy was written and performed by combining Italian commedia dell'arte and neoclassical French comedy to create some of the longest lasting and most influential satiric comedies.[48] Tragedies were similarly victorious in their sense of righting political power, especially poignant because of the recent Restoration of the Crown.[49] They were also imitations of French tragedy, although the French had a larger distinction between comedy and tragedy, whereas the English fudged the lines occasionally and put some comedic parts in their tragedies. Common forms of non-comedic plays were sentimental comedies as well as something that would later be called tragédie bourgeoise, or domestic tragedy—that is, the tragedy of common life—were more popular in England because they appealed more to English sensibilities.[50]
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While theatre troupes were formerly often travelling, the idea of the national theatre gained support in the 18th century, inspired by Ludvig Holberg. The major promoter of the idea of the national theatre in Germany, and also of the Sturm und Drang poets, was Abel Seyler, the owner of the Hamburgische Entreprise and the Seyler Theatre Company.[51]
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Through the 19th century, the popular theatrical forms of Romanticism, melodrama, Victorian burlesque and the well-made plays of Scribe and Sardou gave way to the problem plays of Naturalism and Realism; the farces of Feydeau; Wagner's operatic Gesamtkunstwerk; musical theatre (including Gilbert and Sullivan's operas); F. C. Burnand's, W. S. Gilbert's and Oscar Wilde's drawing-room comedies; Symbolism; proto-Expressionism in the late works of August Strindberg and Henrik Ibsen;[53] and Edwardian musical comedy.
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These trends continued through the 20th century in the realism of Stanislavski and Lee Strasberg, the political theatre of Erwin Piscator and Bertolt Brecht, the so-called Theatre of the Absurd of Samuel Beckett and Eugène Ionesco, American and British musicals, the collective creations of companies of actors and directors such as Joan Littlewood's Theatre Workshop, experimental and postmodern theatre of Robert Wilson and Robert Lepage, the postcolonial theatre of August Wilson or Tomson Highway, and Augusto Boal's Theatre of the Oppressed.
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The first form of Indian theatre was the Sanskrit theatre.[54] It began after the development of Greek and Roman theatre and before the development of theatre in other parts of Asia.[54] It emerged sometime between the 2nd century BCE and the 1st century CE and flourished between the 1st century CE and the 10th, which was a period of relative peace in the history of India during which hundreds of plays were written.[55][41] Japanese forms of Kabuki, Nō, and Kyōgen developed in the 17th century CE.[56] Theatre in the medieval Islamic world included puppet theatre (which included hand puppets, shadow plays and marionette productions) and live passion plays known as ta'ziya, where actors re-enact episodes from Muslim history. In particular, Shia Islamic plays revolved around the shaheed (martyrdom) of Ali's sons Hasan ibn Ali and Husayn ibn Ali. Secular plays were known as akhraja, recorded in medieval adab literature, though they were less common than puppetry and ta'ziya theatre.[57]
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Drama is the specific mode of fiction represented in performance.[58] The term comes from a Greek word meaning "action", which is derived from the verb δράω, dráō, "to do" or "to act". The enactment of drama in theatre, performed by actors on a stage before an audience, presupposes collaborative modes of production and a collective form of reception. The structure of dramatic texts, unlike other forms of literature, is directly influenced by this collaborative production and collective reception.[59] The early modern tragedy Hamlet (1601) by Shakespeare and the classical Athenian tragedy Oedipus Rex (c. 429 BCE) by Sophocles are among the masterpieces of the art of drama.[60] A modern example is Long Day's Journey into Night by Eugene O'Neill (1956).[61]
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Considered as a genre of poetry in general, the dramatic mode has been contrasted with the epic and the lyrical modes ever since Aristotle's Poetics (c. 335 BCE)—the earliest work of dramatic theory.[n] The use of "drama" in the narrow sense to designate a specific type of play dates from the 19th century. Drama in this sense refers to a play that is neither a comedy nor a tragedy—for example, Zola's Thérèse Raquin (1873) or Chekhov's Ivanov (1887). In Ancient Greece however, the word drama encompassed all theatrical plays, tragic, comic, or anything in between.
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Drama is often combined with music and dance: the drama in opera is generally sung throughout; musicals generally include both spoken dialogue and songs; and some forms of drama have incidental music or musical accompaniment underscoring the dialogue (melodrama and Japanese Nō, for example).[o] In certain periods of history (the ancient Roman and modern Romantic) some dramas have been written to be read rather than performed.[p] In improvisation, the drama does not pre-exist the moment of performance; performers devise a dramatic script spontaneously before an audience.[q]
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Music and theatre have had a close relationship since ancient times—Athenian tragedy, for example, was a form of dance-drama that employed a chorus whose parts were sung (to the accompaniment of an aulos—an instrument comparable to the modern clarinet), as were some of the actors' responses and their 'solo songs' (monodies).[62] Modern musical theatre is a form of theatre that also combines music, spoken dialogue, and dance. It emerged from comic opera (especially Gilbert and Sullivan), variety, vaudeville, and music hall genres of the late 19th and early 20th century.[63] After the Edwardian musical comedy that began in the 1890s, the Princess Theatre musicals of the early 20th century, and comedies in the 1920s and 1930s (such as the works of Rodgers and Hammerstein), with Oklahoma! (1943), musicals moved in a more dramatic direction.[r] Famous musicals over the subsequent decades included My Fair Lady (1956), West Side Story (1957), The Fantasticks (1960), Hair (1967), A Chorus Line (1975), Les Misérables (1980), Into the Woods (1986), and The Phantom of the Opera (1986),[64] as well as more contemporary hits including Rent (1994), The Lion King (1997), Wicked (2003), and Hamilton (2015).
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Musical theatre may be produced on an intimate scale Off-Broadway, in regional theatres, and elsewhere, but it often includes spectacle. For instance, Broadway and West End musicals often include lavish costumes and sets supported by multimillion-dollar budgets.
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Theatre productions that use humour as a vehicle to tell a story qualify as comedies. This may include a modern farce such as Boeing Boeing or a classical play such as As You Like It. Theatre expressing bleak, controversial or taboo subject matter in a deliberately humorous way is referred to as black comedy. Black Comedy can have several genres like slapstick humour, dark and sarcastic comedy.
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Tragedy, then, is an imitation of an action that is serious, complete, and of a certain magnitude: in language embellished with each kind of artistic ornament, the several kinds being found in separate parts of the play; in the form of action, not of narrative; through pity and fear effecting the proper purgation of these emotions.
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Aristotle's phrase "several kinds being found in separate parts of the play" is a reference to the structural origins of drama. In it the spoken parts were written in the Attic dialect whereas the choral (recited or sung) ones in the Doric dialect, these discrepancies reflecting the differing religious origins and poetic metres of the parts that were fused into a new entity, the theatrical drama.
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Tragedy refers to a specific tradition of drama that has played a unique and important role historically in the self-definition of Western civilisation.[66][67] That tradition has been multiple and discontinuous, yet the term has often been used to invoke a powerful effect of cultural identity and historical continuity—"the Greeks and the Elizabethans, in one cultural form; Hellenes and Christians, in a common activity," as Raymond Williams puts it.[68] From its obscure origins in the theatres of Athens 2,500 years ago, from which there survives only a fraction of the work of Aeschylus, Sophocles and Euripides, through its singular articulations in the works of Shakespeare, Lope de Vega, Racine, and Schiller, to the more recent naturalistic tragedy of Strindberg, Beckett's modernist meditations on death, loss and suffering, and Müller's postmodernist reworkings of the tragic canon, tragedy has remained an important site of cultural experimentation, negotiation, struggle, and change.[69][70] In the wake of Aristotle's Poetics (335 BCE), tragedy has been used to make genre distinctions, whether at the scale of poetry in general (where the tragic divides against epic and lyric) or at the scale of the drama (where tragedy is opposed to comedy). In the modern era, tragedy has also been defined against drama, melodrama, the tragicomic, and epic theatre.[s]
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Improvisation has been a consistent feature of theatre, with the Commedia dell'arte in the sixteenth century being recognised as the first improvisation form. Popularized by Nobel Prize Winner Dario Fo and troupes such as the Upright Citizens Brigade improvisational theatre continues to evolve with many different streams and philosophies. Keith Johnstone and Viola Spolin are recognized as the first teachers of improvisation in modern times, with Johnstone exploring improvisation as an alternative to scripted theatre and Spolin and her successors exploring improvisation principally as a tool for developing dramatic work or skills or as a form for situational comedy. Spolin also became interested in how the process of learning improvisation was applicable to the development of human potential.[71] Spolin's son, Paul Sills popularized improvisational theatre as a theatrical art form when he founded, as its first director, The Second City in Chicago.
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Having been an important part of human culture for more than 2,500 years, theatre has evolved a wide range of different theories and practices. Some are related to political or spiritual ideologies, while others are based purely on "artistic" concerns. Some processes focus on a story, some on theatre as event, and some on theatre as catalyst for social change. The classical Greek philosopher Aristotle, in his seminal treatise, Poetics (c. 335 BCE) is the earliest-surviving example and its arguments have influenced theories of theatre ever since.[13][14] In it, he offers an account of what he calls "poetry" (a term which in Greek literally means "making" and in this context includes drama—comedy, tragedy, and the satyr play—as well as lyric poetry, epic poetry, and the dithyramb). He examines its "first principles" and identifies its genres and basic elements; his analysis of tragedy constitutes the core of the discussion.[72]
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Aristotle argues that tragedy consists of six qualitative parts, which are (in order of importance) mythos or "plot", ethos or "character", dianoia or "thought", lexis or "diction", melos or "song", and opsis or "spectacle".[73][74] "Although Aristotle's Poetics is universally acknowledged in the Western critical tradition", Marvin Carlson explains, "almost every detail about his seminal work has aroused divergent opinions."[75] Important theatre practitioners of the 20th century include Konstantin Stanislavski, Vsevolod Meyerhold, Jacques Copeau, Edward Gordon Craig, Bertolt Brecht, Antonin Artaud, Joan Littlewood, Peter Brook, Jerzy Grotowski, Augusto Boal, Eugenio Barba, Dario Fo, Viola Spolin, Keith Johnstone and Robert Wilson (director).
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Stanislavski treated the theatre as an art-form that is autonomous from literature and one in which the playwright's contribution should be respected as that of only one of an ensemble of creative artists.[76][77][78][79][t] His innovative contribution to modern acting theory has remained at the core of mainstream western performance training for much of the last century.[80][81][82][83][84] That many of the precepts of his system of actor training seem to be common sense and self-evident testifies to its hegemonic success.[85] Actors frequently employ his basic concepts without knowing they do so.[85] Thanks to its promotion and elaboration by acting teachers who were former students and the many translations of his theoretical writings, Stanislavski's 'system' acquired an unprecedented ability to cross cultural boundaries and developed an international reach, dominating debates about acting in Europe and the United States.[80][86][87][88] Many actors routinely equate his 'system' with the North American Method, although the latter's exclusively psychological techniques contrast sharply with Stanislavski's multivariant, holistic and psychophysical approach, which explores character and action both from the 'inside out' and the 'outside in' and treats the actor's mind and body as parts of a continuum.[89][90]
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Theatre presupposes collaborative modes of production and a collective form of reception. The structure of dramatic texts, unlike other forms of literature, is directly influenced by this collaborative production and collective reception.[59] The production of plays usually involves contributions from a playwright, director, a cast of actors, and a technical production team that includes a scenic or set designer, lighting designer, costume designer, sound designer, stage manager, production manager and technical director. Depending on the production, this team may also include a composer, dramaturg, video designer or fight director.
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Stagecraft is a generic term referring to the technical aspects of theatrical, film, and video production. It includes, but is not limited to, constructing and rigging scenery, hanging and focusing of lighting, design and procurement of costumes, makeup, procurement of props, stage management, and recording and mixing of sound. Stagecraft is distinct from the wider umbrella term of scenography. Considered a technical rather than an artistic field, it relates primarily to the practical implementation of a designer's artistic vision.
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In its most basic form, stagecraft is managed by a single person (often the stage manager of a smaller production) who arranges all scenery, costumes, lighting, and sound, and organizes the cast. At a more professional level, for example in modern Broadway houses, stagecraft is managed by hundreds of skilled carpenters, painters, electricians, stagehands, stitchers, wigmakers, and the like. This modern form of stagecraft is highly technical and specialized: it comprises many sub-disciplines and a vast trove of history and tradition. The majority of stagecraft lies between these two extremes. Regional theatres and larger community theatres will generally have a technical director and a complement of designers, each of whom has a direct hand in their respective designs.
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There are many modern theatre movements which go about producing theatre in a variety of ways. Theatrical enterprises vary enormously in sophistication and purpose. People who are involved vary from novices and hobbyists (in community theatre) to professionals (in Broadway and similar productions). Theatre can be performed with a shoestring budget or on a grand scale with multimillion-dollar budgets. This diversity manifests in the abundance of theatre sub-categories, which include:
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While most modern theatre companies rehearse one piece of theatre at a time, perform that piece for a set "run", retire the piece, and begin rehearsing a new show, repertory companies rehearse multiple shows at one time. These companies are able to perform these various pieces upon request and often perform works for years before retiring them. Most dance companies operate on this repertory system. The Royal National Theatre in London performs on a repertory system.
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Repertory theatre generally involves a group of similarly accomplished actors, and relies more on the reputation of the group than on an individual star actor. It also typically relies less on strict control by a director and less on adherence to theatrical conventions, since actors who have worked together in multiple productions can respond to each other without relying as much on convention or external direction.[91]
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In order to put on a piece of theatre, both a theatre company and a theatre venue are needed. When a theatre company is the sole company in residence at a theatre venue, this theatre (and its corresponding theatre company) are called a resident theatre or a producing theatre, because the venue produces its own work. Other theatre companies, as well as dance companies, who do not have their own theatre venue, perform at rental theatres or at presenting theatres. Both rental and presenting theatres have no full-time resident companies. They do, however, sometimes have one or more part-time resident companies, in addition to other independent partner companies who arrange to use the space when available. A rental theatre allows the independent companies to seek out the space, while a presenting theatre seeks out the independent companies to support their work by presenting them on their stage.
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Some performance groups perform in non-theatrical spaces. Such performances can take place outside or inside, in a non-traditional performance space, and include street theatre, and site-specific theatre. Non-traditional venues can be used to create more immersive or meaningful environments for audiences. They can sometimes be modified more heavily than traditional theatre venues, or can accommodate different kinds of equipment, lighting and sets.[92]
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A touring company is an independent theatre or dance company that travels, often internationally, being presented at a different theatre in each city.
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There are many theatre unions including: Actors' Equity Association (for actors and stage managers), the Stage Directors and Choreographers Society (SDC), and the International Alliance of Theatrical Stage Employees (IATSE, for designers and technicians). Many theatres require that their staff be members of these organizations.
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Theatre or theater[a] is a collaborative form of performing art that uses live performers, typically actors or actresses, to present the experience of a real or imagined event before a live audience in a specific place, often a stage. The performers may communicate this experience to the audience through combinations of gesture, speech, song, music, and dance. Elements of art, such as painted scenery and stagecraft such as lighting are used to enhance the physicality, presence and immediacy of the experience.[1] The specific place of the performance is also named by the word "theatre" as derived from the Ancient Greek θέατρον (théatron, "a place for viewing"), itself from θεάομαι (theáomai, "to see", "to watch", "to observe").
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Modern Western theatre comes, in large measure, from the theatre of ancient Greece, from which it borrows technical terminology, classification into genres, and many of its themes, stock characters, and plot elements. Theatre artist Patrice Pavis defines theatricality, theatrical language, stage writing and the specificity of theatre as synonymous expressions that differentiate theatre from the other performing arts, literature and the arts in general.[2][b]
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Modern theatre includes performances of plays and musical theatre. The art forms of ballet and opera are also theatre and use many conventions such as acting, costumes and staging. They were influential to the development of musical theatre; see those articles for more information.
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The city-state of Athens is where western theatre originated.[3][4][5][c] It was part of a broader culture of theatricality and performance in classical Greece that included festivals, religious rituals, politics, law, athletics and gymnastics, music, poetry, weddings, funerals, and symposia.[6][5][7][8][d]
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Participation in the city-state's many festivals—and mandatory attendance at the City Dionysia as an audience member (or even as a participant in the theatrical productions) in particular—was an important part of citizenship.[10] Civic participation also involved the evaluation of the rhetoric of orators evidenced in performances in the law-court or political assembly, both of which were understood as analogous to the theatre and increasingly came to absorb its dramatic vocabulary.[11][12] The Greeks also developed the concepts of dramatic criticism and theatre architecture.[13][14][15] Actors were either amateur or at best semi-professional.[16] The theatre of ancient Greece consisted of three types of drama: tragedy, comedy, and the satyr play.[17]
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The origins of theatre in ancient Greece, according to Aristotle (384–322 BCE), the first theoretician of theatre, are to be found in the festivals that honoured Dionysus. The performances were given in semi-circular auditoria cut into hillsides, capable of seating 10,000–20,000 people. The stage consisted of a dancing floor (orchestra), dressing room and scene-building area (skene). Since the words were the most important part, good acoustics and clear delivery were paramount. The actors (always men) wore masks appropriate to the characters they represented, and each might play several parts.[18]
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Athenian tragedy—the oldest surviving form of tragedy—is a type of dance-drama that formed an important part of the theatrical culture of the city-state.[3][4][5][19][20][e] Having emerged sometime during the 6th century BCE, it flowered during the 5th century BCE (from the end of which it began to spread throughout the Greek world), and continued to be popular until the beginning of the Hellenistic period.[22][23][4][f]
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No tragedies from the 6th century BCE and only 32 of the more than a thousand that were performed in during the 5th century BCE have survived.[25][26][g] We have complete texts extant by Aeschylus, Sophocles, and Euripides.[27][h] The origins of tragedy remain obscure, though by the 5th century BCE it was institutionalised in competitions (agon) held as part of festivities celebrating Dionysus (the god of wine and fertility).[28][29] As contestants in the City Dionysia's competition (the most prestigious of the festivals to stage drama) playwrights were required to present a tetralogy of plays (though the individual works were not necessarily connected by story or theme), which usually consisted of three tragedies and one satyr play.[30][31][i] The performance of tragedies at the City Dionysia may have begun as early as 534 BCE; official records (didaskaliai) begin from 501 BCE, when the satyr play was introduced.[32][30][j]
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Most Athenian tragedies dramatise events from Greek mythology, though The Persians—which stages the Persian response to news of their military defeat at the Battle of Salamis in 480 BCE—is the notable exception in the surviving drama.[30][k] When Aeschylus won first prize for it at the City Dionysia in 472 BCE, he had been writing tragedies for more than 25 years, yet its tragic treatment of recent history is the earliest example of drama to survive.[30][34] More than 130 years later, the philosopher Aristotle analysed 5th-century Athenian tragedy in the oldest surviving work of dramatic theory—his Poetics (c. 335 BCE).
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Athenian comedy is conventionally divided into three periods, "Old Comedy", "Middle Comedy", and "New Comedy". Old Comedy survives today largely in the form of the eleven surviving plays of Aristophanes, while Middle Comedy is largely lost (preserved only in relatively short fragments in authors such as Athenaeus of Naucratis). New Comedy is known primarily from the substantial papyrus fragments of Menander. Aristotle defined comedy as a representation of laughable people that involves some kind of blunder or ugliness that does not cause pain or disaster.[l]
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In addition to the categories of comedy and tragedy at the City Dionysia, the festival also included the Satyr Play. Finding its origins in rural, agricultural rituals dedicated to Dionysus, the satyr play eventually found its way to Athens in its most well-known form. Satyr's themselves were tied to the god Dionysus as his loyal woodland companions, often engaging in drunken revelry and mischief at his side. The satyr play itself was classified as tragicomedy, erring on the side of the more modern burlesque traditions of the early twentieth century. The plotlines of the plays were typically concerned with the dealings of the pantheon of Gods and their involvement in human affairs, backed by the chorus of Satyrs. However, according to Webster, satyr actors did not always perform typical satyr actions and would break from the acting traditions assigned to the character type of a mythical forest creature.[35]
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Western theatre developed and expanded considerably under the Romans. The Roman historian Livy wrote that the Romans first experienced theatre in the 4th century BCE, with a performance by Etruscan actors.[36] Beacham argues that they had been familiar with "pre-theatrical practices" for some time before that recorded contact.[37] The theatre of ancient Rome was a thriving and diverse art form, ranging from festival performances of street theatre, nude dancing, and acrobatics, to the staging of Plautus's broadly appealing situation comedies, to the high-style, verbally elaborate tragedies of Seneca. Although Rome had a native tradition of performance, the Hellenization of Roman culture in the 3rd century BCE had a profound and energizing effect on Roman theatre and encouraged the development of Latin literature of the highest quality for the stage. The only surviving Roman tragedies, indeed the only plays of any kind from the Roman Empire, are ten dramas attributed to Lucius Annaeus Seneca (4 BCE–65 CE), the Corduba-born Stoic philosopher and tutor of Nero.[38]
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The earliest-surviving fragments of Sanskrit drama date from the 1st century CE.[39][40] The wealth of archeological evidence from earlier periods offers no indication of the existence of a tradition of theatre.[41] The ancient Vedas (hymns from between 1500 and 1000 BCE that are among the earliest examples of literature in the world) contain no hint of it (although a small number are composed in a form of dialogue) and the rituals of the Vedic period do not appear to have developed into theatre.[41] The Mahābhāṣya by Patañjali contains the earliest reference to what may have been the seeds of Sanskrit drama.[42] This treatise on grammar from 140 BCE provides a feasible date for the beginnings of theatre in India.[42]
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The major source of evidence for Sanskrit theatre is A Treatise on Theatre (Nātyaśāstra), a compendium whose date of composition is uncertain (estimates range from 200 BCE to 200 CE) and whose authorship is attributed to Bharata Muni. The Treatise is the most complete work of dramaturgy in the ancient world. It addresses acting, dance, music, dramatic construction, architecture, costuming, make-up, props, the organisation of companies, the audience, competitions, and offers a mythological account of the origin of theatre.[42] In doing so, it provides indications about the nature of actual theatrical practices. Sanskrit theatre was performed on sacred ground by priests who had been trained in the necessary skills (dance, music, and recitation) in a [hereditary process]. Its aim was both to educate and to entertain.
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Under the patronage of royal courts, performers belonged to professional companies that were directed by a stage manager (sutradhara), who may also have acted.[39][42] This task was thought of as being analogous to that of a puppeteer—the literal meaning of "sutradhara" is "holder of the strings or threads".[42] The performers were trained rigorously in vocal and physical technique.[43] There were no prohibitions against female performers; companies were all-male, all-female, and of mixed gender. Certain sentiments were considered inappropriate for men to enact, however, and were thought better suited to women. Some performers played characters their own age, while others played ages different from their own (whether younger or older). Of all the elements of theatre, the Treatise gives most attention to acting (abhinaya), which consists of two styles: realistic (lokadharmi) and conventional (natyadharmi), though the major focus is on the latter.[43][m]
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Its drama is regarded as the highest achievement of Sanskrit literature.[39] It utilised stock characters, such as the hero (nayaka), heroine (nayika), or clown (vidusaka). Actors may have specialised in a particular type. Kālidāsa in the 1st century BCE, is arguably considered to be ancient India's greatest Sanskrit dramatist. Three famous romantic plays written by Kālidāsa are the Mālavikāgnimitram (Mālavikā and Agnimitra), Vikramuurvashiiya (Pertaining to Vikrama and Urvashi), and Abhijñānaśākuntala (The Recognition of Shakuntala). The last was inspired by a story in the Mahabharata and is the most famous. It was the first to be translated into English and German. Śakuntalā (in English translation) influenced Goethe's Faust (1808–1832).[39]
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The next great Indian dramatist was Bhavabhuti (c. 7th century CE). He is said to have written the following three plays: Malati-Madhava, Mahaviracharita and Uttar Ramacharita. Among these three, the last two cover between them the entire epic of Ramayana. The powerful Indian emperor Harsha (606–648) is credited with having written three plays: the comedy Ratnavali, Priyadarsika, and the Buddhist drama Nagananda.
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The Tang dynasty is sometimes known as "The Age of 1000 Entertainments". During this era, Ming Huang formed an acting school known as The Pear Garden to produce a form of drama that was primarily musical. That is why actors are commonly called "Children of the Pear Garden." During the dynasty of Empress Ling, shadow puppetry first emerged as a recognized form of theatre in China. There were two distinct forms of shadow puppetry, Pekingese (northern) and Cantonese (southern). The two styles were differentiated by the method of making the puppets and the positioning of the rods on the puppets, as opposed to the type of play performed by the puppets. Both styles generally performed plays depicting great adventure and fantasy, rarely was this very stylized form of theatre used for political propaganda.
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Cantonese shadow puppets were the larger of the two. They were built using thick leather which created more substantial shadows. Symbolic color was also very prevalent; a black face represented honesty, a red one bravery. The rods used to control Cantonese puppets were attached perpendicular to the puppets' heads. Thus, they were not seen by the audience when the shadow was created. Pekingese puppets were more delicate and smaller. They were created out of thin, translucent leather (usually taken from the belly of a donkey). They were painted with vibrant paints, thus they cast a very colorful shadow. The thin rods which controlled their movements were attached to a leather collar at the neck of the puppet. The rods ran parallel to the bodies of the puppet then turned at a ninety degree angle to connect to the neck. While these rods were visible when the shadow was cast, they laid outside the shadow of the puppet; thus they did not interfere with the appearance of the figure. The rods attached at the necks to facilitate the use of multiple heads with one body. When the heads were not being used, they were stored in a muslin book or fabric lined box. The heads were always removed at night. This was in keeping with the old superstition that if left intact, the puppets would come to life at night. Some puppeteers went so far as to store the heads in one book and the bodies in another, to further reduce the possibility of reanimating puppets. Shadow puppetry is said to have reached its highest point of artistic development in the eleventh century before becoming a tool of the government.
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In the Song dynasty, there were many popular plays involving acrobatics and music. These developed in the Yuan dynasty into a more sophisticated form known as zaju, with a four- or five-act structure. Yuan drama spread across China and diversified into numerous regional forms, one of the best known of which is Peking Opera which is still popular today.
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Xiangsheng is a certain traditional Chinese comedic performance in the forms of monologue or dialogue.
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Theatre took on many alternate forms in the West between the 15th and 19th centuries, including commedia dell'arte and melodrama. The general trend was away from the poetic drama of the Greeks and the Renaissance and toward a more naturalistic prose style of dialogue, especially following the Industrial Revolution.[44]
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Theatre took a big pause during 1642 and 1660 in England because of the Puritan Interregnum. Viewing theatre as something sinful, the Puritans ordered the closure of London theatres in 1642. This stagnant period ended once Charles II came back to the throne in 1660 in the Restoration. Theatre (among other arts) exploded, with influence from French culture, since Charles had been exiled in France in the years previous to his reign.
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One of the big changes was the new theatre house. Instead of the type of the Elizabethan era, such as the Globe Theatre, round with no place for the actors to really prep for the next act and with no "theatre manners", the theatre house became transformed into a place of refinement, with a stage in front and stadium seating facing it. Since seating was no longer all the way around the stage, it became prioritized—some seats were obviously better than others. The king would have the best seat in the house: the very middle of the theatre, which got the widest view of the stage as well as the best way to see the point of view and vanishing point that the stage was constructed around. Philippe Jacques de Loutherbourg was one of the most influential set designers of the time because of his use of floor space and scenery.
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Because of the turmoil before this time, there was still some controversy about what should and should not be put on the stage. Jeremy Collier, a preacher, was one of the heads in this movement through his piece A Short View of the Immorality and Profaneness of the English Stage. The beliefs in this paper were mainly held by non-theatre goers and the remainder of the Puritans and very religious of the time. The main question was if seeing something immoral on stage affects behavior in the lives of those who watch it, a controversy that is still playing out today.[46]
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The seventeenth century had also introduced women to the stage, which was considered inappropriate earlier. These women were regarded as celebrities (also a newer concept, thanks to ideas on individualism that arose in the wake of Renaissance Humanism), but on the other hand, it was still very new and revolutionary that they were on the stage, and some said they were unladylike, and looked down on them. Charles II did not like young men playing the parts of young women, so he asked that women play their own parts.[47] Because women were allowed on the stage, playwrights had more leeway with plot twists, like women dressing as men, and having narrow escapes from morally sticky situations as forms of comedy.
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Comedies were full of the young and very much in vogue, with the storyline following their love lives: commonly a young roguish hero professing his love to the chaste and free minded heroine near the end of the play, much like Sheridan's The School for Scandal. Many of the comedies were fashioned after the French tradition, mainly Molière, again hailing back to the French influence brought back by the King and the Royals after their exile. Molière was one of the top comedic playwrights of the time, revolutionizing the way comedy was written and performed by combining Italian commedia dell'arte and neoclassical French comedy to create some of the longest lasting and most influential satiric comedies.[48] Tragedies were similarly victorious in their sense of righting political power, especially poignant because of the recent Restoration of the Crown.[49] They were also imitations of French tragedy, although the French had a larger distinction between comedy and tragedy, whereas the English fudged the lines occasionally and put some comedic parts in their tragedies. Common forms of non-comedic plays were sentimental comedies as well as something that would later be called tragédie bourgeoise, or domestic tragedy—that is, the tragedy of common life—were more popular in England because they appealed more to English sensibilities.[50]
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While theatre troupes were formerly often travelling, the idea of the national theatre gained support in the 18th century, inspired by Ludvig Holberg. The major promoter of the idea of the national theatre in Germany, and also of the Sturm und Drang poets, was Abel Seyler, the owner of the Hamburgische Entreprise and the Seyler Theatre Company.[51]
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Through the 19th century, the popular theatrical forms of Romanticism, melodrama, Victorian burlesque and the well-made plays of Scribe and Sardou gave way to the problem plays of Naturalism and Realism; the farces of Feydeau; Wagner's operatic Gesamtkunstwerk; musical theatre (including Gilbert and Sullivan's operas); F. C. Burnand's, W. S. Gilbert's and Oscar Wilde's drawing-room comedies; Symbolism; proto-Expressionism in the late works of August Strindberg and Henrik Ibsen;[53] and Edwardian musical comedy.
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These trends continued through the 20th century in the realism of Stanislavski and Lee Strasberg, the political theatre of Erwin Piscator and Bertolt Brecht, the so-called Theatre of the Absurd of Samuel Beckett and Eugène Ionesco, American and British musicals, the collective creations of companies of actors and directors such as Joan Littlewood's Theatre Workshop, experimental and postmodern theatre of Robert Wilson and Robert Lepage, the postcolonial theatre of August Wilson or Tomson Highway, and Augusto Boal's Theatre of the Oppressed.
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The first form of Indian theatre was the Sanskrit theatre.[54] It began after the development of Greek and Roman theatre and before the development of theatre in other parts of Asia.[54] It emerged sometime between the 2nd century BCE and the 1st century CE and flourished between the 1st century CE and the 10th, which was a period of relative peace in the history of India during which hundreds of plays were written.[55][41] Japanese forms of Kabuki, Nō, and Kyōgen developed in the 17th century CE.[56] Theatre in the medieval Islamic world included puppet theatre (which included hand puppets, shadow plays and marionette productions) and live passion plays known as ta'ziya, where actors re-enact episodes from Muslim history. In particular, Shia Islamic plays revolved around the shaheed (martyrdom) of Ali's sons Hasan ibn Ali and Husayn ibn Ali. Secular plays were known as akhraja, recorded in medieval adab literature, though they were less common than puppetry and ta'ziya theatre.[57]
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Drama is the specific mode of fiction represented in performance.[58] The term comes from a Greek word meaning "action", which is derived from the verb δράω, dráō, "to do" or "to act". The enactment of drama in theatre, performed by actors on a stage before an audience, presupposes collaborative modes of production and a collective form of reception. The structure of dramatic texts, unlike other forms of literature, is directly influenced by this collaborative production and collective reception.[59] The early modern tragedy Hamlet (1601) by Shakespeare and the classical Athenian tragedy Oedipus Rex (c. 429 BCE) by Sophocles are among the masterpieces of the art of drama.[60] A modern example is Long Day's Journey into Night by Eugene O'Neill (1956).[61]
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Considered as a genre of poetry in general, the dramatic mode has been contrasted with the epic and the lyrical modes ever since Aristotle's Poetics (c. 335 BCE)—the earliest work of dramatic theory.[n] The use of "drama" in the narrow sense to designate a specific type of play dates from the 19th century. Drama in this sense refers to a play that is neither a comedy nor a tragedy—for example, Zola's Thérèse Raquin (1873) or Chekhov's Ivanov (1887). In Ancient Greece however, the word drama encompassed all theatrical plays, tragic, comic, or anything in between.
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Drama is often combined with music and dance: the drama in opera is generally sung throughout; musicals generally include both spoken dialogue and songs; and some forms of drama have incidental music or musical accompaniment underscoring the dialogue (melodrama and Japanese Nō, for example).[o] In certain periods of history (the ancient Roman and modern Romantic) some dramas have been written to be read rather than performed.[p] In improvisation, the drama does not pre-exist the moment of performance; performers devise a dramatic script spontaneously before an audience.[q]
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Music and theatre have had a close relationship since ancient times—Athenian tragedy, for example, was a form of dance-drama that employed a chorus whose parts were sung (to the accompaniment of an aulos—an instrument comparable to the modern clarinet), as were some of the actors' responses and their 'solo songs' (monodies).[62] Modern musical theatre is a form of theatre that also combines music, spoken dialogue, and dance. It emerged from comic opera (especially Gilbert and Sullivan), variety, vaudeville, and music hall genres of the late 19th and early 20th century.[63] After the Edwardian musical comedy that began in the 1890s, the Princess Theatre musicals of the early 20th century, and comedies in the 1920s and 1930s (such as the works of Rodgers and Hammerstein), with Oklahoma! (1943), musicals moved in a more dramatic direction.[r] Famous musicals over the subsequent decades included My Fair Lady (1956), West Side Story (1957), The Fantasticks (1960), Hair (1967), A Chorus Line (1975), Les Misérables (1980), Into the Woods (1986), and The Phantom of the Opera (1986),[64] as well as more contemporary hits including Rent (1994), The Lion King (1997), Wicked (2003), and Hamilton (2015).
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Musical theatre may be produced on an intimate scale Off-Broadway, in regional theatres, and elsewhere, but it often includes spectacle. For instance, Broadway and West End musicals often include lavish costumes and sets supported by multimillion-dollar budgets.
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Theatre productions that use humour as a vehicle to tell a story qualify as comedies. This may include a modern farce such as Boeing Boeing or a classical play such as As You Like It. Theatre expressing bleak, controversial or taboo subject matter in a deliberately humorous way is referred to as black comedy. Black Comedy can have several genres like slapstick humour, dark and sarcastic comedy.
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Tragedy, then, is an imitation of an action that is serious, complete, and of a certain magnitude: in language embellished with each kind of artistic ornament, the several kinds being found in separate parts of the play; in the form of action, not of narrative; through pity and fear effecting the proper purgation of these emotions.
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Aristotle's phrase "several kinds being found in separate parts of the play" is a reference to the structural origins of drama. In it the spoken parts were written in the Attic dialect whereas the choral (recited or sung) ones in the Doric dialect, these discrepancies reflecting the differing religious origins and poetic metres of the parts that were fused into a new entity, the theatrical drama.
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Tragedy refers to a specific tradition of drama that has played a unique and important role historically in the self-definition of Western civilisation.[66][67] That tradition has been multiple and discontinuous, yet the term has often been used to invoke a powerful effect of cultural identity and historical continuity—"the Greeks and the Elizabethans, in one cultural form; Hellenes and Christians, in a common activity," as Raymond Williams puts it.[68] From its obscure origins in the theatres of Athens 2,500 years ago, from which there survives only a fraction of the work of Aeschylus, Sophocles and Euripides, through its singular articulations in the works of Shakespeare, Lope de Vega, Racine, and Schiller, to the more recent naturalistic tragedy of Strindberg, Beckett's modernist meditations on death, loss and suffering, and Müller's postmodernist reworkings of the tragic canon, tragedy has remained an important site of cultural experimentation, negotiation, struggle, and change.[69][70] In the wake of Aristotle's Poetics (335 BCE), tragedy has been used to make genre distinctions, whether at the scale of poetry in general (where the tragic divides against epic and lyric) or at the scale of the drama (where tragedy is opposed to comedy). In the modern era, tragedy has also been defined against drama, melodrama, the tragicomic, and epic theatre.[s]
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Improvisation has been a consistent feature of theatre, with the Commedia dell'arte in the sixteenth century being recognised as the first improvisation form. Popularized by Nobel Prize Winner Dario Fo and troupes such as the Upright Citizens Brigade improvisational theatre continues to evolve with many different streams and philosophies. Keith Johnstone and Viola Spolin are recognized as the first teachers of improvisation in modern times, with Johnstone exploring improvisation as an alternative to scripted theatre and Spolin and her successors exploring improvisation principally as a tool for developing dramatic work or skills or as a form for situational comedy. Spolin also became interested in how the process of learning improvisation was applicable to the development of human potential.[71] Spolin's son, Paul Sills popularized improvisational theatre as a theatrical art form when he founded, as its first director, The Second City in Chicago.
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Having been an important part of human culture for more than 2,500 years, theatre has evolved a wide range of different theories and practices. Some are related to political or spiritual ideologies, while others are based purely on "artistic" concerns. Some processes focus on a story, some on theatre as event, and some on theatre as catalyst for social change. The classical Greek philosopher Aristotle, in his seminal treatise, Poetics (c. 335 BCE) is the earliest-surviving example and its arguments have influenced theories of theatre ever since.[13][14] In it, he offers an account of what he calls "poetry" (a term which in Greek literally means "making" and in this context includes drama—comedy, tragedy, and the satyr play—as well as lyric poetry, epic poetry, and the dithyramb). He examines its "first principles" and identifies its genres and basic elements; his analysis of tragedy constitutes the core of the discussion.[72]
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Aristotle argues that tragedy consists of six qualitative parts, which are (in order of importance) mythos or "plot", ethos or "character", dianoia or "thought", lexis or "diction", melos or "song", and opsis or "spectacle".[73][74] "Although Aristotle's Poetics is universally acknowledged in the Western critical tradition", Marvin Carlson explains, "almost every detail about his seminal work has aroused divergent opinions."[75] Important theatre practitioners of the 20th century include Konstantin Stanislavski, Vsevolod Meyerhold, Jacques Copeau, Edward Gordon Craig, Bertolt Brecht, Antonin Artaud, Joan Littlewood, Peter Brook, Jerzy Grotowski, Augusto Boal, Eugenio Barba, Dario Fo, Viola Spolin, Keith Johnstone and Robert Wilson (director).
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Stanislavski treated the theatre as an art-form that is autonomous from literature and one in which the playwright's contribution should be respected as that of only one of an ensemble of creative artists.[76][77][78][79][t] His innovative contribution to modern acting theory has remained at the core of mainstream western performance training for much of the last century.[80][81][82][83][84] That many of the precepts of his system of actor training seem to be common sense and self-evident testifies to its hegemonic success.[85] Actors frequently employ his basic concepts without knowing they do so.[85] Thanks to its promotion and elaboration by acting teachers who were former students and the many translations of his theoretical writings, Stanislavski's 'system' acquired an unprecedented ability to cross cultural boundaries and developed an international reach, dominating debates about acting in Europe and the United States.[80][86][87][88] Many actors routinely equate his 'system' with the North American Method, although the latter's exclusively psychological techniques contrast sharply with Stanislavski's multivariant, holistic and psychophysical approach, which explores character and action both from the 'inside out' and the 'outside in' and treats the actor's mind and body as parts of a continuum.[89][90]
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Theatre presupposes collaborative modes of production and a collective form of reception. The structure of dramatic texts, unlike other forms of literature, is directly influenced by this collaborative production and collective reception.[59] The production of plays usually involves contributions from a playwright, director, a cast of actors, and a technical production team that includes a scenic or set designer, lighting designer, costume designer, sound designer, stage manager, production manager and technical director. Depending on the production, this team may also include a composer, dramaturg, video designer or fight director.
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Stagecraft is a generic term referring to the technical aspects of theatrical, film, and video production. It includes, but is not limited to, constructing and rigging scenery, hanging and focusing of lighting, design and procurement of costumes, makeup, procurement of props, stage management, and recording and mixing of sound. Stagecraft is distinct from the wider umbrella term of scenography. Considered a technical rather than an artistic field, it relates primarily to the practical implementation of a designer's artistic vision.
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In its most basic form, stagecraft is managed by a single person (often the stage manager of a smaller production) who arranges all scenery, costumes, lighting, and sound, and organizes the cast. At a more professional level, for example in modern Broadway houses, stagecraft is managed by hundreds of skilled carpenters, painters, electricians, stagehands, stitchers, wigmakers, and the like. This modern form of stagecraft is highly technical and specialized: it comprises many sub-disciplines and a vast trove of history and tradition. The majority of stagecraft lies between these two extremes. Regional theatres and larger community theatres will generally have a technical director and a complement of designers, each of whom has a direct hand in their respective designs.
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There are many modern theatre movements which go about producing theatre in a variety of ways. Theatrical enterprises vary enormously in sophistication and purpose. People who are involved vary from novices and hobbyists (in community theatre) to professionals (in Broadway and similar productions). Theatre can be performed with a shoestring budget or on a grand scale with multimillion-dollar budgets. This diversity manifests in the abundance of theatre sub-categories, which include:
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While most modern theatre companies rehearse one piece of theatre at a time, perform that piece for a set "run", retire the piece, and begin rehearsing a new show, repertory companies rehearse multiple shows at one time. These companies are able to perform these various pieces upon request and often perform works for years before retiring them. Most dance companies operate on this repertory system. The Royal National Theatre in London performs on a repertory system.
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Repertory theatre generally involves a group of similarly accomplished actors, and relies more on the reputation of the group than on an individual star actor. It also typically relies less on strict control by a director and less on adherence to theatrical conventions, since actors who have worked together in multiple productions can respond to each other without relying as much on convention or external direction.[91]
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In order to put on a piece of theatre, both a theatre company and a theatre venue are needed. When a theatre company is the sole company in residence at a theatre venue, this theatre (and its corresponding theatre company) are called a resident theatre or a producing theatre, because the venue produces its own work. Other theatre companies, as well as dance companies, who do not have their own theatre venue, perform at rental theatres or at presenting theatres. Both rental and presenting theatres have no full-time resident companies. They do, however, sometimes have one or more part-time resident companies, in addition to other independent partner companies who arrange to use the space when available. A rental theatre allows the independent companies to seek out the space, while a presenting theatre seeks out the independent companies to support their work by presenting them on their stage.
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Some performance groups perform in non-theatrical spaces. Such performances can take place outside or inside, in a non-traditional performance space, and include street theatre, and site-specific theatre. Non-traditional venues can be used to create more immersive or meaningful environments for audiences. They can sometimes be modified more heavily than traditional theatre venues, or can accommodate different kinds of equipment, lighting and sets.[92]
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A touring company is an independent theatre or dance company that travels, often internationally, being presented at a different theatre in each city.
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There are many theatre unions including: Actors' Equity Association (for actors and stage managers), the Stage Directors and Choreographers Society (SDC), and the International Alliance of Theatrical Stage Employees (IATSE, for designers and technicians). Many theatres require that their staff be members of these organizations.
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The Beatles were an English rock band formed in Liverpool in 1960. The group, whose best-known line-up comprised John Lennon, Paul McCartney, George Harrison and Ringo Starr, are regarded as the most influential band of all time.[1] They were integral to the development of 1960s counterculture and popular music's recognition as an art form.[2] Rooted in skiffle, beat and 1950s rock and roll, their sound incorporated elements of classical music and traditional pop in innovative ways; the band later explored music styles ranging from ballads and Indian music to psychedelia and hard rock. As pioneers in recording, songwriting and artistic presentation, the Beatles revolutionised many aspects of the music industry and were often publicised as leaders of the era's youth and sociocultural movements.[3]
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Led by primary songwriters Lennon and McCartney, the Beatles built their reputation playing clubs in Liverpool and Hamburg over three years from 1960, initially with Stuart Sutcliffe playing bass. The core trio of Lennon, McCartney and Harrison, together since 1958, went through a succession of drummers, including Pete Best, before asking Starr to join them in 1962. Manager Brian Epstein moulded them into a professional act, and producer George Martin guided and developed their recordings, greatly expanding their domestic success after their first hit, "Love Me Do", in late 1962. As their popularity grew into the intense fan frenzy dubbed "Beatlemania", the band acquired the nickname "the Fab Four", with Epstein, Martin and other members of the band's entourage sometimes given the informal title of "fifth Beatle".
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By early 1964, the Beatles were international stars, leading the "British Invasion" of the United States pop market and breaking numerous sales records. They soon made their film debut with A Hard Day's Night (1964). From 1965 onwards, they produced records of greater complexity, including the albums Rubber Soul (1965), Revolver (1966) and Sgt. Pepper's Lonely Hearts Club Band (1967), and enjoyed further commercial success with The Beatles (also known as "the White Album", 1968) and Abbey Road (1969). In 1968, they founded Apple Corps, a multi-armed multimedia corporation that continues to oversee projects related to the band's legacy. After the group's break-up in 1970, all four members enjoyed success as solo artists. Lennon was shot and killed in December 1980, and Harrison died of lung cancer in November 2001. McCartney and Starr remain musically active.
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The Beatles are the best-selling music act of all time, with estimated sales of 600 million units worldwide.[4] They are the best-selling act in the US, with certified sales of 183 million units. They hold the record for most number-one albums on the UK Albums Chart, most number-one hits on the Billboard Hot 100 chart, and most singles sold in the UK. The group were inducted into the Rock and Roll Hall of Fame in 1988, and all four main members were inducted individually between 1994 and 2015. In 2008, the group topped Billboard's list of the all-time most successful artists on the Billboard Hot 100. The band received seven Grammy Awards, four Brit Awards, an Academy Award (for Best Original Song Score for the 1970 film Let It Be) and fifteen Ivor Novello Awards. Time magazine named them among the 20th century's 100 most important people.
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In March 1957, John Lennon, then aged sixteen, formed a skiffle group with several friends from Quarry Bank High School in Liverpool. They briefly called themselves the Blackjacks, before changing their name to the Quarrymen after discovering that another local group were already using the name.[5] Fifteen-year-old Paul McCartney joined them as a rhythm guitarist shortly after he and Lennon met that July.[6] In February 1958, McCartney invited his friend George Harrison to watch the band. The fifteen-year-old auditioned for Lennon, impressing him with his playing, but Lennon initially thought Harrison was too young for the band. After a month of Harrison's persistence, during a second meeting (arranged by McCartney), he performed the lead guitar part of the instrumental song "Raunchy" on the upper deck of a Liverpool bus,[7] and they enlisted him as their lead guitarist.[8][9]
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By January 1959, Lennon's Quarry Bank friends had left the group, and he began his studies at the Liverpool College of Art.[10] The three guitarists, billing themselves as Johnny and the Moondogs,[11] were playing rock and roll whenever they could find a drummer.[12] Lennon's art school friend Stuart Sutcliffe, who had just sold one of his paintings and was persuaded to purchase a bass guitar with the proceeds, joined in January 1960, and it was he who suggested changing the band's name to Beatals, as a tribute to Buddy Holly and the Crickets.[13][14] They used this name until May, when they became the Silver Beetles, before undertaking a brief tour of Scotland as the backing group for pop singer and fellow Liverpudlian Johnny Gentle. By early July, they had refashioned themselves as the Silver Beatles, and by the middle of August shortened the name to the Beatles.[15]
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Allan Williams, the Beatles' unofficial manager, arranged a residency for them in Hamburg, and for this they auditioned and hired drummer Pete Best in mid-August 1960. The band, now a five-piece, departed Liverpool for Hamburg four days later, contracted to club owner Bruno Koschmider for what would be a 31⁄2-month residency.[16] Beatles historian Mark Lewisohn writes: "They pulled into Hamburg at dusk on 17 August, the time when the red-light area comes to life ... flashing neon lights screamed out the various entertainment on offer, while scantily clad women sat unabashed in shop windows waiting for business opportunities."[17]
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Koschmider had converted a couple of strip clubs in the district into music venues, and he initially placed the Beatles at the Indra Club. After closing Indra due to noise complaints, he moved them to the Kaiserkeller in October.[18] When he learned they had been performing at the rival Top Ten Club in breach of their contract, he gave the band one month's termination notice,[19] and reported the underage Harrison, who had obtained permission to stay in Hamburg by lying to the German authorities about his age.[20] The authorities arranged for Harrison's deportation in late November.[21] One week later, Koschmider had McCartney and Best arrested for arson after they set fire to a condom in a concrete corridor; the authorities deported them.[22] Lennon returned to Liverpool in early December, while Sutcliffe remained in Hamburg until late February with his German fiancée Astrid Kirchherr,[23] who took the first semi-professional photos of the Beatles.[24]
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During the next two years, the Beatles were resident for periods in Hamburg, where they used Preludin both recreationally and to maintain their energy through all-night performances.[25] In 1961, during their second Hamburg engagement, Kirchherr cut Sutcliffe's hair in the "exi" (existentialist) style, later adopted by the other Beatles.[26][27] When Sutcliffe decided to leave the band early that year and resume his art studies in Germany, McCartney took up the bass.[28] Producer Bert Kaempfert contracted what was now a four-piece group until June 1962, and he used them as Tony Sheridan's backing band on a series of recordings for Polydor Records.[14][29] As part of the sessions, the Beatles were signed to Polydor for one year.[30] Credited to "Tony Sheridan & the Beat Brothers", the single "My Bonnie", recorded in June 1961 and released four months later, reached number 32 on the Musikmarkt chart.[31]
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After the Beatles completed their second Hamburg residency, they enjoyed increasing popularity in Liverpool with the growing Merseybeat movement. However, they were also growing tired of the monotony of numerous appearances at the same clubs night after night.[32] In November 1961, during one of the group's frequent performances at The Cavern Club, they encountered Brian Epstein, a local record-store owner and music columnist.[33] He later recalled: "I immediately liked what I heard. They were fresh, and they were honest, and they had what I thought was a sort of presence ... [a] star quality."[34]
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Epstein courted the band over the next couple of months, and they appointed him as their manager in January 1962.[35] Throughout early and mid-1962, Epstein sought to free the Beatles from their contractual obligations to Bert Kaempfert Productions. He eventually negotiated a one-month-early release from their contract in exchange for one last recording session in Hamburg.[36] Tragedy greeted them on their return to Germany in April, when a distraught Kirchherr met them at the airport with news of Sutcliffe's death the previous day from what was later determined as a brain haemorrhage.[37]
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Epstein began negotiations with record labels for a recording contract. To secure a UK record contract, Epstein negotiated an early end to the band's contract with Polydor, in exchange for more recordings backing Tony Sheridan.[38] After a New Year's Day audition, Decca Records rejected the band with the comment "Guitar groups are on the way out, Mr. Epstein."[39] However, three months later, producer George Martin signed the Beatles to EMI's Parlophone label.[37]
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Martin's first recording session with the Beatles took place at EMI's Abbey Road Studios in London on 6 June 1962.[40] Martin immediately complained to Epstein about Best's poor drumming and suggested they use a session drummer in his place.[41] Already contemplating Best's dismissal,[42] the Beatles replaced him in mid-August with Ringo Starr, who left Rory Storm and the Hurricanes to join them.[40] A 4 September session at EMI yielded a recording of "Love Me Do" featuring Starr on drums, but a dissatisfied Martin hired drummer Andy White for the band's third session a week later, which produced recordings of "Love Me Do", "Please Please Me" and "P.S. I Love You".[40]
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Martin initially selected the Starr version of "Love Me Do" for the band's first single, though subsequent re-pressings featured the White version, with Starr on tambourine.[40] Released in early October, "Love Me Do" peaked at number seventeen on the Record Retailer chart.[43] Their television debut came later that month with a live performance on the regional news programme People and Places.[44] After Martin suggested rerecording "Please Please Me" at a faster tempo,[45] a studio session in late November yielded that recording,[46] of which Martin accurately predicted, "You've just made your first No. 1."[47]
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In December 1962, the Beatles concluded their fifth and final Hamburg residency.[48] By 1963, they had agreed that all four band members would contribute vocals to their albums – including Starr, despite his restricted vocal range, to validate his standing in the group.[49] Lennon and McCartney had established a songwriting partnership, and as the band's success grew, their dominant collaboration limited Harrison's opportunities as a lead vocalist.[50] Epstein, to maximise the Beatles' commercial potential, encouraged them to adopt a professional approach to performing.[51] Lennon recalled him saying, "Look, if you really want to get in these bigger places, you're going to have to change – stop eating on stage, stop swearing, stop smoking ...."[39] Lennon said: "We used to dress how we liked, on and off stage. He'd tell us that jeans were not particularly smart and could we possibly manage to wear proper trousers, but he didn't want us suddenly looking square. He'd let us have our own sense of individuality."[39]
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On 11 February 1963, the Beatles recorded ten songs during a single studio session for their debut LP, Please Please Me. The album was supplemented by the four tracks already released on their first two singles. Martin originally considered recording the Beatles' debut LP live at The Cavern Club, but after deciding that the building's acoustics were inadequate, he elected to simulate a "live" album with minimal production in "a single marathon session at Abbey Road".[53] After the moderate success of "Love Me Do", the single "Please Please Me" met with a more emphatic reception. Released in January 1963, two months ahead of the album of the same name, the song reached number one on every UK chart except Record Retailer, where it peaked at number two.[54]
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Recalling how the Beatles "rushed to deliver a debut album, bashing out Please Please Me in a day", AllMusic's Stephen Thomas Erlewine comments, "Decades after its release, the album still sounds fresh, precisely because of its intense origins."[55] Lennon said little thought went into composition at the time; he and McCartney were "just writing songs à la Everly Brothers, à la Buddy Holly, pop songs with no more thought of them than that – to create a sound. And the words were almost irrelevant."[56]
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Released in March 1963, the album initiated a run during which eleven of their twelve studio albums released in the United Kingdom through to 1970 reached number one.[59] The band's third single, "From Me to You", came out in April and was also a chart-topping hit, starting an almost unbroken string of seventeen British number-one singles for the Beatles, including all but one of the eighteen they released over the next six years.[60] Issued in August, the band's fourth single, "She Loves You", achieved the fastest sales of any record in the UK up to that time, selling three-quarters of a million copies in under four weeks.[61] It became their first single to sell a million copies, and remained the biggest-selling record in the UK until 1978.[62][nb 1]
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Their commercial success brought increased media exposure, to which the Beatles responded with an irreverent and comical attitude that defied the expectations of pop musicians at the time, inspiring even more interest.[63] The band toured the UK three times in the first half of the year: a four-week tour that began in February, the Beatles' first nationwide, preceded three-week tours in March and May–June.[64] As their popularity spread, a frenzied adulation of the group took hold. Greeted with riotous enthusiasm by screaming fans, the press dubbed the phenomenon "Beatlemania".[65] Although not billed as tour leaders, the Beatles overshadowed American acts Tommy Roe and Chris Montez during the February engagements and assumed top billing "by audience demand", something no British act had previously accomplished while touring with artists from the US.[66] A similar situation arose during their May–June tour with Roy Orbison.[67]
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In late October, the Beatles began a five-day tour of Sweden, their first time abroad since the final Hamburg engagement of December 1962.[69] On their return to the UK on 31 October, several hundred screaming fans greeted them in heavy rain at Heathrow Airport. Around 50 to 100 journalists and photographers, as well as representatives from the BBC, also joined the airport reception, the first of more than 100 such events.[70] The next day, the band began its fourth tour of Britain within nine months, this one scheduled for six weeks.[71] In mid-November, as Beatlemania intensified, police resorted to using high-pressure water hoses to control the crowd before a concert in Plymouth.[72]
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Please Please Me maintained the top position on the Record Retailer chart for 30 weeks, only to be displaced by its follow-up, With the Beatles,[73] which EMI released on 22 November to record advance orders of 270,000 copies. The LP topped a half-million albums sold in one week.[74] Recorded between July and October, With the Beatles made better use of studio production techniques than its predecessor.[75] It held the top spot for 21 weeks with a chart life of 40 weeks.[76] Erlewine described the LP as "a sequel of the highest order – one that betters the original".[77]
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In a reversal of then standard practice, EMI released the album ahead of the impending single "I Want to Hold Your Hand", with the song excluded to maximise the single's sales.[78] The album caught the attention of music critic William Mann of The Times, who suggested that Lennon and McCartney were "the outstanding English composers of 1963".[75] The newspaper published a series of articles in which Mann offered detailed analyses of the music, lending it respectability.[79] With the Beatles became the second album in UK chart history to sell a million copies, a figure previously reached only by the 1958 South Pacific soundtrack.[80] When writing the sleeve notes for the album, the band's press officer, Tony Barrow, used the superlative the "fabulous foursome", which the media widely adopted as "the Fab Four".[81]
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EMI's American subsidiary, Capitol Records, hindered the Beatles' releases in the United States for more than a year by initially declining to issue their music, including their first three singles. Concurrent negotiations with the independent US label Vee-Jay led to the release of some of the songs in 1963, but not all.[82] Vee-Jay finished preparation for the album Introducing... The Beatles, culled from most of the songs of Parlophone's Please Please Me, but a management shake-up led to the album not being released.[nb 2] Then when it surfaced that the label did not report royalties on their sales, the licence Vee-Jay signed with EMI was voided.[84] A new licence was granted to the Swan label for the single "She Loves You". The record received some airplay in the Tidewater area of Virginia by Gene Loving of radio station WGH and was featured on the "Rate-a-Record" segment of American Bandstand, but it failed to catch on nationally.[85]
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Epstein arranged for a $40,000 US marketing campaign. American chart success began after disc jockey Carroll James of AM radio station WWDC, in Washington, DC, obtained a copy of the British single "I Want to Hold Your Hand" in mid-December 1963 and began playing it on-air.[86] Taped copies of the song soon circulated among other radio stations throughout the US. This caused an increase in demand, leading Capitol to bring forward the release of "I Want to Hold Your Hand" by three weeks.[87] Issued on 26 December, with the band's previously scheduled debut there just weeks away, "I Want to Hold Your Hand" sold a million copies, becoming a number-one hit in the US by mid-January.[88] In its wake, Vee-Jay released Introducing... The Beatles[89] to go along with Capitol's debut album, Meet the Beatles!, while Swan reactivated production of "She Loves You".[90]
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On 7 February 1964, the Beatles left the UK with an estimated 4,000 fans gathered at Heathrow, waving and screaming as the aircraft took off.[91] Upon landing at New York's John F. Kennedy Airport, an uproarious crowd estimated at 3,000 greeted them.[92] They gave their first live US television performance two days later on The Ed Sullivan Show, watched by approximately 73 million viewers in over 23 million households,[93] or 34 percent of the American population. Biographer Jonathan Gould writes that, according to the Nielsen rating service, it was "the largest audience that had ever been recorded for an American television program".[94] The next morning, the Beatles awoke to a largely negative critical consensus in the US,[95] but a day later at their first US concert, Beatlemania erupted at the Washington Coliseum.[96] Back in New York the following day, the Beatles met with another strong reception during two shows at Carnegie Hall.[93] The band flew to Florida, where they appeared on the weekly Ed Sullivan Show a second time, before another 70 million viewers, before returning to the UK on 22 February.[97]
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The Beatles' first visit to the US took place when the nation were still mourning the assassination of President John F. Kennedy the previous November.[98] Commentators often suggest that for many, particularly the young, the Beatles' performances reignited the sense of excitement and possibility that momentarily faded in the wake of the assassination, and helped make way for the revolutionary social changes to come in the decade.[99] Their hairstyle, unusually long for the era and mocked by many adults,[14] became an emblem of rebellion to the burgeoning youth culture.[100]
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The group's popularity generated unprecedented interest in British music, and many other UK acts subsequently made their American debuts, successfully touring over the next three years in what was termed the British Invasion.[101] The Beatles' success in the US opened the door for a successive string of British beat groups and pop acts such as the Dave Clark Five, the Animals, Petula Clark, the Kinks, and the Rolling Stones to achieve success in America.[102] During the week of 4 April 1964, the Beatles held twelve positions on the Billboard Hot 100 singles chart, including the top five.[103][nb 3]
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Capitol Records' lack of interest throughout 1963 did not go unnoticed, and a competitor, United Artists Records, encouraged their film division to offer the Beatles a three-motion-picture deal, primarily for the commercial potential of the soundtracks in the US.[105] Directed by Richard Lester, A Hard Day's Night involved the band for six weeks in March–April 1964 as they played themselves in a musical comedy.[106] The film premiered in London and New York in July and August, respectively, and was an international success, with some critics drawing a comparison with the Marx Brothers.[107]
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United Artists released a full soundtrack album for the North American market, combining Beatles songs and Martin's orchestral score; elsewhere, the group's third studio LP, A Hard Day's Night, contained songs from the film on side one and other new recordings on side two.[108] According to Erlewine, the album saw them "truly coming into their own as a band. All of the disparate influences on their first two albums coalesced into a bright, joyous, original sound, filled with ringing guitars and irresistible melodies."[109] That "ringing guitar" sound was primarily the product of Harrison's 12-string electric Rickenbacker, a prototype given to him by the manufacturer, which made its debut on the record.[110][nb 4]
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Touring internationally in June and July, the Beatles staged 37 shows over 27 days in Denmark, the Netherlands, Hong Kong, Australia and New Zealand.[111][nb 5] In August and September they returned to the US, with a 30-concert tour of 23 cities.[113] Generating intense interest once again, the month-long tour attracted between 10,000 and 20,000 fans to each 30-minute performance in cities from San Francisco to New York.[113]
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In August, journalist Al Aronowitz arranged for the Beatles to meet Bob Dylan.[114] Visiting the band in their New York hotel suite, Dylan introduced them to cannabis.[115] Gould points out the musical and cultural significance of this meeting, before which the musicians' respective fanbases were "perceived as inhabiting two separate subcultural worlds": Dylan's audience of "college kids with artistic or intellectual leanings, a dawning political and social idealism, and a mildly bohemian style" contrasted with their fans, "veritable 'teenyboppers' – kids in high school or grade school whose lives were totally wrapped up in the commercialised popular culture of television, radio, pop records, fan magazines, and teen fashion. To many of Dylan's followers in the folk music scene, the Beatles were seen as idolaters, not idealists."[116] Within six months of the meeting, according to Gould, "Lennon would be making records on which he openly imitated Dylan's nasal drone, brittle strum, and introspective vocal persona"; and six months after that, Dylan began performing with a backing band and electric instrumentation, and "dressed in the height of Mod fashion".[117] As a result, Gould continues, the traditional division between folk and rock enthusiasts "nearly evaporated", as the Beatles' fans began to mature in their outlook and Dylan's audience embraced the new, youth-driven pop culture.[117]
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During the 1964 US tour, the group were confronted with racial segregation in the country at the time.[118][119] When informed that the venue for their 11 September concert, the Gator Bowl in Jacksonville, Florida, was segregated, the Beatles said they would refuse to perform unless the audience was integrated.[120][118][119] Lennon stated: "We never play to segregated audiences and we aren't going to start now ... I'd sooner lose our appearance money."[118] City officials relented and agreed to allow an integrated show.[118] The group also cancelled their reservations at the whites-only Hotel George Washington in Jacksonville.[119] For their subsequent US tours in 1965 and 1966, the Beatles included clauses in contracts stipulating that shows be integrated.[119][121]
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According to Gould, Beatles for Sale, the Beatles' fourth studio LP, evidenced a growing conflict between the commercial pressures of their global success and their creative ambitions.[122] They had intended the album, recorded between August and October 1964,[123] to continue the format established by A Hard Day's Night which, unlike their first two LPs, contained only original songs.[122] They had nearly exhausted their backlog of songs on the previous album, however, and given the challenges constant international touring posed to their songwriting efforts, Lennon admitted, "Material's becoming a hell of a problem".[124] As a result, six covers from their extensive repertoire were chosen to complete the album. Released in early December, its eight original compositions stood out, demonstrating the growing maturity of the Lennon–McCartney songwriting partnership.[122]
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In early 1965, following a dinner with Lennon, Harrison and their wives, Harrison's dentist, John Riley, secretly added LSD to their coffee.[125] Lennon described the experience: "It was just terrifying, but it was fantastic. I was pretty stunned for a month or two."[126] He and Harrison subsequently became regular users of the drug, joined by Starr on at least one occasion. Harrison's use of psychedelic drugs encouraged his path to meditation and Hinduism. He commented: "For me, it was like a flash. The first time I had acid, it just opened up something in my head that was inside of me, and I realized a lot of things. I didn't learn them because I already knew them, but that happened to be the key that opened the door to reveal them. From the moment I had that, I wanted to have it all the time – these thoughts about the yogis and the Himalayas, and Ravi's music."[127][128] McCartney was initially reluctant to try it, but eventually did so in late 1966.[129] He became the first Beatle to discuss LSD publicly, declaring in a magazine interview that "it opened my eyes" and "made me a better, more honest, more tolerant member of society".[130]
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Controversy erupted in June 1965 when Queen Elizabeth II appointed all four Beatles Members of the Order of the British Empire (MBE) after Prime Minister Harold Wilson nominated them for the award.[131] In protest – the honour was at that time primarily bestowed upon military veterans and civic leaders – some conservative MBE recipients returned their insignia.[132]
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In July, the Beatles' second film, Help!, was released, again directed by Lester. Described as "mainly a relentless spoof of Bond",[133] it inspired a mixed response among both reviewers and the band. McCartney said: "Help! was great but it wasn't our film – we were sort of guest stars. It was fun, but basically, as an idea for a film, it was a bit wrong."[134] The soundtrack was dominated by Lennon, who wrote and sang lead on most of its songs, including the two singles: "Help!" and "Ticket to Ride".[135] The accompanying Help! album, the group's fifth studio LP, mirrored A Hard Day's Night by featuring soundtrack songs on side one and additional songs from the same sessions on side two.[136] The LP contained all original material save for two covers, "Act Naturally" and "Dizzy Miss Lizzy"; they were the last covers the band would include on an album, except for Let It Be's brief rendition of the traditional Liverpool folk song "Maggie Mae".[137] The band expanded their use of vocal overdubs on Help! and incorporated classical instruments into some arrangements, including a string quartet on the pop ballad "Yesterday".[138] Composed by and sung by McCartney – none of the other Beatles perform on the recording[139] – "Yesterday" inspired the most cover versions of any song ever written.[140] With Help!, the Beatles became the first rock group to be nominated for a Grammy Award for Album of the Year.[141]
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The group's third US tour opened with a performance before a world-record crowd of 55,600 at New York's Shea Stadium on 15 August 1965 – "perhaps the most famous of all Beatles' concerts", in Lewisohn's description.[142] A further nine successful concerts followed in other American cities. At a show in Atlanta, the Beatles gave one of the first live performances ever to make use of a foldback system of on-stage monitor speakers.[143] Towards the end of the tour, they met with Elvis Presley, a foundational musical influence on the band, who invited them to his home in Beverly Hills.[144][145] September saw the launch of an American Saturday-morning cartoon series, The Beatles, that echoed A Hard Day's Night's slapstick antics over its two-year original run.[146] The series was a historical milestone as the first weekly television series to feature animated versions of real, living people.[147]
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In mid-October 1965, the Beatles entered the recording studio; for the first time when making an album, they had an extended period without other major commitments.[148] Until this time, according to George Martin, "we had been making albums rather like a collection of singles. Now we were really beginning to think about albums as a bit of art on their own."[149] Released in December, Rubber Soul was hailed by critics as a major step forward in the maturity and complexity of the band's music.[150] Their thematic reach was beginning to expand as they embraced deeper aspects of romance and philosophy.[151] NEMS executive Peter Brown attributed the new musical direction to "the Beatles' now habitual use of marijuana",[152] an assertion confirmed by the band – Lennon referred to it as "the pot album",[153] and Starr said: "Grass was really influential in a lot of our changes, especially with the writers. And because they were writing different material, we were playing differently."[153] After Help!'s foray into the world of classical music with flutes and strings, Harrison's introduction of a sitar on "Norwegian Wood (This Bird Has Flown)" marked a further progression outside the traditional boundaries of popular music. As their lyrics grew more artful, fans began to study them for deeper meaning. Of "Norwegian Wood" Lennon commented: "I was trying to be sophisticated in writing about an affair ... but in such a smokescreen way that you couldn't tell."[154]
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While many of Rubber Soul's prominent songs were the product of Lennon and McCartney's collaborative songwriting,[155] it also featured distinct compositions from each,[156] though they continued to share official credit. The song "In My Life", of which each later claimed lead authorship, is considered a highlight of the entire Lennon–McCartney catalogue.[157] Harrison called Rubber Soul his "favourite album"[153] and Starr referred to it as "the departure record".[158] McCartney has said, "We'd had our cute period, and now it was time to expand."[159] However, recording engineer Norman Smith later stated that the studio sessions revealed signs of growing conflict within the group – "the clash between John and Paul was becoming obvious", he wrote, and "as far as Paul was concerned, George could do no right".[160] In 2003, Rolling Stone ranked Rubber Soul fifth among "The 500 Greatest Albums of All Time",[161] and AllMusic's Richie Unterberger describes it as "one of the classic folk-rock records".[162]
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Capitol Records, from December 1963 when it began issuing Beatles recordings for the US market, exercised complete control over format,[82] compiling distinct US albums from the band's recordings and issuing songs of their choosing as singles.[163][nb 6] In June 1966, Yesterday and Today, one of Capitol's compilation albums, caused an uproar with its cover, which portrayed the grinning Beatles dressed in butcher's overalls, accompanied by raw meat and mutilated plastic baby dolls. According to Beatles biographer Bill Harry, it has been incorrectly suggested that this was meant as a satirical response to the way Capitol had "butchered" the US versions of the band's albums.[165] Thousands of copies of the LP had a new cover pasted over the original; an unpeeled "first-state" copy fetched $10,500 at a December 2005 auction.[166] In England, meanwhile, Harrison met sitar maestro Ravi Shankar, who agreed to train him on the instrument.[167]
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During a tour of the Philippines the month after the Yesterday and Today furore, the Beatles unintentionally snubbed the nation's first lady, Imelda Marcos, who had expected them to attend a breakfast reception at the Presidential Palace.[168] When presented with the invitation, Epstein politely declined on the band members' behalf, as it had never been his policy to accept such official invitations.[169] They soon found that the Marcos regime was unaccustomed to taking no for an answer. The resulting riots endangered the group and they escaped the country with difficulty.[170] Immediately afterwards, the band members visited India for the first time.[171]
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– John Lennon, 1966[172]
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Almost as soon as they returned home, the Beatles faced a fierce backlash from US religious and social conservatives (as well as the Ku Klux Klan) over a comment Lennon had made in a March interview with British reporter Maureen Cleave.[173] "Christianity will go", Lennon had said. "It will vanish and shrink. I needn't argue about that; I'm right and I will be proved right ... Jesus was alright but his disciples were thick and ordinary. It's them twisting it that ruins it for me."[174] His comments went virtually unnoticed in England, but when US teenage fan magazine Datebook printed them five months later, it sparked a controversy with Christians in America's conservative Bible Belt region.[173] The Vatican issued a protest, and bans on Beatles' records were imposed by Spanish and Dutch stations and South Africa's national broadcasting service.[175] Epstein accused Datebook of having taken Lennon's words out of context. At a press conference Lennon pointed out, "If I'd said television was more popular than Jesus, I might have got away with it."[176] He claimed that he was referring to how other people viewed their success, but at the prompting of reporters, he concluded: "If you want me to apologise, if that will make you happy, then okay, I'm sorry."[176]
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Released in August, a week before the Beatles' final tour, Revolver marked another artistic step forward for the group.[177] The album featured sophisticated songwriting, studio experimentation, and a greatly expanded repertoire of musical styles, ranging from innovative classical string arrangements to psychedelic rock.[177] Abandoning the customary group photograph, its Aubrey Beardsley-inspired cover – designed by Klaus Voormann, a friend of the band since their Hamburg days – was a monochrome collage and line drawing caricature of the group.[177] The album was preceded by the single "Paperback Writer", backed by "Rain".[178] Short promotional films were made for both songs; described by cultural historian Saul Austerlitz as "among the first true music videos",[179] they aired on The Ed Sullivan Show and Top of the Pops in June.[180]
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Among the experimental songs that Revolver featured was "Tomorrow Never Knows", the lyrics for which Lennon drew from Timothy Leary's The Psychedelic Experience: A Manual Based on the Tibetan Book of the Dead. Its creation involved eight tape decks distributed about the EMI building, each staffed by an engineer or band member, who randomly varied the movement of a tape loop while Martin created a composite recording by sampling the incoming data.[181] McCartney's "Eleanor Rigby" made prominent use of a string octet; Gould describes it as "a true hybrid, conforming to no recognisable style or genre of song".[182] Harrison was developing as a songwriter, and three of his compositions earned a place on the record.[183] In 2003, Rolling Stone ranked Revolver as the third greatest album of all time.[161]
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As preparations were made for a tour of the US, the Beatles knew that their music would hardly be heard. Having originally used Vox AC30 amplifiers, they later acquired more powerful 100-watt amplifiers, specially designed by Vox for them as they moved into larger venues in 1964, but these were still inadequate. Struggling to compete with the volume of sound generated by screaming fans, the band had grown increasingly bored with the routine of performing live.[184] Recognising that their shows were no longer about the music, they decided to make the August tour their last.[185]
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The band performed none of their new songs on the tour.[186] In Chris Ingham's description, they were very much "studio creations ... and there was no way a four-piece rock 'n' roll group could do them justice, particularly through the desensitising wall of the fans' screams. 'Live Beatles' and 'Studio Beatles' had become entirely different beasts."[187] The band's concert at San Francisco's Candlestick Park on 29 August was their last commercial concert.[188] It marked the end of four years dominated by almost nonstop touring that included over 1,400 concert appearances internationally.[189]
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Freed from the burden of touring, the Beatles embraced an increasingly experimental approach as they recorded Sgt. Pepper's Lonely Hearts Club Band, beginning in late November 1966.[191] According to engineer Geoff Emerick, the album's recording took over 700 hours.[192] He recalled the band's insistence "that everything on Sgt. Pepper had to be different. We had microphones right down in the bells of brass instruments and headphones turned into microphones attached to violins. We used giant primitive oscillators to vary the speed of instruments and vocals and we had tapes chopped to pieces and stuck together upside down and the wrong way around."[193] Parts of "A Day in the Life" featured a 40-piece orchestra.[193] The sessions initially yielded the non-album double A-side single "Strawberry Fields Forever"/"Penny Lane" in February 1967;[194] the Sgt. Pepper LP followed with a rush-release in May.[195] The musical complexity of the records, created using relatively primitive four-track recording technology, astounded contemporary artists.[190] Among music critics, acclaim for the album was virtually universal.[196] Gould writes:
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The overwhelming consensus is that the Beatles had created a popular masterpiece: a rich, sustained, and overflowing work of collaborative genius whose bold ambition and startling originality dramatically enlarged the possibilities and raised the expectations of what the experience of listening to popular music on record could be. On the basis of this perception, Sgt. Pepper became the catalyst for an explosion of mass enthusiasm for album-formatted rock that would revolutionise both the aesthetics and the economics of the record business in ways that far outstripped the earlier pop explosions triggered by the Elvis phenomenon of 1956 and the Beatlemania phenomenon of 1963.[197]
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In the wake of Sgt. Pepper, the underground and mainstream press widely publicised the Beatles as leaders of youth culture, as well as "lifestyle revolutionaries".[3] The album was the first major pop/rock LP to include its complete lyrics, which appeared on the back cover.[198][199] Those lyrics were the subject of critical analysis; for instance, in late 1967 the album was the subject of a scholarly inquiry by American literary critic and professor of English Richard Poirier, who observed that his students were "listening to the group's music with a degree of engagement that he, as a teacher of literature, could only envy".[200][nb 7] The elaborate cover also attracted considerable interest and study.[201] A collage designed by pop artists Peter Blake and Jann Haworth, it depicted the group as the fictional band referred to in the album's title track[202] standing in front of a crowd of famous people.[203] The heavy moustaches worn by the group reflected the growing influence of hippie style,[204] while cultural historian Jonathan Harris describes their "brightly coloured parodies of military uniforms" as a knowingly "anti-authoritarian and anti-establishment" display.[205]
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Sgt. Pepper topped the UK charts for 23 consecutive weeks, with a further four weeks at number one in the period through to February 1968.[206] With 2.5 million copies sold within three months of its release,[207] Sgt. Pepper's initial commercial success exceeded that of all previous Beatles albums.[208] It sustained its immense popularity into the 21st century while breaking numerous sales records.[209] In 2003, Rolling Stone ranked Sgt. Pepper foremost on its list of the greatest albums of all time.[161]
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Two Beatles film projects were conceived within weeks of completing Sgt. Pepper: Magical Mystery Tour, a one-hour television film, and Yellow Submarine, an animated feature-length film produced by United Artists.[210] The group began recording music for the former in late April 1967, but the project then lay dormant as they focused on recording songs for the latter.[211] On 25 June, the Beatles performed their forthcoming single "All You Need Is Love" to an estimated 350 million viewers on Our World, the first live global television link.[212] Released a week later, during the Summer of Love, the song was adopted as a flower power anthem.[213] The Beatles' use of psychedelic drugs was at its height during that summer.[214] In July and August, the group pursued interests related to similar utopian-based ideology, including a week-long investigation into the possibility of starting an island-based commune off the coast of Greece.[215]
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On 24 August, the group were introduced to Maharishi Mahesh Yogi in London. The next day, they travelled to Bangor for his Transcendental Meditation retreat. On 27 August, their manager's assistant, Peter Brown, phoned to inform them that Epstein had died.[216] The coroner ruled the death an accidental carbitol overdose, although it was widely rumoured to be a suicide.[217][nb 8] His death left the group disoriented and fearful about the future.[219] Lennon recalled: "We collapsed. I knew that we were in trouble then. I didn't really have any misconceptions about our ability to do anything other than play music, and I was scared. I thought, 'We've fuckin' had it now.'"[220] Harrison's then-wife Pattie Boyd remembered that "Paul and George were in complete shock. I don't think it could have been worse if they had heard that their own fathers had dropped dead."[221] During a band meeting in September, McCartney recommended that the band proceed with Magical Mystery Tour.[211]
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The Magical Mystery Tour soundtrack was released in the UK as a six-track double extended play (EP) in early December 1967.[82][222] It was the first example of a double EP in the UK.[223][224] The record carried on the psychedelic vein of Sgt. Pepper,[225] however, in line with the band's wishes, the packaging reinforced the idea that the release was a film soundtrack rather than a follow-up to Sgt. Pepper.[222] In the US, the soundtrack appeared as an identically titled LP that also included five tracks from the band's recent singles.[104] In its first three weeks, the album set a record for the highest initial sales of any Capitol LP, and it is the only Capitol compilation later to be adopted in the band's official canon of studio albums.[226]
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Magical Mystery Tour first aired on Boxing Day to an audience of approximately 15 million.[227] Largely directed by McCartney, the film was the band's first critical failure in the UK.[228] It was dismissed as "blatant rubbish" by the Daily Express; the Daily Mail called it "a colossal conceit"; and The Guardian labelled the film "a kind of fantasy morality play about the grossness and warmth and stupidity of the audience".[229] Gould describes it as "a great deal of raw footage showing a group of people getting on, getting off, and riding on a bus".[229] Although the viewership figures were respectable, its slating in the press led US television networks to lose interest in broadcasting the film.[230]
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The group were less involved with Yellow Submarine, which only featured the band appearing as themselves for a short live-action segment.[231] Premiering in July 1968, the film featured cartoon versions of the band members and a soundtrack with eleven of their songs, including four unreleased studio recordings that made their debut in the film.[232] Critics praised the film for its music, humour and innovative visual style.[233] A soundtrack LP was issued seven months later; it contained those four new songs, the title track (already issued on Revolver), "All You Need Is Love" (already issued as a single and on the US Magical Mystery Tour LP) and seven instrumental pieces composed by Martin.[234]
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In February 1968, the Beatles travelled to Maharishi Mahesh Yogi's ashram in Rishikesh, India, to take part in a three-month meditation "Guide Course". Their time in India marked one of the band's most prolific periods, yielding numerous songs, including a majority of those on their next album.[235] However, Starr left after only ten days, unable to stomach the food, and McCartney eventually grew bored and departed a month later.[236] For Lennon and Harrison, creativity turned to question when an electronics technician known as Magic Alex suggested that the Maharishi was attempting to manipulate them.[237] When he alleged that the Maharishi had made sexual advances to women attendees, a persuaded Lennon left abruptly just two months into the course, bringing an unconvinced Harrison and the remainder of the group's entourage with him.[236] In anger, Lennon wrote a scathing song titled "Maharishi", renamed "Sexy Sadie" to avoid potential legal issues. McCartney said, "We made a mistake. We thought there was more to him than there was."[237]
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In May, Lennon and McCartney traveled to New York for the public unveiling of the Beatles' new business venture, Apple Corps.[238] It was initially formed several months earlier as part of a plan to create a tax-effective business structure, but the band then desired to extend the corporation to other pursuits, including record distribution, peace activism, and education.[239] McCartney described Apple as "rather like a Western communism".[240] The enterprise drained the group financially with a series of unsuccessful projects[241] handled largely by members of the Beatles' entourage, who were given their jobs regardless of talent and experience.[242] Among its numerous subsidiaries were Apple Electronics, established to foster technological innovations with Magic Alex at the head, and Apple Retailing, which opened the short-lived Apple Boutique in London.[243] Harrison later said, "Basically, it was chaos ... John and Paul got carried away with the idea and blew millions, and Ringo and I just had to go along with it."[240]
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From late May to mid-October 1968, the group recorded what became The Beatles, a double LP commonly known as "the White Album" for its virtually featureless cover.[245] During this time, relations between the members grew openly divisive.[246] Starr quit for two weeks, leaving his bandmates to record "Back in the U.S.S.R." and "Dear Prudence" as a trio.[247] Lennon had lost interest in collaborating with McCartney,[248] whose contribution "Ob-La-Di, Ob-La-Da" he scorned as "granny music shit".[249] Tensions were further aggravated by Lennon's romantic preoccupation with avant-garde artist Yoko Ono, whom he insisted on bringing to the sessions despite the group's well-established understanding that girlfriends were not allowed in the studio.[250] McCartney has recalled that the album "wasn't a pleasant one to make".[251] He and Lennon identified the sessions as the start of the band's break-up.[252][253]
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With the record, the band executed a wider range of musical styles[254] and broke with their recent tradition of incorporating several musical styles in one song by keeping each piece of music consistently faithful to a select genre.[255] During the sessions, the group upgraded to an eight-track tape console, which made it easier for them to layer tracks piecemeal, while the members often recorded independently of each other, affording the album a reputation as a collection of solo recordings rather than a unified group effort.[256] Describing the double album, Lennon later said: "Every track is an individual track; there isn't any Beatle music on it. [It's] John and the band, Paul and the band, George and the band."[257] The sessions also produced the Beatles' longest song yet, "Hey Jude", released in August as a non-album single with "Revolution".[258]
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Issued in November, the White Album was the band's first Apple Records album release, although EMI continued to own their recordings.[259] The record attracted more than 2 million advance orders, selling nearly 4 million copies in the US in little over a month, and its tracks dominated the playlists of American radio stations.[260] Its lyric content was the focus of much analysis by the counterculture.[261] Despite its popularity, reviewers were largely confused by the album's content, and it failed to inspire the level of critical writing that Sgt. Pepper had.[260] General critical opinion eventually turned in favour of the White Album, and in 2003, Rolling Stone ranked it as the tenth greatest album of all time.[161]
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Although Let It Be was the Beatles' final album release, it was largely recorded before Abbey Road. The project's impetus came from an idea Martin attributes to McCartney, who suggested they "record an album of new material and rehearse it, then perform it before a live audience for the very first time – on record and on film".[262] Originally intended for a one-hour television programme to be called Beatles at Work, in the event much of the album's content came from studio work beginning in January 1969, many hours of which were captured on film by director Michael Lindsay-Hogg.[262][263] Martin said that the project was "not at all a happy recording experience. It was a time when relations between the Beatles were at their lowest ebb."[262] Lennon described the largely impromptu sessions as "hell ... the most miserable ... on Earth", and Harrison, "the low of all-time".[264] Irritated by McCartney and Lennon, Harrison walked out for five days. Upon returning, he threatened to leave the band unless they "abandon[ed] all talk of live performance" and instead focused on finishing a new album, initially titled Get Back, using songs recorded for the TV special.[265] He also demanded they cease work at Twickenham Film Studios, where the sessions had begun, and relocate to the newly finished Apple Studio. His bandmates agreed, and it was decided to salvage the footage shot for the TV production for use in a feature film.[266]
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To alleviate tensions within the band and improve the quality of their live sound, Harrison invited keyboardist Billy Preston to participate in the last nine days of sessions.[267] Preston received label billing on the "Get Back" single – the only musician ever to receive that acknowledgment on an official Beatles release.[268] After the rehearsals, the band could not agree on a location to film a concert, rejecting several ideas, including a boat at sea, a lunatic asylum, the Tunisian desert, and the Colosseum.[262] Ultimately, what would be their final live performance was filmed on the rooftop of the Apple Corps building at 3 Savile Row, London, on 30 January 1969.[269] Five weeks later, engineer Glyn Johns, whom Lewisohn describes as Get Back's "uncredited producer", began work assembling an album, given "free rein" as the band "all but washed their hands of the entire project".[270]
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New strains developed between the band members regarding the appointment of a financial adviser, the need for which had become evident without Epstein to manage business affairs. Lennon, Harrison and Starr favoured Allen Klein, who had managed the Rolling Stones and Sam Cooke;[271] McCartney wanted Lee and John Eastman – father and brother, respectively, of Linda Eastman,[272] whom McCartney married on 12 March.[273] Agreement could not be reached, so both Klein and the Eastmans were temporarily appointed: Klein as the Beatles' business manager and the Eastmans as their lawyers.[274][275] Further conflict ensued, however, and financial opportunities were lost.[271] On 8 May, Klein was named sole manager of the band,[276] the Eastmans having previously been dismissed as the Beatles' lawyers. McCartney refused to sign the management contract with Klein, but he was out-voted by the other Beatles.[277]
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Martin stated that he was surprised when McCartney asked him to produce another album, as the Get Back sessions had been "a miserable experience" and he had "thought it was the end of the road for all of us".[278] The primary recording sessions for Abbey Road began on 2 July.[279] Lennon, who rejected Martin's proposed format of a "continuously moving piece of music", wanted his and McCartney's songs to occupy separate sides of the album.[280] The eventual format, with individually composed songs on the first side and the second consisting largely of a medley, was McCartney's suggested compromise.[280] Emerick noted that the replacement of the studio's valve mixing console with a transistorised one yielded a less punchy sound, leaving the group frustrated at the thinner tone and lack of impact and contributing to its "kinder, gentler" feel relative to their previous albums.[281]
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On 4 July, the first solo single by a Beatle was released: Lennon's "Give Peace a Chance", credited to the Plastic Ono Band. The completion and mixing of "I Want You (She's So Heavy)" on 20 August was the last occasion on which all four Beatles were together in the same studio.[282] On 8 September, while Starr was in hospital, the other band members met to discuss recording a new album. They considered a different approach to songwriting by ending the Lennon–McCartney pretense and having four compositions apiece from Lennon, McCartney and Harrison, with two from Starr and a lead single around Christmas.[283] On 20 September, Lennon announced his departure to the rest of the group but agreed to withhold a public announcement to avoid undermining sales of the forthcoming album.[284]
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Released on 26 September, Abbey Road sold four million copies within three months and topped the UK charts for a total of seventeen weeks.[285] Its second track, the ballad "Something", was issued as a single – the only Harrison composition that appeared as a Beatles A-side.[286] Abbey Road received mixed reviews, although the medley met with general acclaim.[285] Unterberger considers it "a fitting swan song for the group", containing "some of the greatest harmonies to be heard on any rock record".[287] Musicologist and author Ian MacDonald calls the album "erratic and often hollow", despite the "semblance of unity and coherence" offered by the medley.[288] Martin singled it out as his favourite Beatles album; Lennon said it was "competent" but had "no life in it".[281]
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For the still unfinished Get Back album, one last song, Harrison's "I Me Mine", was recorded on 3 January 1970. Lennon, in Denmark at the time, did not participate.[289] In March, rejecting the work Johns had done on the project, now retitled Let It Be, Klein gave the session tapes to American producer Phil Spector, who had recently produced Lennon's solo single "Instant Karma!"[290] In addition to remixing the material, Spector edited, spliced and overdubbed several of the recordings that had been intended as "live". McCartney was unhappy with the producer's approach and particularly dissatisfied with the lavish orchestration on "The Long and Winding Road", which involved a fourteen-voice choir and 36-piece instrumental ensemble.[291] McCartney's demands that the alterations to the song be reverted were ignored,[292] and he publicly announced his departure from the band on 10 April, a week before the release of his first, self-titled solo album.[291][293]
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On 8 May 1970, Let It Be was released. Its accompanying single, "The Long and Winding Road", was the Beatles' last; it was released in the US, but not in the UK.[178] The Let It Be documentary film followed later that month, and would win the 1970 Academy Award for Best Original Song Score.[294] Sunday Telegraph critic Penelope Gilliatt called it "a very bad film and a touching one ... about the breaking apart of this reassuring, geometrically perfect, once apparently ageless family of siblings".[295] Several reviewers stated that some of the performances in the film sounded better than their analogous album tracks.[296] Describing Let It Be as the "only Beatles album to occasion negative, even hostile reviews", Unterberger calls it "on the whole underrated"; he singles out "some good moments of straight hard rock in 'I've Got a Feeling' and 'Dig a Pony'", and praises "Let It Be", "Get Back", and "the folky 'Two of Us', with John and Paul harmonising together".[297]
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McCartney filed suit for the dissolution of the Beatles' contractual partnership on 31 December 1970.[298] Legal disputes continued long after their break-up, and the dissolution was not formalised until 29 December 1974,[299] when Lennon signed the paperwork terminating the partnership while on vacation with his family at Walt Disney World Resort in Florida.[300]
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Lennon, McCartney, Harrison and Starr all released solo albums in 1970. Their solo records sometimes involved one or more of the others;[301] Starr's Ringo (1973) was the only album to include compositions and performances by all four ex-Beatles, albeit on separate songs. With Starr's participation, Harrison staged the Concert for Bangladesh in New York City in August 1971.[302] Other than an unreleased jam session in 1974, later bootlegged as A Toot and a Snore in '74, Lennon and McCartney never recorded together again.[303]
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Two double-LP sets of the Beatles' greatest hits, compiled by Klein, 1962–1966 and 1967–1970, were released in 1973, at first under the Apple Records imprint.[304] Commonly known as the "Red Album" and "Blue Album", respectively, each has earned a Multi-Platinum certification in the US and a Platinum certification in the UK.[305][306] Between 1976 and 1982, EMI/Capitol released a wave of compilation albums without input from the ex-Beatles, starting with the double-disc compilation Rock 'n' Roll Music.[307] The only one to feature previously unreleased material was The Beatles at the Hollywood Bowl (1977); the first officially issued concert recordings by the group, it contained selections from two shows they played during their 1964 and 1965 US tours.[308][nb 9]
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The music and enduring fame of the Beatles were commercially exploited in various other ways, again often outside their creative control. In April 1974, the musical John, Paul, George, Ringo ... and Bert, written by Willy Russell and featuring singer Barbara Dickson, opened in London. It included, with permission from Northern Songs, eleven Lennon-McCartney compositions and one by Harrison, "Here Comes the Sun". Displeased with the production's use of his song, Harrison withdrew his permission to use it.[310] Later that year, the off-Broadway musical Sgt. Pepper's Lonely Hearts Club Band on the Road opened.[311] All This and World War II (1976) was an unorthodox nonfiction film that combined newsreel footage with covers of Beatles songs by performers ranging from Elton John and Keith Moon to the London Symphony Orchestra.[312] The Broadway musical Beatlemania, an unauthorised nostalgia revue, opened in early 1977 and proved popular, spinning off five separate touring productions.[313] In 1979, the band sued the producers, settling for several million dollars in damages.[313] Sgt. Pepper's Lonely Hearts Club Band (1978), a musical film starring the Bee Gees and Peter Frampton, was a commercial failure and an "artistic fiasco", according to Ingham.[314]
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Accompanying the wave of Beatles nostalgia and persistent reunion rumours in the US during the 1970s, several entrepreneurs made public offers to the Beatles for a reunion concert.[315] Promoter Bill Sargent first offered the Beatles $10 million for a reunion concert in 1974. He raised his offer to $30 million in January 1976 and then to $50 million the following month.[316][317] On 24 April 1976, during a broadcast of Saturday Night Live, producer Lorne Michaels offered the Beatles $3,000 to reunite on the show. Lennon and McCartney were watching the live broadcast at Lennon's apartment at the Dakota in New York, which was within driving distance of the NBC studio where the show was being broadcast. The former bandmates briefly entertained the idea of going to the studio and surprising Michaels by accepting his offer, but decided not to.[318]
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In December 1980, Lennon was shot and killed outside his New York City apartment. Harrison rewrote the lyrics of his song "All Those Years Ago" in Lennon's honour. With Starr on drums and McCartney and his wife, Linda, contributing backing vocals, the song was released as a single in May 1981.[319] McCartney's own tribute, "Here Today", appeared on his Tug of War album in April 1982.[320] In 1987, Harrison's Cloud Nine album included "When We Was Fab", a song about the Beatlemania era.[321]
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When the Beatles' studio albums were released on CD by EMI and Apple Corps in 1987, their catalogue was standardised throughout the world, establishing a canon of the twelve original studio LPs as issued in the UK plus the US LP version of Magical Mystery Tour.[322] All the remaining material from the singles and EPs that had not appeared on these thirteen studio albums was gathered on the two-volume compilation Past Masters (1988). Except for the Red and Blue albums, EMI deleted all its other Beatles compilations – including the Hollywood Bowl record – from its catalogue.[308]
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In 1988, the Beatles were inducted into the Rock and Roll Hall of Fame, their first year of eligibility. Harrison and Starr attended the ceremony with Lennon's widow, Yoko Ono, and his two sons, Julian and Sean.[323][324] McCartney declined to attend, citing unresolved "business differences" that would make him "feel like a complete hypocrite waving and smiling with them at a fake reunion".[324] The following year, EMI/Capitol settled a decade-long lawsuit filed by the band over royalties, clearing the way to commercially package previously unreleased material.[325][326]
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Live at the BBC, the first official release of unissued Beatles performances in seventeen years, appeared in 1994.[327] That same year McCartney, Harrison and Starr collaborated on the Anthology project. Anthology was the culmination of work begun in 1970, when Apple Corps director Neil Aspinall, their former road manager and personal assistant, had started to gather material for a documentary with the working title The Long and Winding Road.[328] Documenting their history in the band's own words, the Anthology project included the release of several unissued Beatles recordings. McCartney, Harrison and Starr also added new instrumental and vocal parts to songs recorded as demos by Lennon in the late 1970s.[329]
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During 1995–96, the project yielded a television miniseries, an eight-volume video set, and three two-CD/three-LP box sets featuring artwork by Klaus Voormann. Two songs based on Lennon demos, "Free as a Bird" and "Real Love", were issued as new Beatles singles. The releases were commercially successful and the television series was viewed by an estimated 400 million people.[330] In 1999, to coincide with the re-release of the 1968 film Yellow Submarine, an expanded soundtrack album, Yellow Submarine Songtrack, was issued.[331]
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The Beatles' 1, a compilation album of the band's British and American number-one hits, was released on 13 November 2000. It became the fastest-selling album of all time, with 3.6 million sold in its first week[332] and 13 million within a month.[333] It topped albums charts in at least 28 countries.[334] The compilation had sold 31 million copies globally by April 2009.[335]
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Harrison died from metastatic lung cancer in November 2001.[336][337][338] McCartney and Starr were among the musicians who performed at the Concert for George, organised by Eric Clapton and Harrison's widow, Olivia. The tribute event took place at the Royal Albert Hall on the first anniversary of Harrison's death.[339]
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In 2003, Let It Be... Naked, a reconceived version of the Let It Be album, with McCartney supervising production, was released. One of the main differences from the Spector-produced version was the omission of the original string arrangements.[340] It was a top ten hit in both Britain and America. The US album configurations from 1964 to 1965 were released as box sets in 2004 and 2006; The Capitol Albums, Volume 1 and Volume 2 included both stereo and mono versions based on the mixes that were prepared for vinyl at the time of the music's original American release.[341]
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As a soundtrack for Cirque du Soleil's Las Vegas Beatles stage revue, Love, George Martin and his son Giles remixed and blended 130 of the band's recordings to create what Martin called "a way of re-living the whole Beatles musical lifespan in a very condensed period".[342] The show premiered in June 2006, and the Love album was released that November.[343] In April 2009, Starr performed three songs with McCartney at a benefit concert held at New York's Radio City Music Hall and organised by McCartney.[344]
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On 9 September 2009, the Beatles' entire back catalogue was reissued following an extensive digital remastering process that lasted four years.[322] Stereo editions of all twelve original UK studio albums, along with Magical Mystery Tour and the Past Masters compilation, were released on compact disc both individually and as a box set.[345] A second collection, The Beatles in Mono, included remastered versions of every Beatles album released in true mono along with the original 1965 stereo mixes of Help! and Rubber Soul (both of which Martin had remixed for the 1987 editions).[346] The Beatles: Rock Band, a music video game in the Rock Band series, was issued on the same day.[347] In December 2009, the band's catalogue was officially released in FLAC and MP3 format in a limited edition of 30,000 USB flash drives.[348]
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Owing to a long-running royalty disagreement, the Beatles were among the last major artists to sign deals with online music services.[349] Residual disagreement emanating from Apple Corps' dispute with Apple, Inc., iTunes' owners, over the use of the name "Apple" was also partly responsible for the delay, although in 2008, McCartney stated that the main obstacle to making the Beatles' catalogue available online was that EMI "want[s] something we're not prepared to give them".[350] In 2010, the official canon of thirteen Beatles studio albums, Past Masters, and the "Red" and "Blue" greatest-hits albums were made available on iTunes.[351]
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In 2012, EMI's recorded music operations were sold to Universal Music Group. In order for Universal Music to acquire EMI, the European Union, for antitrust reasons, forced EMI to spin off assets including Parlophone. Universal was allowed to keep the Beatles' recorded music catalogue, managed by Capitol Records under its Capitol Music Group division.[352] The entire original Beatles album catalogue was also reissued on vinyl in 2012; available either individually or as a box set.[353]
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In 2013, a second volume of BBC recordings, titled On Air – Live at the BBC Volume 2, was released. That December saw the release of another 59 Beatles recordings on iTunes. The set, titled The Beatles Bootleg Recordings 1963, had the opportunity to gain a 70-year copyright extension conditional on the songs being published at least once before the end of 2013. Apple Records released the recordings on 17 December to prevent them from going into the public domain and had them taken down from iTunes later that same day. Fan reactions to the release were mixed, with one blogger saying "the hardcore Beatles collectors who are trying to obtain everything will already have these."[354][355]
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On 26 January 2014, McCartney and Starr performed together at the 56th Annual Grammy Awards, held at the Staples Center in Los Angeles.[356] The following day, The Night That Changed America: A Grammy Salute to The Beatles television special was taped in the Los Angeles Convention Center's West Hall. It aired on 9 February, the exact date of – and at the same time, and on the same network as – the original broadcast of the Beatles' first US television appearance on The Ed Sullivan Show, 50 years earlier. The special included performances of Beatles songs by current artists as well as by McCartney and Starr, archival footage, and interviews with the two surviving ex-Beatles carried out by David Letterman at the Ed Sullivan Theater.[357][358] In December 2015, the Beatles released their catalogue for streaming on various streaming music services including Spotify and Apple Music.[359]
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In September 2016, the documentary film The Beatles: Eight Days a Week was released. Directed by Ron Howard, it chronicled the Beatles' career during their touring years from 1962 to 1966, from their performances in Liverpool's the Cavern Club in 1961 to their final concert in San Francisco in 1966. The film was released theatrically on 15 September in the UK and the US, and started streaming on Hulu on 17 September. It received several awards and nominations, including for Best Documentary at the 70th British Academy Film Awards and the Outstanding Documentary or Nonfiction Special at the 69th Primetime Creative Arts Emmy Awards.[360] An expanded, remixed and remastered version of The Beatles at the Hollywood Bowl was released on 9 September, to coincide with the release of the film.[361][362]
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On 18 May 2017, Sirius XM Radio launched a 24/7 radio channel, The Beatles Channel. A week later, Sgt. Pepper's Lonely Hearts Club Band was reissued with new stereo mixes and unreleased material for the album's 50th anniversary.[363] Similar box sets were released for The Beatles in November 2018,[364] and Abbey Road in September 2019.[365] On the first week of October 2019, Abbey Road returned to number one on the UK Albums Chart. The Beatles broke their own record for the album with the longest gap between topping the charts as Abbey Road hit the top spot 50 years after its original release.[366]
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In August 2021, The Beatles: Get Back, a new documentary film directed by Peter Jackson utilising footage captured for what became the Let It Be film, will be released by Walt Disney Studios Motion Pictures in the US and Canada, with a global release to follow.[367]
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In Icons of Rock: An Encyclopedia of the Legends Who Changed Music Forever, Scott Schinder and Andy Schwartz describe the Beatles' musical evolution:
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In their initial incarnation as cheerful, wisecracking moptops, the Fab Four revolutionised the sound, style, and attitude of popular music and opened rock and roll's doors to a tidal wave of British rock acts. Their initial impact would have been enough to establish the Beatles as one of their era's most influential cultural forces, but they didn't stop there. Although their initial style was a highly original, irresistibly catchy synthesis of early American rock and roll and R&B, the Beatles spent the rest of the 1960s expanding rock's stylistic frontiers, consistently staking out new musical territory on each release. The band's increasingly sophisticated experimentation encompassed a variety of genres, including folk-rock, country, psychedelia, and baroque pop, without sacrificing the effortless mass appeal of their early work.[368]
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In The Beatles as Musicians, Walter Everett describes Lennon and McCartney's contrasting motivations and approaches to composition: "McCartney may be said to have constantly developed – as a means to entertain – a focused musical talent with an ear for counterpoint and other aspects of craft in the demonstration of a universally agreed-upon common language that he did much to enrich. Conversely, Lennon's mature music is best appreciated as the daring product of a largely unconscious, searching but undisciplined artistic sensibility."[369]
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Ian MacDonald describes McCartney as "a natural melodist – a creator of tunes capable of existing apart from their harmony". His melody lines are characterised as primarily "vertical", employing wide, consonant intervals which express his "extrovert energy and optimism". Conversely, Lennon's "sedentary, ironic personality" is reflected in a "horizontal" approach featuring minimal, dissonant intervals and repetitive melodies which rely on their harmonic accompaniment for interest: "Basically a realist, he instinctively kept his melodies close to the rhythms and cadences of speech, colouring his lyrics with bluesy tone and harmony rather than creating tunes that made striking shapes of their own."[370] MacDonald praises Harrison's lead guitar work for the role his "characterful lines and textural colourings" play in supporting Lennon and McCartney's parts, and describes Starr as "the father of modern pop/rock drumming".[371]
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The band's earliest influences include Elvis Presley, Carl Perkins, Little Richard and Chuck Berry.[372] During the Beatles' co-residency with Little Richard at the Star-Club in Hamburg, from April to May 1962, he advised them on the proper technique for performing his songs.[373] Of Presley, Lennon said, "Nothing really affected me until I heard Elvis. If there hadn't been Elvis, there would not have been the Beatles."[374] Other early influences include Buddy Holly, Eddie Cochran, Roy Orbison[375] and the Everly Brothers.[376]
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The Beatles continued to absorb influences long after their initial success, often finding new musical and lyrical avenues by listening to their contemporaries, including Bob Dylan, the Who, Frank Zappa, the Lovin' Spoonful, the Byrds and the Beach Boys, whose 1966 album Pet Sounds amazed and inspired McCartney.[377][378][379][380] Referring to the Beach Boys' creative leader, Martin later stated: "No one made a greater impact on the Beatles than Brian [Wilson]."[381] Ravi Shankar, with whom Harrison studied for six weeks in India in late 1966, had a significant effect on his musical development during the band's later years.[382]
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Originating as a skiffle group, the Beatles quickly embraced 1950s rock and roll and helped pioneer the Merseybeat genre,[383] and their repertoire ultimately expanded to include a broad variety of pop music.[384] Reflecting the range of styles they explored, Lennon said of Beatles for Sale, "You could call our new one a Beatles country-and-western LP",[385] while Gould credits Rubber Soul as "the instrument by which legions of folk-music enthusiasts were coaxed into the camp of pop".[386]
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Although the 1965 song "Yesterday" was not the first pop record to employ orchestral strings, it marked the group's first recorded use of classical music elements. Gould observes: "The more traditional sound of strings allowed for a fresh appreciation of their talent as composers by listeners who were otherwise allergic to the din of drums and electric guitars."[387] They continued to experiment with string arrangements to various effect; Sgt. Pepper's "She's Leaving Home", for instance, is "cast in the mold of a sentimental Victorian ballad", Gould writes, "its words and music filled with the clichés of musical melodrama".[388]
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The band's stylistic range expanded in another direction with their 1966 B-side "Rain", described by Martin Strong as "the first overtly psychedelic Beatles record".[389] Other psychedelic numbers followed, such as "Tomorrow Never Knows" (recorded before "Rain"), "Strawberry Fields Forever", "Lucy in the Sky with Diamonds" and "I Am the Walrus". The influence of Indian classical music was evident in Harrison's "The Inner Light", "Love You To" and "Within You Without You" – Gould describes the latter two as attempts "to replicate the raga form in miniature".[390]
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Innovation was the most striking feature of their creative evolution, according to music historian and pianist Michael Campbell: "'A Day in the Life' encapsulates the art and achievement of the Beatles as well as any single track can. It highlights key features of their music: the sound imagination, the persistence of tuneful melody, and the close coordination between words and music. It represents a new category of song – more sophisticated than pop ... and uniquely innovative. There literally had never before been a song – classical or vernacular – that had blended so many disparate elements so imaginatively."[391] Philosophy professor Bruce Ellis Benson agrees: "the Beatles ... give us a wonderful example of how such far-ranging influences as Celtic music, rhythm and blues, and country and western could be put together in a new way."[392]
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Author Dominic Pedler describes the way they crossed musical styles: "Far from moving sequentially from one genre to another (as is sometimes conveniently suggested) the group maintained in parallel their mastery of the traditional, catchy chart hit while simultaneously forging rock and dabbling with a wide range of peripheral influences from country to vaudeville. One of these threads was their take on folk music, which would form such essential groundwork for their later collisions with Indian music and philosophy."[393] As the personal relationships between the band members grew increasingly strained, their individual tastes became more apparent. The minimalistic cover artwork for the White Album contrasted with the complexity and diversity of its music, which encompassed Lennon's "Revolution 9" (whose musique concrète approach was influenced by Yoko Ono), Starr's country song "Don't Pass Me By", Harrison's rock ballad "While My Guitar Gently Weeps", and the "proto-metal roar" of McCartney's "Helter Skelter".[394]
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George Martin's close involvement in his role as producer made him one of the leading candidates for the informal title of the "fifth Beatle".[395] He applied his classical musical training in various ways, and functioned as "an informal music teacher" to the progressing songwriters, according to Gould.[396] Martin suggested to a sceptical McCartney that the arrangement of "Yesterday" should feature a string quartet accompaniment, thereby introducing the Beatles to a "hitherto unsuspected world of classical instrumental colour", in MacDonald's description.[397] Their creative development was also facilitated by Martin's willingness to experiment in response to their suggestions, such as adding "something baroque" to a particular recording.[398] In addition to scoring orchestral arrangements for recordings, Martin often performed on them, playing instruments including piano, organ and brass.[399]
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Collaborating with Lennon and McCartney required Martin to adapt to their different approaches to songwriting and recording. MacDonald comments, "while [he] worked more naturally with the conventionally articulate McCartney, the challenge of catering to Lennon's intuitive approach generally spurred him to his more original arrangements, of which 'Being for the Benefit of Mr. Kite!' is an outstanding example."[400] Martin said of the two composers' distinct songwriting styles and his stabilising influence:
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Compared with Paul's songs, all of which seemed to keep in some sort of touch with reality, John's had a psychedelic, almost mystical quality ... John's imagery is one of the best things about his work – 'tangerine trees', 'marmalade skies', 'cellophane flowers' ... I always saw him as an aural Salvador Dalí, rather than some drug-ridden record artist. On the other hand, I would be stupid to pretend that drugs didn't figure quite heavily in the Beatles' lives at that time ... they knew that I, in my schoolmasterly role, didn't approve ... Not only was I not into it myself, I couldn't see the need for it; and there's no doubt that, if I too had been on dope, Pepper would never have been the album it was. Perhaps it was the combination of dope and no dope that worked, who knows?[401]
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Harrison echoed Martin's description of his stabilising role: "I think we just grew through those years together, him as the straight man and us as the loonies; but he was always there for us to interpret our madness – we used to be slightly avant-garde on certain days of the week, and he would be there as the anchor person, to communicate that through the engineers and on to the tape."[402]
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Making innovative use of technology while expanding the possibilities of recorded music, the Beatles urged experimentation by Martin and his recording engineers. Seeking ways to put chance occurrences to creative use, accidental guitar feedback, a resonating glass bottle, a tape loaded the wrong way round so that it played backwards – any of these might be incorporated into their music.[403] Their desire to create new sounds on every new recording, combined with Martin's arranging abilities and the studio expertise of EMI staff engineers Norman Smith, Ken Townsend and Geoff Emerick, all contributed significantly to their records from Rubber Soul and, especially, Revolver onwards.[403]
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Along with innovative studio techniques such as sound effects, unconventional microphone placements, tape loops, double tracking and vari-speed recording, the Beatles augmented their songs with instruments that were unconventional in rock music at the time. These included string and brass ensembles as well as Indian instruments such as the sitar in "Norwegian Wood" and the swarmandal in "Strawberry Fields Forever".[404] They also used novel electronic instruments such as the Mellotron, with which McCartney supplied the flute voices on the "Strawberry Fields Forever" intro,[405] and the clavioline, an electronic keyboard that created the unusual oboe-like sound on "Baby, You're a Rich Man".[406]
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Former Rolling Stone associate editor Robert Greenfield compared the Beatles to Picasso, as "artists who broke through the constraints of their time period to come up with something that was unique and original ... [I]n the form of popular music, no one will ever be more revolutionary, more creative and more distinctive ..."[347] The British poet Philip Larkin described their work as "an enchanting and intoxicating hybrid of Negro rock-and-roll with their own adolescent romanticism", and "the first advance in popular music since the War".[408]
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They not only sparked the British Invasion of the US,[409] they became a globally influential phenomenon as well.[410] From the 1920s, the US had dominated popular entertainment culture throughout much of the world, via Hollywood films, jazz, the music of Broadway and Tin Pan Alley and, later, the rock and roll that first emerged in Memphis, Tennessee.[333] The Beatles are regarded as British cultural icons, with young adults from abroad naming the band among a group of people that they most associated with UK culture.[411][412]
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Their musical innovations and commercial success inspired musicians worldwide.[410] Many artists have acknowledged the Beatles' influence and enjoyed chart success with covers of their songs.[413] On radio, their arrival marked the beginning of a new era; in 1968 the programme director of New York's WABC radio station forbade his DJs from playing any "pre-Beatles" music, marking the defining line of what would be considered oldies on American radio.[414] They helped to redefine the album as something more than just a few hits padded out with "filler",[415] and they were primary innovators of the modern music video.[416] The Shea Stadium show with which they opened their 1965 North American tour attracted an estimated 55,600 people,[142] then the largest audience in concert history; Spitz describes the event as a "major breakthrough ... a giant step toward reshaping the concert business".[417] Emulation of their clothing and especially their hairstyles, which became a mark of rebellion, had a global impact on fashion.[100]
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According to Gould, the Beatles changed the way people listened to popular music and experienced its role in their lives. From what began as the Beatlemania fad, the group's popularity grew into what was seen as an embodiment of sociocultural movements of the decade. As icons of the 1960s counterculture, Gould continues, they became a catalyst for bohemianism and activism in various social and political arenas, fuelling movements such as women's liberation, gay liberation and environmentalism.[418] According to Peter Lavezzoli, after the "more popular than Jesus" controversy in 1966, the Beatles felt considerable pressure to say the right things and "began a concerted effort to spread a message of wisdom and higher consciousness".[167]
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Other commentators such as Mikal Gilmore and Todd Leopold have traced the inception of their socio-cultural impact earlier, interpreting even the Beatlemania period, particularly on their first visit to the US, as a key moment in the development of generational awareness.[98][419] Referring to their appearance on the Ed Sullivan Show Leopold states: "In many ways, the Sullivan appearance marked the beginning of a cultural revolution ... The Beatles were like aliens dropped into the United States of 1964."[419] According to Gilmore:
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Elvis Presley had shown us how rebellion could be fashioned into eye-opening style; the Beatles were showing us how style could have the impact of cultural revelation – or at least how a pop vision might be forged into an unimpeachable consensus.[98]
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Established in 2009, Global Beatles Day is an annual holiday on 25 June each year that honours and celebrates the ideals of the Beatles.[420] The date was chosen to commemorate the date the group participated in the BBC programme Our World in 1967, performing "All You Need Is Love" broadcast to an international audience.[421]
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In 1965, Queen Elizabeth II appointed Lennon, McCartney, Harrison and Starr Members of the Order of the British Empire (MBE).[131] The Beatles won the 1971 Academy Award for Best Original Song Score for the film Let It Be (1970).[294] The recipients of seven Grammy Awards[422] and fifteen Ivor Novello Awards,[423] the Beatles have six Diamond albums, as well as 20 Multi-Platinum albums, 16 Platinum albums and six Gold albums in the US.[305] In the UK, the Beatles have four Multi-Platinum albums, four Platinum albums, eight Gold albums and one Silver album.[306] They were inducted into the Rock and Roll Hall of Fame in 1988.[323]
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The best-selling band in history, the Beatles have sold more than 800 million physical and digital albums as of 2013[update].[424][nb 10] They have had more number-one albums on the UK charts, fifteen,[426] and sold more singles in the UK, 21.9 million, than any other act.[427] In 2004, Rolling Stone ranked the Beatles as the most significant and influential rock music artists of the last 50 years.[428] They ranked number one on Billboard magazine's list of the all-time most successful Hot 100 artists, released in 2008 to celebrate the US singles chart's 50th anniversary.[429] As of 2017[update], they hold the record for most number-one hits on the Billboard Hot 100, with twenty.[430] The Recording Industry Association of America certifies that the Beatles have sold 178 million units in the US, more than any other artist.[431] They were collectively included in Time magazine's compilation of the 20th century's 100 most influential people.[432] In 2014, they received the Grammy Lifetime Achievement Award.[433]
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On 16 January each year, beginning in 2001, people celebrate World Beatles Day under UNESCO. This date has direct relation to the opening of The Cavern Club in 1957.[434][435]
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Five asteroids, 4147 Lennon, 4148 McCartney, 4149 Harrison, 4150 Starr and 8749 Beatles are named after the Beatles. In 2007, the Beatles became the first band to feature on a series of UK postage stamps issued by the Royal Mail.[436]
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Principal members
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Early members
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Touring musician
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The Beatles have a core catalogue consisting of 13 studio albums and one compilation.[437]
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Core catalogue
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See also
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Through 1969, the Beatles' catalogue was published almost exclusively by Northern Songs Ltd, a company formed in February 1963 by music publisher Dick James specifically for Lennon and McCartney, though it later acquired songs by other artists. The company was organised with James and his partner, Emmanuel Silver, owning a controlling interest, variously described as 51% or 50% plus one share. McCartney had 20%. Reports again vary concerning Lennon's portion – 19 or 20% – and Brian Epstein's – 9 or 10% – which he received in lieu of a 25% band management fee.[438][439][440] In 1965, the company went public. Five million shares were created, of which the original principals retained 3.75 million. James and Silver each received 937,500 shares (18.75% of 5 million); Lennon and McCartney each received 750,000 shares (15%); and Epstein's management company, NEMS Enterprises, received 375,000 shares (7.5%). Of the 1.25 million shares put up for sale, Harrison and Starr each acquired 40,000.[441] At the time of the stock offering, Lennon and McCartney renewed their three-year publishing contracts, binding them to Northern Songs until 1973.[442]
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Harrison created Harrisongs to represent his Beatles compositions, but signed a three-year contract with Northern Songs that gave it the copyright to his work through March 1968, which included "Taxman" and "Within You Without You".[443] The songs on which Starr received co-writing credit before 1968, such as "What Goes On" and "Flying", were also Northern Songs copyrights.[444] Harrison did not renew his contract with Northern Songs when it ended, signing instead with Apple Publishing while retaining the copyright to his work from that point on. Harrison thus owns the rights to his later Beatles songs such as "While My Guitar Gently Weeps" and "Something". That year, as well, Starr created Startling Music, which holds the rights to his Beatles compositions, "Don't Pass Me By" and "Octopus's Garden".[445][446]
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In March 1969, James arranged to sell his and his partner's shares of Northern Songs to the British broadcasting company Associated Television (ATV), founded by impresario Lew Grade, without first informing the Beatles. The band then made a bid to gain a controlling interest by attempting to work out a deal with a consortium of London brokerage firms that had accumulated a 14% holding.[447] The deal collapsed over the objections of Lennon, who declared, "I'm sick of being fucked about by men in suits sitting on their fat arses in the City."[448] By the end of May, ATV had acquired a majority stake in Northern Songs, controlling nearly the entire Lennon–McCartney catalogue, as well as any future material until 1973.[449] In frustration, Lennon and McCartney sold their shares to ATV in late October 1969.[450]
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In 1981, financial losses by ATV's parent company, Associated Communications Corporation (ACC), led it to attempt to sell its music division. According to authors Brian Southall and Rupert Perry, Grade contacted McCartney, offering ATV Music and Northern Songs for $30 million.[451] According to an account McCartney gave in 1995, he met with Grade and explained he was interested solely in the Northern Songs catalogue if Grade were ever willing to "separate off" that portion of ATV Music. Soon afterwards, Grade offered to sell him Northern Songs for £20 million, giving the ex-Beatle "a week or so" to decide. By McCartney's account, he and Ono countered with a £5 million bid that was rejected.[452] According to reports at the time, Grade refused to separate Northern Songs and turned down an offer of £21–25 million from McCartney and Ono for Northern Songs. In 1982, ACC was acquired in a takeover by Australian business magnate Robert Holmes à Court for £60 million.[453]
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In 1985, Michael Jackson purchased ATV for a reported $47.5 million. The acquisition gave him control over the publishing rights to more than 200 Beatles songs, as well as 40,000 other copyrights.[454] In 1995, in a deal that earned him a reported $110 million, Jackson merged his music publishing business with Sony, creating a new company, Sony/ATV Music Publishing, in which he held a 50% stake. The merger made the new company, then valued at over half a billion dollars, the third-largest music publisher in the world.[455] In 2016, Sony acquired Jackson's share of Sony/ATV from the Jackson estate for $750 million.[456]
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Despite the lack of publishing rights to most of their songs, Lennon's estate and McCartney continue to receive their respective shares of the writers' royalties, which together are 331⁄3% of total commercial proceeds in the US and which vary elsewhere around the world between 50 and 55%.[457] Two of Lennon and McCartney's earliest songs – "Love Me Do" and "P.S. I Love You" – were published by an EMI subsidiary, Ardmore & Beechwood, before they signed with James. McCartney acquired their publishing rights from Ardmore[458] in 1978,[459] and they are the only two Beatles songs owned by McCartney's company MPL Communications.[460] On 18 January 2017, McCartney filed a suit in the United States district court against Sony/ATV Music Publishing seeking to reclaim ownership of his share of the Lennon–McCartney song catalogue beginning in 2018. Under US copyright law, for works published before 1978 the author can reclaim copyrights assigned to a publisher after 56 years.[461][462] McCartney and Sony agreed to a confidential settlement in June 2017.[463][464]
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Fictionalised
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Documentaries and filmed performances
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1963
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1964
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1965
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1966
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