--- license: cc-by-nc-4.0 library_name: diffusers base_model: stabilityai/stable-diffusion-3-medium tags: - lora - text-to-image inference: False --- # ⚡ Flash Diffusion: FlashSD3 ⚡ Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.* This model is a **90.4M** LoRA distilled version of [SD3](https://huggingface.co/stabilityai/stable-diffusion-3-medium) model that is able to generate 1024x1024 images in **4 steps**. See our [live demo](https://huggingface.co/spaces/jasperai/flash-sd3) and official [Github repo](https://github.com/gojasper/flash-diffusion).

# How to use? The model can be used using the `StableDiffusion3Pipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. First, you need to install a specific version of `diffusers` by runniung ```bash pip install git+https://github.com/initml/diffusers.git@clement/feature/flash_sd3 ``` Then, you can ru the following to generate an image ```python import torch from diffusers import StableDiffusion3Pipeline, SD3Transformer2DModel, FlashFlowMatchEulerDiscreteScheduler from peft import PeftModel # Load LoRA transformer = SD3Transformer2DModel.from_pretrained( "stabilityai/stable-diffusion-3-medium-diffusers", subfolder="transformer", torch_dtype=torch.float16, ) transformer = PeftModel.from_pretrained(transformer, "jasperai/flash-sd3") # Pipeline pipe = StableDiffusion3Pipeline.from_pretrained( "stabilityai/stable-diffusion-3-medium-diffusers", transformer=transformer, torch_dtype=torch.float16, text_encoder_3=None, tokenizer_3=None ) # Scheduler pipe.scheduler = FlashFlowMatchEulerDiscreteScheduler.from_pretrained( "stabilityai/stable-diffusion-3-medium-diffusers", subfolder="scheduler", ) pipe.to("cuda") prompt = "A raccoon trapped inside a glass jar full of colorful candies, the background is steamy with vivid colors." image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] ```

## Citation If you find this work useful or use it in your research, please consider citing us ```bibtex @misc{chadebec2024flash, title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation}, author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin}, year={2024}, eprint={2406.02347}, archivePrefix={arXiv}, primaryClass={cs.CV} } ``` ## License This model is released under the the Creative Commons BY-NC license.