--- base_model: stabilityai/stable-diffusion-xl-base-1.0 library_name: diffusers license: openrail++ tags: - text-to-image - text-to-image - diffusers-training - diffusers - lora - template:sd-lora - stable-diffusion-xl - stable-diffusion-xl-diffusers instance_prompt: a photo of sks dog widget: - text: A photo of sks dog in a bucket output: url: image_0.png - text: A photo of sks dog in a bucket output: url: image_1.png - text: A photo of sks dog in a bucket output: url: image_2.png - text: A photo of sks dog in a bucket output: url: image_3.png --- # SDXL LoRA DreamBooth - preetham/labra_model2 ## Model description These are preetham/labra_model2 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. The weights were trained using [DreamBooth](https://dreambooth.github.io/). LoRA for the text encoder was enabled: False. Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. ## Trigger words You should use a photo of sks dog to trigger the image generation. ## Download model Weights for this model are available in Safetensors format. [Download](preetham/labra_model2/tree/main) them in the Files & versions tab. ## Intended uses & limitations #### How to use ```python # TODO: add an example code snippet for running this diffusion pipeline ``` #### Limitations and bias [TODO: provide examples of latent issues and potential remediations] ## Training details [TODO: describe the data used to train the model]