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import torch | |
from PIL import Image | |
import argparse | |
import os, json, random | |
import pandas as pd | |
import matplotlib.pyplot as plt | |
import glob, re | |
from safetensors.torch import load_file | |
import matplotlib.image as mpimg | |
import copy | |
import gc | |
from transformers import CLIPTextModel, CLIPTokenizer | |
import diffusers | |
from diffusers import DiffusionPipeline | |
from diffusers import AutoencoderKL, DDPMScheduler, DiffusionPipeline, UNet2DConditionModel, LMSDiscreteScheduler | |
from diffusers.loaders import AttnProcsLayers | |
from diffusers.models.attention_processor import LoRAAttnProcessor, AttentionProcessor | |
from typing import Any, Dict, List, Optional, Tuple, Union | |
from lora import LoRANetwork, DEFAULT_TARGET_REPLACE, UNET_TARGET_REPLACE_MODULE_CONV | |
from diffusers.pipelines.stable_diffusion_xl import StableDiffusionXLPipelineOutput | |
import inspect | |
import os | |
from typing import Any, Callable, Dict, List, Optional, Tuple, Union | |
from diffusers.pipelines import StableDiffusionXLPipeline | |
import random | |
import torch | |
from transformers import CLIPTextModel, CLIPTextModelWithProjection, CLIPTokenizer | |
import re | |
import argparse | |
def flush(): | |
torch.cuda.empty_cache() | |
gc.collect() | |
def call( | |
self, | |
prompt: Union[str, List[str]] = None, | |
prompt_2: Optional[Union[str, List[str]]] = None, | |
height: Optional[int] = None, | |
width: Optional[int] = None, | |
num_inference_steps: int = 50, | |
denoising_end: Optional[float] = None, | |
guidance_scale: float = 5.0, | |
negative_prompt: Optional[Union[str, List[str]]] = None, | |
negative_prompt_2: Optional[Union[str, List[str]]] = None, | |
num_images_per_prompt: Optional[int] = 1, | |
eta: float = 0.0, | |
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, | |
latents: Optional[torch.FloatTensor] = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
output_type: Optional[str] = "pil", | |
return_dict: bool = True, | |
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None, | |
callback_steps: int = 1, | |
cross_attention_kwargs: Optional[Dict[str, Any]] = None, | |
guidance_rescale: float = 0.0, | |
original_size: Optional[Tuple[int, int]] = None, | |
crops_coords_top_left: Tuple[int, int] = (0, 0), | |
target_size: Optional[Tuple[int, int]] = None, | |
negative_original_size: Optional[Tuple[int, int]] = None, | |
negative_crops_coords_top_left: Tuple[int, int] = (0, 0), | |
negative_target_size: Optional[Tuple[int, int]] = None, | |
network=None, | |
start_noise=None, | |
scale=None, | |
unet=None, | |
): | |
r""" | |
Function invoked when calling the pipeline for generation. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. | |
instead. | |
prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is | |
used in both text-encoders | |
height (`int`, *optional*, defaults to unet.config.sample_size * self.vae_scale_factor): | |
The height in pixels of the generated image. This is set to 1024 by default for the best results. | |
Anything below 512 pixels won't work well for | |
[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | |
and checkpoints that are not specifically fine-tuned on low resolutions. | |
width (`int`, *optional*, defaults to unet.config.sample_size * self.vae_scale_factor): | |
The width in pixels of the generated image. This is set to 1024 by default for the best results. | |
Anything below 512 pixels won't work well for | |
[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | |
and checkpoints that are not specifically fine-tuned on low resolutions. | |
num_inference_steps (`int`, *optional*, defaults to 50): | |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the | |
expense of slower inference. | |
denoising_end (`float`, *optional*): | |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be | |
completed before it is intentionally prematurely terminated. As a result, the returned sample will | |
still retain a substantial amount of noise as determined by the discrete timesteps selected by the | |
scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a | |
"Mixture of Denoisers" multi-pipeline setup, as elaborated in [**Refining the Image | |
Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output) | |
guidance_scale (`float`, *optional*, defaults to 5.0): | |
Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598). | |
`guidance_scale` is defined as `w` of equation 2. of [Imagen | |
Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale > | |
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, | |
usually at the expense of lower image quality. | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
less than `1`). | |
negative_prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and | |
`text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders | |
num_images_per_prompt (`int`, *optional*, defaults to 1): | |
The number of images to generate per prompt. | |
eta (`float`, *optional*, defaults to 0.0): | |
Corresponds to parameter eta (Ξ·) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to | |
[`schedulers.DDIMScheduler`], will be ignored for others. | |
generator (`torch.Generator` or `List[torch.Generator]`, *optional*): | |
One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html) | |
to make generation deterministic. | |
latents (`torch.FloatTensor`, *optional*): | |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image | |
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents | |
tensor will ge generated by sampling using the supplied random `generator`. | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
provided, text embeddings will be generated from `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
argument. | |
pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. | |
If not provided, pooled text embeddings will be generated from `prompt` input argument. | |
negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` | |
input argument. | |
output_type (`str`, *optional*, defaults to `"pil"`): | |
The output format of the generate image. Choose between | |
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`. | |
return_dict (`bool`, *optional*, defaults to `True`): | |
Whether or not to return a [`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] instead | |
of a plain tuple. | |
callback (`Callable`, *optional*): | |
A function that will be called every `callback_steps` steps during inference. The function will be | |
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`. | |
callback_steps (`int`, *optional*, defaults to 1): | |
The frequency at which the `callback` function will be called. If not specified, the callback will be | |
called at every step. | |
cross_attention_kwargs (`dict`, *optional*): | |
A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under | |
`self.processor` in | |
[diffusers.models.attention_processor](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py). | |
guidance_rescale (`float`, *optional*, defaults to 0.7): | |
Guidance rescale factor proposed by [Common Diffusion Noise Schedules and Sample Steps are | |
Flawed](https://arxiv.org/pdf/2305.08891.pdf) `guidance_scale` is defined as `Ο` in equation 16. of | |
[Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf). | |
Guidance rescale factor should fix overexposure when using zero terminal SNR. | |
original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
If `original_size` is not the same as `target_size` the image will appear to be down- or upsampled. | |
`original_size` defaults to `(width, height)` if not specified. Part of SDXL's micro-conditioning as | |
explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): | |
`crops_coords_top_left` can be used to generate an image that appears to be "cropped" from the position | |
`crops_coords_top_left` downwards. Favorable, well-centered images are usually achieved by setting | |
`crops_coords_top_left` to (0, 0). Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
For most cases, `target_size` should be set to the desired height and width of the generated image. If | |
not specified it will default to `(width, height)`. Part of SDXL's micro-conditioning as explained in | |
section 2.2 of [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
negative_original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's | |
micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
negative_crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): | |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's | |
micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
negative_target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
To negatively condition the generation process based on a target image resolution. It should be as same | |
as the `target_size` for most cases. Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
Examples: | |
Returns: | |
[`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] or `tuple`: | |
[`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput`] if `return_dict` is True, otherwise a | |
`tuple`. When returning a tuple, the first element is a list with the generated images. | |
""" | |
# 0. Default height and width to unet | |
height = height or self.default_sample_size * self.vae_scale_factor | |
width = width or self.default_sample_size * self.vae_scale_factor | |
original_size = original_size or (height, width) | |
target_size = target_size or (height, width) | |
# 1. Check inputs. Raise error if not correct | |
self.check_inputs( | |
prompt, | |
prompt_2, | |
height, | |
width, | |
callback_steps, | |
negative_prompt, | |
negative_prompt_2, | |
prompt_embeds, | |
negative_prompt_embeds, | |
pooled_prompt_embeds, | |
negative_pooled_prompt_embeds, | |
) | |
# 2. Define call parameters | |
if prompt is not None and isinstance(prompt, str): | |
batch_size = 1 | |
elif prompt is not None and isinstance(prompt, list): | |
batch_size = len(prompt) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
device = self._execution_device | |
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2) | |
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1` | |
# corresponds to doing no classifier free guidance. | |
do_classifier_free_guidance = guidance_scale > 1.0 | |
# 3. Encode input prompt | |
text_encoder_lora_scale = ( | |
cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None | |
) | |
( | |
prompt_embeds, | |
negative_prompt_embeds, | |
pooled_prompt_embeds, | |
negative_pooled_prompt_embeds, | |
) = self.encode_prompt( | |
prompt=prompt, | |
prompt_2=prompt_2, | |
device=device, | |
num_images_per_prompt=num_images_per_prompt, | |
do_classifier_free_guidance=do_classifier_free_guidance, | |
negative_prompt=negative_prompt, | |
negative_prompt_2=negative_prompt_2, | |
prompt_embeds=prompt_embeds, | |
negative_prompt_embeds=negative_prompt_embeds, | |
pooled_prompt_embeds=pooled_prompt_embeds, | |
negative_pooled_prompt_embeds=negative_pooled_prompt_embeds, | |
lora_scale=text_encoder_lora_scale, | |
) | |
# 4. Prepare timesteps | |
self.scheduler.set_timesteps(num_inference_steps, device=device) | |
timesteps = self.scheduler.timesteps | |
# 5. Prepare latent variables | |
num_channels_latents = unet.config.in_channels | |
latents = self.prepare_latents( | |
batch_size * num_images_per_prompt, | |
num_channels_latents, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
latents, | |
) | |
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline | |
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) | |
# 7. Prepare added time ids & embeddings | |
add_text_embeds = pooled_prompt_embeds | |
add_time_ids = self._get_add_time_ids( | |
original_size, crops_coords_top_left, target_size, dtype=prompt_embeds.dtype | |
) | |
if negative_original_size is not None and negative_target_size is not None: | |
negative_add_time_ids = self._get_add_time_ids( | |
negative_original_size, | |
negative_crops_coords_top_left, | |
negative_target_size, | |
dtype=prompt_embeds.dtype, | |
) | |
else: | |
negative_add_time_ids = add_time_ids | |
if do_classifier_free_guidance: | |
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds], dim=0) | |
add_text_embeds = torch.cat([negative_pooled_prompt_embeds, add_text_embeds], dim=0) | |
add_time_ids = torch.cat([negative_add_time_ids, add_time_ids], dim=0) | |
prompt_embeds = prompt_embeds.to(device) | |
add_text_embeds = add_text_embeds.to(device) | |
add_time_ids = add_time_ids.to(device).repeat(batch_size * num_images_per_prompt, 1) | |
# 8. Denoising loop | |
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0) | |
# 7.1 Apply denoising_end | |
if denoising_end is not None and isinstance(denoising_end, float) and denoising_end > 0 and denoising_end < 1: | |
discrete_timestep_cutoff = int( | |
round( | |
self.scheduler.config.num_train_timesteps | |
- (denoising_end * self.scheduler.config.num_train_timesteps) | |
) | |
) | |
num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps))) | |
timesteps = timesteps[:num_inference_steps] | |
with self.progress_bar(total=num_inference_steps) as progress_bar: | |
for i, t in enumerate(timesteps): | |
if t>start_noise: | |
network.set_lora_slider(scale=0) | |
else: | |
network.set_lora_slider(scale=scale) | |
# expand the latents if we are doing classifier free guidance | |
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents | |
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) | |
# predict the noise residual | |
added_cond_kwargs = {"text_embeds": add_text_embeds, "time_ids": add_time_ids} | |
with network: | |
noise_pred = unet( | |
latent_model_input, | |
t, | |
encoder_hidden_states=prompt_embeds, | |
cross_attention_kwargs=cross_attention_kwargs, | |
added_cond_kwargs=added_cond_kwargs, | |
return_dict=False, | |
)[0] | |
# perform guidance | |
if do_classifier_free_guidance: | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
if do_classifier_free_guidance and guidance_rescale > 0.0: | |
# Based on 3.4. in https://arxiv.org/pdf/2305.08891.pdf | |
noise_pred = rescale_noise_cfg(noise_pred, noise_pred_text, guidance_rescale=guidance_rescale) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] | |
# call the callback, if provided | |
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): | |
progress_bar.update() | |
if callback is not None and i % callback_steps == 0: | |
callback(i, t, latents) | |
if not output_type == "latent": | |
# make sure the VAE is in float32 mode, as it overflows in float16 | |
needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast | |
if needs_upcasting: | |
self.upcast_vae() | |
latents = latents.to(next(iter(self.vae.post_quant_conv.parameters())).dtype) | |
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0] | |
# cast back to fp16 if needed | |
if needs_upcasting: | |
self.vae.to(dtype=torch.float16) | |
else: | |
image = latents | |
if not output_type == "latent": | |
# apply watermark if available | |
if self.watermark is not None: | |
image = self.watermark.apply_watermark(image) | |
image = self.image_processor.postprocess(image, output_type=output_type) | |
# Offload all models | |
# self.maybe_free_model_hooks() | |
if not return_dict: | |
return (image,) | |
return StableDiffusionXLPipelineOutput(images=image) | |
def sorted_nicely( l ): | |
convert = lambda text: float(text) if text.replace('-','').replace('.','').isdigit() else text | |
alphanum_key = lambda key: [convert(c) for c in re.split('(-?[0-9]+.?[0-9]+?)', key) ] | |
return sorted(l, key = alphanum_key) | |
def flush(): | |
torch.cuda.empty_cache() | |
gc.collect() | |
if __name__=='__main__': | |
device = 'cuda:0' | |
StableDiffusionXLPipeline.__call__ = call | |
pipe = StableDiffusionXLPipeline.from_pretrained('stabilityai/stable-diffusion-xl-base-1.0') | |
# pipe.__call__ = call | |
pipe = pipe.to(device) | |
parser = argparse.ArgumentParser( | |
prog = 'generateImages', | |
description = 'Generate Images using Diffusers Code') | |
parser.add_argument('--model_name', help='name of model', type=str, required=True) | |
parser.add_argument('--prompts_path', help='path to csv file with prompts', type=str, required=True) | |
parser.add_argument('--negative_prompts', help='negative prompt', type=str, required=False, default=None) | |
parser.add_argument('--save_path', help='folder where to save images', type=str, required=True) | |
parser.add_argument('--base', help='version of stable diffusion to use', type=str, required=False, default='1.4') | |
parser.add_argument('--guidance_scale', help='guidance to run eval', type=float, required=False, default=7.5) | |
parser.add_argument('--image_size', help='image size used to train', type=int, required=False, default=512) | |
parser.add_argument('--till_case', help='continue generating from case_number', type=int, required=False, default=1000000) | |
parser.add_argument('--from_case', help='continue generating from case_number', type=int, required=False, default=0) | |
parser.add_argument('--num_samples', help='number of samples per prompt', type=int, required=False, default=5) | |
parser.add_argument('--ddim_steps', help='ddim steps of inference used to train', type=int, required=False, default=50) | |
parser.add_argument('--rank', help='rank of the LoRA', type=int, required=False, default=4) | |
parser.add_argument('--start_noise', help='what time stamp to flip to edited model', type=int, required=False, default=750) | |
args = parser.parse_args() | |
lora_weight = args.model_name | |
csv_path = args.prompts_path | |
save_path = args.save_path | |
start_noise = args.start_noise | |
from_case = args.from_case | |
till_case = args.till_case | |
weight_dtype = torch.float16 | |
num_images_per_prompt = 1 | |
scales = [-2, -1, 0, 1, 2] | |
scales = [-1, -.5, 0, .5, 1] | |
scales = [-2] | |
df = pd.read_csv(csv_path) | |
for scale in scales: | |
os.makedirs(f'{save_path}/{os.path.basename(lora_weight)}/{scale}', exist_ok=True) | |
prompts = list(df['prompt']) | |
seeds = list(df['evaluation_seed']) | |
case_numbers = list(df['case_number']) | |
pipe = StableDiffusionXLPipeline.from_pretrained('stabilityai/stable-diffusion-xl-base-1.0',torch_dtype=torch.float16,) | |
# pipe.__call__ = call | |
pipe = pipe.to(device) | |
unet = pipe.unet | |
if 'full' in lora_weight: | |
train_method = 'full' | |
elif 'noxattn' in lora_weight: | |
train_method = 'noxattn' | |
else: | |
train_method = 'noxattn' | |
network_type = "c3lier" | |
if train_method == 'xattn': | |
network_type = 'lierla' | |
modules = DEFAULT_TARGET_REPLACE | |
if network_type == "c3lier": | |
modules += UNET_TARGET_REPLACE_MODULE_CONV | |
import os | |
model_name = lora_weight | |
name = os.path.basename(model_name) | |
rank = 1 | |
alpha = 4 | |
if 'rank4' in lora_weight: | |
rank = 4 | |
if 'rank8' in lora_weight: | |
rank = 8 | |
if 'alpha1' in lora_weight: | |
alpha = 1.0 | |
network = LoRANetwork( | |
unet, | |
rank=rank, | |
multiplier=1.0, | |
alpha=alpha, | |
train_method=train_method, | |
).to(device, dtype=weight_dtype) | |
network.load_state_dict(torch.load(lora_weight)) | |
for idx, prompt in enumerate(prompts): | |
seed = seeds[idx] | |
case_number = case_numbers[idx] | |
if not (case_number>=from_case and case_number<=till_case): | |
continue | |
if os.path.exists(f'{save_path}/{os.path.basename(lora_weight)}/{scale}/{case_number}_{idx}.png'): | |
continue | |
print(prompt, seed) | |
for scale in scales: | |
generator = torch.manual_seed(seed) | |
images = pipe(prompt, num_images_per_prompt=args.num_samples, num_inference_steps=50, generator=generator, network=network, start_noise=start_noise, scale=scale, unet=unet).images | |
for idx, im in enumerate(images): | |
im.save(f'{save_path}/{os.path.basename(lora_weight)}/{scale}/{case_number}_{idx}.png') | |
del unet, network, pipe | |
unet = None | |
network = None | |
pipe = None | |
torch.cuda.empty_cache() | |
flush() |