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0704.0313 | Possibility of Gapless Spin Liquid State by One-dimensionalization | arXiv:0704.0313v1 [cond-mat.str-el] 3 Apr 2007
Typeset with jpsj2.cls <ver.1.2> Letter
Possibility of Gapless Spin Liquid State by One-dimensionalization
Yuta Hayashi∗ and Masao Ogata
Department of Physics, University of Tokyo, Hongo, Bunkyo-ku, Tokyo, 113-0033
Motivated by the observation of a gapless spin liquid state in κ-(BEDT-TTF)2Cu2(CN)3, we
analyze the anisotropic triangular lattice S = 1/2 Heisenberg model with the resonating valence
bond mean-field approximation. Paying attention to the small quasi-one-dimensional anisotropy
of the material, we take an approach from one-dimensional (1D) chains coupled with frustrating
zig-zag bonds. By calculating one-particle excitation spectra changing anisotropy parameter
J ′/J from the decoupled 1D chains to the isotropic triangular lattice, we find almost gapless
excitations in the wide range from the 1D limit. This one-dimensionalization by frustration is
considered to be a candidate for the mechanism of the gapless spin liquid state.
KEYWORDS: gapless spin liquid, κ-(BEDT-TTF)2Cu2(CN)3, anisotropic triangular lattice, frustration,
one-dimensionalization
Organic conductors are one of the fascinating materi-
als which have low-dimensionality and relatively strong
electron correlations. So far, various physical states
have been observed and investigated intensively.1 Among
them, magnetism in the Mott insulating phase next to
the unconventional superconductivity has been attract-
ing considerable attention. This phase is observed in the
family of κ-(BEDT-TTF)2X, where BEDT-TTF (ET)
denotes bis(ethylenedithio)-tetrathiafulvalene and X rep-
resents a monovalent anion. Similarities to that of high-
Tc cuprates are worthy of note. Another stimulating
problem concerning magnetism is ground state proper-
ties of geometrically frustrated spin systems such as a tri-
angular lattice and a Kagomé lattice. These two intrigu-
ing issues meet in a material κ-(ET)2Cu2(CN)3, which
is a Mott insulator having a nearly isotropic triangular
lattice, and it has been in the spotlight of late.
According to 1H NMR measurements at ambient pres-
sure,2 κ-(ET)2Cu2(CN)3 shows no indication of long-
range magnetic order (LRMO) down to 32mK. This
is 4 orders of magnitude below the exchange constant
J ∼ 250K estimated from the temperature dependence
of susceptibility. Recently, a similar result has been ob-
tained by zero-field muon spin relaxation measurements,
which have observed no LRMO down to 20mK.3 These
results suggest that a quantum spin liquid state is real-
ized in the ground state. On the other hand, the static
susceptibility remains finite down to 1.9K, and spin-
lattice relaxation rate 1/T1 shows power-law temperature
dependence below 1K. These imply that almost gapless
spin excitation exists. This fact is a significant feature of
the spin liquid phase observed in this material.
Since Anderson’s proposal of a resonating valence
bond (RVB) state,5 enormous number of studies have
been made on the triangular lattice spin system. It is
now a general view that the ground state of the isotropic
triangular lattice Heisenberg model has LRMO, such as
the 120◦ structure.6–9 On the other hand, if one ne-
glects the LRMO and assumes a disordered ground state,
the mean-field theory of RVB state gives a spin-gap
∗E-mail address: [email protected]
Table I. Anisotropy of effective transfer integrals in κ-(ET)2X.
The definition of t and t′ are not as usual (see the text).
Anion X t′/t
Cu2(CN)3 0.94
Cu(NCS)2 1.19
Cu[N(CN)2]Br 1.33
Cu[N(CN)2]Cl 1.47
Cu(CN)[N(CN)2] 1.47
Ag(CN)2·H2O 1.67
I3 1.72
state with dx2−y2+idxy-wave symmetry, which is called
“d+id state”.10–12 This RVB state, describing an insu-
lating spin system, corresponds to a projected BCS state
at half-filling in which doubly occupied states are ex-
cluded. Thus, the existing theories show that the ground
state has LRMO in general, and if the magnetic order
is destroyed in some reason, the d+id fullgap state will
appear. If we regard the Mott insulating phase of κ-
(ET)2Cu2(CN)3 in low temperatures as an isotropic tri-
angular lattice spin system, the results of NMR and sus-
ceptibility measurements, which suggest neither LRMO
nor spin gap, cannot be explained.
In this letter, we pay attention to small anisotropy of
κ-(ET)2Cu2(CN)3 and propose a new possibility for un-
derstanding its gapless spin liquid state. As shown in Ta-
ble I, only κ-(ET)2Cu2(CN)3 has an opposite anisotropy
among the family of κ-(ET)2X studied in the past. Here,
the effective transfer integrals t and t′ are defined in-
versely to the conventional way; t = 0 corresponds to
the square lattice, and t′ = 0 the decoupled chains.
Therefore, κ-(ET)2Cu2(CN)3 has quasi-one-dimensional
(Q1D) anisotropy rather than an isotropic triangular lat-
tice. Considering that the pure 1D spin system has no
LRMO and gapless spin excitation, it is likely that this
Q1D anisotropy is concerned with the formation of the
gapless spin liquid state in κ-(ET)2Cu2(CN)3.
Based on the above consideration, we study the
Heisenberg model on an anisotropic triangular lattice,
which is equivalent to 1D chains coupled with zig-zag
http://arxiv.org/abs/0704.0313v1
2 J. Phys. Soc. Jpn. Letter Author Name
bonds as shown in Fig. 1. The Hamiltonian is given by
<i,i′>
JSi · Si′ +
<i,j>
J ′Si · Sj , (1)
where <i, i′> and <i, j> represent the summation over
intrachain and interchain nearest-neighbor pairs with an-
tiferromagnetic coupling constant J and J ′, respectively
(see Fig. 1). We investigate the anisotropy parameter
range J ′/J = 0.0-1.0, in which the model interpolates
between the decoupled chains (J ′ = 0) and the isotropic
triangular lattice (J ′ = J).
In the following, we consider a projected BCS state
defined as
∣p-BCS
, (2)
where PG is the Gutzwiller projection operator which ex-
cludes double occupancy and
is a BCS mean-field
wave function. Since it is difficult to treat the Gutzwiller
projection analytically, we apply an RVB mean-field ap-
proximation to the Hamiltonian (1) and calculate the
one-particle excitation spectra. To put it more con-
cretely, we introduce mean fields ∆ij ≡
ci↑cj↓
, ξij ≡
and obtain its excitation spectrum by
diagonalizing the mean-field Hamiltonian. This approxi-
mation is equivalent to the “Gutzwiller approximation”
which replaces the effect of the Gutzwiller projection op-
erator with the statistical weight gs as
p-BCS
∣Si ·Sj
∣p-BCS
∣Si ·Sj
. (3)
In the simplest Gutzwiller approximation, the statisti-
cal weight is given as gs = 4/(1 + δ)
2 where δ is the
density of holes,15 and in the case of half-filling (δ = 0),
gs = 4. Although double occupancy is no longer excluded
from wave functions in this approximation, it is known
in the research of high-Tc superconductivity that the
RVB mean-field (Gutzwiller) approximation gives quali-
tatively good results.
The spin operators Si ·Sj in the Hamiltonian (1) can
be rewritten by the fermion operators as
Si · Sj =
ci↑ − c†i↓ci↓
cj↑ − c†j↓cj↓
cj↑ + c
. (4)
Fig. 1. The anisotropic triangular lattice Heisenberg model with
intrachain coupling J and interchain zig-zag coupling J ′.
τ1, τ2, τ3 are lattice vectors.
By introducing the mean fields, we can rewrite the
Hamiltonian as
HMF =
ck↑+ c
+ h.c.
except for constant terms. Here, ξk and ∆k are given by
ξk ≡ −3Jξτ1cos(k · τ 1)
− 3J ′
ξτ 2cos(k · τ 2) + ξτ 3cos(k · τ 3)
, (6)
∆k ≡ 3J∆τ1cos(k · τ 1)
+ 3J ′
∆τ 2cos(k · τ 2) + ∆τ3cos(k · τ 3)
, (7)
where τ 1 = (1, 0), τ 2 = (1/2,
3/2), τ 3 = (1/2,−
as shown in Fig. 1, and
ci+τ↑
ci+τ↓
, ∆τ ≡
ci↑ci+τ↓
. (8)
On the analogy of BCS theory, we obtain self-consistent
equations at zero temperature
ξτ i = −
eik·τ i
∆τ i =
e−ik·τ i
with a quasiparticle excitation spectrum
+ |∆k|2. (10)
We determine the order parameters ∆τ i , ξτ i (i = 1, 2, 3)
by solving self-consistent equations (9) numerically, and
obtain the one-particle excitation spectrum Ek.
Firstly, we verify our method in 1D limit (J ′/J = 0).
According to the exact solution, the ground state is a
spin disordered state and the excitation spectrum is “des
Cloizeaux-Pearson mode” with S = 1.16 In the present
RVB mean-field theory, the one-particle excitation spec-
trum becomes
Ek = 3J
+ |∆τ 1 |
2 |cos kx| (11)
in the 1D limit. This clearly realizes gapless excitations
at kx = ±π/2. Note that this one-particle excitation
describes a spin singlet breaking, i.e. S = 1/2 spinon
excitation, whereas the des Cloizeaux-Pearson mode de-
scribes S = 1 spin-wave (magnon) excitation. Thus, two-
spinon excitations with kx = π/2 and kx = −π/2 form an
S = 1 magnon with kx = 0. This means that the present
gapless excitation spectrum obtained in the RVB mean-
field theory is consistent with the exact des Cloizeaux-
Pearson mode.
Nextly, we show the results of 0 ≤ J ′/J ≤ 1 case,
focusing on the following parameters
+ |∆τ 1 |
D23 ≡
+ |∆τ2 |
+ |∆τ3 |
Because of the SU(2) degeneracy at half-filling,10, 15 these
parameters are determined uniquely regardless of the de-
generate ground states. Actually, the excitation spectrum
J. Phys. Soc. Jpn. Letter Author Name 3
can be written as
= 9J2D21 cos
+ 9J ′2D223
+ cos2
Therefore, D1, D23 determine the dispersion relations
along the chains (τ 1) and between the chains (τ 2,τ 3),
respectively. Their J ′/J dependence calculated in the
system size L = 1200 (N = L2) are plotted in Fig. 2.
A notable feature is that D23 remains very small com-
pared to D1, in spite of the comparatively large J
to J ′/J ∼ 0.25. When D23 = 0 the system is a pure 1D
chain. Indeed, when J ′/J = 0, the right-hand side of the
self-consistent equations of ξτ 2 , ξτ3 , ∆τ2 , ∆τ 3 become
all equal to zero. As we show later, D23 is very small
for J ′/J . 0.25 and vanishes when J ′/J → 0. This in-
dicates that there are scarcely any correlations between
spins of different chains, and practically 1D state is real-
ized. As J ′/J approaches unity, D23 gradually increases
and becomes equal to D1.
Finally, we show in Fig. 3 the J ′/J dependence of the
one-particle excitation spectra Ek in (12). We find that
the structure of excitation spectra in 0 ≤ J ′/J . 0.25
has little difference from that of the decoupled chains
(J ′/J = 0.0). As a result, almost gapless excitations are
realized in this wide parameter range. This means that
practically 1D state is realized, which is also expected
from the behavior of D23 in Fig. 2. When J
′/J exceeds
0.25, the excitation gap gradually increases globally in
the first Brillouin zone (1BZ). However, the shape of the
whole spectrum is almost unchanged until the J ′/J be-
comes as large as about 0.6. Moreover, focusing on the
lowest energy excitations (dark areas in the contour plot
shown in Fig. 3), their locations in the 1BZ do not deviate
from those in the 1D limit (kx = ±π/2) for J ′/J . 0.8.
Additionally, when kx = ±π/2, the excitation spectrum
Ek is independent of ky, i.e., Ek = 3J
′D23. This is be-
cause the frustration of two interchain couplings (corre-
sponding to the lattice vector τ 2 and τ 3) cancel the ky
dependence. This fact is rather important, since it indi-
cates that the excited quasiparticles along the kx = ±π/2
lines feel free to move along the ky direction. This is the
same condition as in the 1D limit, except for the exis-
Fig. 2. Anisotropy dependence of D1 and D23 for L = 1200. Note
that D23 is very small compared to D1 in a wide range 0 ≤
J ′/J .0.25.
tence of a finite energy gap.
Figure 4 shows the minimum gap energy in the 1BZ
as a function of anisotropy J ′/J , changing the system
size L. We can see the almost gapless excitations in the
wide parameter range 0 ≤ J ′/J . 0.25, as is already
expected. It is quite natural that this behavior is simi-
Fig. 3. Anisotropy dependence of the one-particle excitation
spectra. Contour plots of the spectra are on the left, and sections
along ky = 0 line are on the right. The hexagons with broken
lines represent 1BZ of the triangular lattice. Up to J ′/J ∼ 0.25,
the spectra for each anisotropy are hardly distinguishable, and
the one-dimensionality strongly remains for large J ′/J .
4 J. Phys. Soc. Jpn. Letter Author Name
lar to that of D23, considering that the minimum energy
excitations are located along kx = ±π/2 for J ′/J . 0.6.
By plotting the same data for various system size, L,
in a semi-log scale (Fig. 4), we can see a discontinuous
jump for every size. We find that this critical value J ′c/J
vanishes very slowly as (lnL)−1. Thus, the discontinu-
ity is an artifact of finite-size calculation. We also find
that the minimum gap energy is finite when infinitesimal
J ′ is introduced. Actually, we can fit the J ′ dependence
as aJ ′ exp(−bJ/J ′)17 for J ′/J . 0.6 as shown in Fig.
4. Considering that the minimum gap energy is already
about 3 orders of magnitude below J at J ′/J ∼ 0.25, it
can be said that almost gapless excitation is realized in
0 ≤ J ′/J . 0.25. This result is fairly suggestive com-
pared with the previous series expansion18 and linear
spin wave19, 20 studies, all of which suggest a spin dis-
ordered state in the parameter range J ′/J . 0.25.
From the above results, we conclude that there is a
strong tendency to form a 1D-like excitation spectrum
for the triangular lattice spin system with anisotropy
0 ≤ J ′/J . 0.6. Furthermore, even if the anisotropy
is as large as 0.6 . J ′/J . 0.8, we can still expect 1D-
like behavior for quasiparticles except for the existence of
the excitation gap. Let us here discuss the relation to κ-
(ET)2Cu2(CN)3. The anisotropy of spin exchange inter-
actions in this material can be estimated from J = 4t2/U
(U being the onsite Coulomb repulsion) as J ′/J ∼ 0.89.
At this anisotropy, a rather large excitation gap exists as
shown in Fig. 4. We consider two possibilities to under-
stand the gaplessness. One is that the small gap region in
Fig. 4 expands to large values of J ′/J by some factors not
considered in the present model. For example, If long-
distance exchange interactions, quantum fluctuation or
multiple spin exchange effect14 (higher order terms of
the Heisenberg model) suppress not only LRMO but also
the spin gap, we can reproduce the gapless spin liquid
state at large J ′/J . These possibilities remain as future
problems. Another possibility is that the anisotropy J ′/J
of κ-(ET)2Cu2(CN)3 deviates from the above estimation
Fig. 4. (Color Online) Anisotropy dependence of the minimum
gap energy in the 1BZ (right axis) for L=60(diamond), 120(plus),
300(square), 600(cross) and 1200(triangle). The semi-log plots of
the same quantity are also shown (left axis). The solid line is a
fitted exponential function aJ ′ exp(−bJ/J ′), where a = 3.50 and
b = 1.61. We find that the observed critical behavior is an artifact
of finite size calculation (see the text).
due to, for example, a finite U effect.21 If it is in the
range J ′/J < 0.25, the excitation gap is sufficiently small
and the susceptibility behavior (finite at 1.9K whereas
J ∼ 250K) can be explained.
In summary, we analyzed an anisotropic triangular lat-
tice Heisenberg model using RVB mean-field approxima-
tion in order to investigate the physical origin of the gap-
less spin liquid state observed in κ-(ET)2Cu2(CN)3. We
payed attention to the Q1D anisotropy of this material,
and took an approach from the 1D limit. As a result of
calculations, we found that a practically 1D state with
almost gapless excitations is realized in the wide range of
the anisotropy parameter 0 ≤ J ′/J . 0.25. Furthermore,
one-dimensionality remained strongly even in J ′/J >
0.25 due to the geometrical frustration of interchain cou-
plings. We consider this “one-dimensionalization by frus-
tration” as a candidate for the mechanism of the gapless
spin liquid state, although the full understanding has not
yet been achieved.
This work was partly supported by a Grant-in-Aid
for Scientific Research on Priority Areas of Molecular
Conductors (No. 15073210) from the Ministry of Edu-
cation, Culture, Sports, Science and Technology, Japan,
and also by a Next Generation Supercomputing Project,
Nanoscience Program, MEXT, Japan.
1) For a review, see T.Ishiguro, K.Yamaji and G.Saito: Organic
Superconductors (Springer-Verlag, Berlin, 1998), 2nd ed.
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B 50 (1994) 10048.
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(1993) 1629.
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(1999) 3899.
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11) G.Baskaran: Phys. Rev. Lett. 91 (2003) 097003.
12) T.Watanabe, H.Yokoyama, Y.Tanaka, J.Inoue and M.Ogata:
J. Phys. Soc. Jpn 73 (2004) 3404.
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Technol. 1 (1988) 36.
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17) We would like to thank T.Misawa for pointing out this possi-
bility.
18) W.Zheng, R.H.McKenzie and R.R.P.Singh: Phys. Rev. B 59
(1999) 14367.
19) J.Merino, R.H.McKenzie, J.B.Marston and C.H.Chung: J.
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|
0704.0314 | Extra dimensions and Lorentz invariance violation | Extra dimensions and Lorentz invariance violation
Viktor Baukh∗ and Alexander Zhuk†
Department of Theoretical Physics and Astronomical Observatory,
Odessa National University, 2 Dvoryanskaya St., Odessa 65026, Ukraine
Tina Kahniashvili‡
CCPP, New York University, 4 Washington Place, New York, NY 10003, USA
National Abastumani Astrophysical Observatory, 2A Kazbegi Ave, Tbilisi, GE-0160 Georgia
We consider effective model where photons interact with scalar field corresponding to conformal
excitations of the internal space (geometrical moduli/gravexcitons). We demonstrate that this
interaction results in a modified dispersion relation for photons, and consequently, the photon group
velocity depends on the energy implying the propagation time delay effect. We suggest to use the
experimental bounds of the time delay of gamma ray bursts (GRBs) photons propagation as an
additional constrain for the gravexciton parameters.
PACS numbers: 04.50.+h, 11.25.Mj, 98.80.-k
Lorentz invariance (LI) of physical laws is one of the
corner stone of modern physics. There is a number of ex-
periments confirming this symmetry at energies we can
approach now. For example, on a classical level, the ro-
tation invariance has been tested in Michelson-Morley
experiments, and the boost invariance has been tested
in Kennedy-Torhndike experiments [1]. Although, up
to now, LI is well established experimentally, we can-
not say surely that at higher energies it is still valid.
Moreover, modern astrophysical and cosmological data
(e.g. UHECR, dark matter, dark energy, etc) indicate
for a possible LI violation (LV). To resolve these chal-
lenges, there are number of attempts to create new phys-
ical models, such as M/string theory, Kaluza-Klein mod-
els, brane-world models, etc. [1].
In this paper we investigate LV test related to photon
dispersion measure (PhDM). This test is based on the
LV effect of a phenomenological energy-dependent speed
of photon [2, 3, 4, 5, 6, 7, 8], for recent studies see Ref.
[9] and references therein.
The formalism that we use is based on the analogy
with electromagnetic waves propagation in a magnetized
medium, and extends previous works [8, 10, 11]. In our
model, instead of propagation in a magnetized medium,
the electromagnetic waves are propagating in vacuum
filled with a scalar field ψ. LV occurs because of an in-
teraction term f(ψ)F 2 where F is an amplitude of the
electromagnetic field. Such an interaction might have
different origins. In the string theory ψ could be a dila-
ton field [12, 13]. The field ψ could be associated with
geometrical moduli. In brane-world models the similar
term describes an interaction between the bulk dilaton
and the Standard Model fields on the brane [14]. In
Ref. [15], such an interaction was obtained in N = 4
∗Electronic address: bauch˙[email protected]
†Electronic address: [email protected]
‡Electronic address: [email protected]
super-gravity in four dimensions. In Kaluza-Klein mod-
els the term f(ψ)F 2 has the pure geometrical origin, and
it appears in the effective, dimensionally reduced, four
dimensional action (see e.g. [16, 17]). In particular, in
reduced Einstein-Yang-Mills theories, the function f(ψ)
coincides (up to a numerical prefactor) with the volume
of the internal space. Phenomenological (exactly solv-
able) models with spherical symmetries were considered
in Refs. [18]. To be more specific, we consider the model
which is based on the reduced Einstein-Yang-Mills the-
ory [17], where the term ∝ ψF 2 describes the interaction
between the conformal excitations of the internal space
(gravexcitons) and photons. It is clear that the similar
LV effect exists for all types of interactions of the form
f(ψ)F 2 mentioned above.
Obviously, the interaction term f(ψ)F 2 modifies the
Maxwell equations, and, consequently, results in a mod-
ified dispersion relation for photons. We show that
this modification has rather specific form. For example,
we demonstrate that refractive indices for the left and
right circularly polarized waves coincide with each other.
Thus, rotational invariance is preserved. However, the
speed of the electromagnetic wave’s propagation in vac-
uum differs from the speed of light c. This difference
implies the time delay effect which can be measured via
high-energy GRB photons propagation over cosmological
distances (see e.g. Ref. [9]). It is clear that gravexcitons
should not overclose the Universe and should not result
in variations of the fine structure constant. These de-
mands lead to a certain constrains for gravexcitons (see
Refs. [17, 19]). We use the time delay effect, caused by
the interaction between photons and gravexcitons, to get
additional bounds on the parameters of gravexcitons.
The starting point of our investigation is the Abelian
part of D-dimensional action of the Einstein-Yang-Mills
theory:
SEM = −
|g|FMNFMN , (1)
http://arxiv.org/abs/0704.0314v4
mailto:[email protected]
mailto:[email protected]
mailto:[email protected]
where the D-dimensional metric, g = gMN (X)dX
dXN = g(0)(x)µνdx
µ ⊗ dxν + a21(x)g(1), is defined on
the product manifold M = M0 × M1. Here, M0 is
the (D0 = d0 + 1)-dimensional external space. The d1-
dimensional internal space M1 has a constant curvature
with the scale factor a1(x) ≡ LPl expβ1(x). Dimensional
reduction of the action (1) results in the following effec-
tive D0-dimensional action [17]
S̄EM = −
|g̃(0)| [(1−Dκ0ψ)FµνFµν ] , (2)
which is written in the Einstein frame with the D0-
dimensional metric, g̃
µν = (exp d1β̄
1)−2/(D0−2)g
Here, κ0ψ ≡ −β̄1
(D0 − 2)/d1(D − 2) ≪ 1 and β̄1 ≡
β1 − β10 are small fluctuations of the internal space scale
factor over the stable background β10 (0 subscript de-
notes the present day value). These internal space scale-
factor small fluctuations/oscillations have the form of
a scalar field (so called gravexciton [20]) with a mass
mψ defined by the curvature of the effective potential
(see for detail [20]). Action (2) is defined under the
approximation κ0ψ < 1 that obviously holds for the
condition1 ψ < MPl. κ
0 = 8π/M
Pl is four dimen-
sional gravitational constant, MPl is the Plank mass,
D = 2
d1/[(D0 − 1)(D − 1)] is a model dependent con-
stant. The Lagrangian density for the scalar field ψ reads:
|g̃(0)|(−g̃µνψ,µψ,ν−m2ψψψ)/2. For simplicity we
assume that g̃0 is the flat Friedman-Lemaitre-Robertson-
Walker (FLRW) metric with the scale factor a(t).
Let’s consider Eq. (2). It is worth of noting that the
D0-dimensional field strength tensor, Fµν , is gauge in-
variant.2 Secondly, action (2) is conformally invariant in
the case when D0 = 4. The transform to the Einstein
frame does not break gauge invariance of the action (2),
and the electromagnetic field is antisymmetric as usual,
Fµν = ∂µAν − ∂νAµ. Varying (2) with respect to the
electromagnetic vector potential,
−g (1−Dκ0ψ)Fµν
= 0. (3)
The second term in the round brackets Dκ0ψFµν reflects
the interaction between photons and the scalar field ψ,
and as we show below, it is responsible for LV. In par-
ticular, coupling between photons and the scalar field ψ
makes the speed of photons different from the standard
speed of light. Eq. (3) together with Bianchi identity
(which is preserved in the considered model due to gauge-
invariance of the tensor, Fµν [17]) defines a complete set
1 In the brane-world model the prefactor κ0 in the expression for
κ0ψ is replaced by the parameter proportional to M
[14].
Thus, the smallness condition holds for ψ < MEW .
2 Eq. (2) can be rewritten in the more familiar form S̄EM =
−(1/2)
|g̃(0)|F̄µν F̄
µν [17]. The field strength tensor
F̄µν is not gauge invariant here.
of the generalized Maxwell equations. As we noted, ac-
tion (2) is conformally invariant in the 4D dimensional
space-time. So, it is convenient to present the flat FLRW
metric g̃0 in the conformally flat form: g̃0µν = a
2ηµν ,
where ηµν is the Minkowski metric.
Using the standard definition of the electromagnetic
field tensor, Fµν , we obtain the complete set of the
Maxwell equations in vacuum,
∇ ·B = 0 , (4)
∇ ·E = Dκ0
1−Dκ0ψ
(∇ψ ·E) , (5)
∇×B = ∂E
− Dκ0ψ̇
1−Dκ0ψ
1−Dκ0ψ
[∇ψ ×B] , (6)
∇×E = −∂B
, (7)
where all operations are performed in the Minkowski
space-time, η denotes conformal time related to physi-
cal time t as dt = a(η)dη, and an overdot represents a
derivative with respect to conformal time η.
Eqs. (4) and (7) correspond to Bianchi identity, and
since it is preserved, Eqs. (4) and (7) keep their usual
forms. Eqs. (5) and (6) are modified due to interactions
between photons and gravexcitons (∝ κ0ψ). These mod-
ifications have simple physical meaning: the interaction
between photons and the scalar field ψ acts as an effective
electric charge eeff . This effective charge is proportional
to the scalar product of the ψ field gradient and the E
field, and it vanishes for an homogeneous ψ field. The
modification of Eq. (6) corresponds to an effective cur-
rent Jeff , which depends on both electric and magnetic
fields. This effective current is determined by variations
of the ψ field over the time (ψ̇) and space (∇ψ). For
the case of a homogeneous ψ field the effective current is
still present and LV takes place. The modified Maxwell
equations are conformally invariant. To account for the
expansion of the Universe we rescale the field components
asB,E → B,E a2 [21].
To obtain a dispersion relation for photons, we use
the Fourier transform between position and wavenumber
spaces as,
F(k, ω) =
dη d3x e−i(ωη−k·x)F(x, η) ,
F(x, η) =
(2π)4
dω d3kei(ωη−k·x)F(k, ω) . (8)
Here, F is a vector function describing either the elec-
tric or the magnetic field, ω is the angular frequency of
the electro-magnetic wave measured today, and k is the
wave-vector. We assume that the field ψ is an oscilla-
tory field with the frequency ωψ and the momentum q,
so ψ(x, η) = Cei(ωψη−q·x) , C = const . Eq. (4) implies
B ⊥ k. Without loosing of generality, and for simplic-
ity of description we assume that the wave-vector k is
oriented along the z axis. Using Eq. (7) we get E ⊥ B.
A linearly polarized wave can be expressed as a super-
position of left (L, −) and right (R, +) circularly polar-
ized (LCP and RCP) waves. Using the polarization basis
of Sec. 1.1.3 of Ref. [22], we derive E± = (Ex± iEy)/
Rewriting Eqs. (4) - (7) in the components,3 for LCP
and RCP waves we get,
(1 − n2+)E+ = 0, (1− n2(−))E
− = 0 , (9)
where n+ and n− are refractive indices for RCP and LCP
electromagnetic waves
n2+ =
k2 [1−Dκ0ψ(1 + qz/k)]
ω2 [1−Dκ0ψ(1 + ωψ/ω)]
= n2− . (10)
In the case when LI is preserved the electromagnetic
waves propagating in vacuum have n+ = n− = n =
k/ω ≡ 1. For the electromagnetic waves propagating in
the magnetized plasma, k/ω 6= 1, and the difference be-
tween the LCP and RCP refractive indices describes the
Faraday rotation effect, α ∝ ω(n+ − n−) [23]. In the
considered model, since n+ = n− the rotation effect is
absent, but the speed of electromagnetic waves propaga-
tion in vacuum differs from the speed of light c (see also
Ref. [24] for LV induced by electromagnetic field cou-
pling to other generic field). This difference implies the
propagation time delay effect, ∆t = ∆l(1−∂k/∂ω) (∆l is
a propagation distance), ∆t is the difference between the
photon travel time and that for a ”photon” which travels
at the speed of light c. Here, t is physical synchronous
time. This formula does not take into account the evo-
lution of the Universe. However, it is easy to show that
the effect of the Universe expansion is negligibly small.
Solving the dispersion relation as a square equation,
we obtain
ω2ψ − q2z
(Dκ0ψ)2
, (11)
where ± signs correspond to photons forward and back-
ward directions respectively.
The modified inverse group velocity (11) shows that
the LV effect can be measured if we know the gravexciton
frequency ωψ, z-component of the momentum qz and its
amplitude ψ. For our estimates, we assume that ψ is
the oscillatory field, satisfying (in local Lorentz frame)
the dispersion relation, ω2ψ = m
ψ + q
2, where mψ is the
mass of gravexcitons4. Unfortunately, we do not have
3 We have defined the system of 6 equations with respect to 6
components of the vectors E and B. This system has non-trivial
solutions only if its determinant is nonzero. From this condition
we get the dispersion relation. The Faraday rotation effect is
absent if the matrix has a diagonal form.
4 To get physical values of the corresponding parameters we should
rescale them by the scale factor a.
any information concerning parameters of gravexcitons
(some estimates can be found in [17, 19]). Thus, we
intend to use possible LV effects (supposing it is caused
by interaction between photons and gravexcitons) to set
limits on gravexciton parameters. For example, we can
easily get the following estimate for the upper limit of
the amplitude of gravexciton oscillations:
|ψ| ≈ 1√
MPl , (12)
where for ω and mψ we can use their physical values.
In the case of GRB with ω ∼ 1021 ÷ 1022Hz ∼ 10−4 ÷
10−3GeV and ∆l ∼ 3 ÷ 5 × 109y ∼ 1017sec the typical
upper limit for the time delay is ∆t ∼ 10−4sec [9]. For
these values the upper limit on gravexciton amplitude of
oscillations is5
|κ0ψ| ≈
10−13GeV
. (13)
This estimate shows that our approximation κ0ψ < 1
works for gravexciton masses mψ > 10
−13GeV. Future
measurements of the time-delay effect for GRBs at fre-
quencies ω ∼ 1 − 10GeV would increase significantly
the limit up to mψ > 10
−9GeV. On the other hand,
Cavendish-type experiments [26, 27]) exclude fifth force
particles with masses mψ . 1/(10
−2cm) ∼ 10−12GeV
which is rather close to our lower bound for ψ field
masses. Respectively we slightly shift the considered
mass lower limit to be mψ ≥ 10−12GeV. These masses
considerably higher than the mass corresponding to the
equality between the energy densities of the matter and
radiation (matter/radiation equality), meq ∼ Heq ∼
10−37GeV, where Heq is the Hubble ”constant” at mat-
ter/radiation equality. It means that such ψ-particles
start to oscillate during the radiation dominated epoch
(see appendix). Another bound on the ψ-particles masses
comes from the condition of their stability. With re-
spect to decay ψ → γγ the life-time of ψ-particles is
τ ∼ (MPl/mψ)3tPl [17], and the stability conditions re-
quires that the decay time should be greater than the age
of the Universe. According this we consider light gravex-
citons with masses mψ ≤ 10−21MPl ∼ 10−2GeV ∼ 20me
(where me is the electron mass).
As an additional restriction arises from the condi-
tion that such cosmological gravexcitons should not
overclose the observable Universe. This reads mψ .
meq(MPl/ψin)
4 which implies the following restriction
for the amplitude of the initial oscillations: ψin .
(meq/mψ)
MPL << MPl [19]. Thus, for the range of
masses 10−12GeV ≤ mψ ≤ 10−2GeV, we obtain respec-
tively ψin . 10
−6MPl and ψin . 10
−9MPl. According to
5 We thank R. Lehnert to point that in addition of the time de-
lay effect the Cherenkov effect could be used to constrain the
electromagnetic field and ψ field coupling strength [25].
Eq. (A.3), we can also get the estimate for the amplitude
of oscillations of the considered gravexciton at the present
time. Together with the non-overcloseness condition,
we obtain from this expression that |κ0ψ| ∼ 10−43 for
mψ ∼ 10−12GeV and ψin ∼ 10−6MPl and |κ0ψ| ∼ 10−53
for mψ ∼ 10−2GeV and ψin ∼ 10−9MPl. Obviously, it is
much less than the upper limit (13). Note, as we men-
tioned above, gravexcitons with masses mψ & 10
−2GeV
can start to decay at the present epoch. However, taking
into account the estimate |κ0ψ| ∼ 10−53, we can easily
get that their energy density ρψ ∼ (|κ0ψ|2/8π)M2Plm2ψ ∼
10−55g/cm3 is much less than the present energy density
of the radiation ργ ∼ 10−34g/cm3. Thus, ρψ contributes
negligibly in ργ . Otherwise, the gravexcitons with masses
mψ & 10
−2GeV should be observed at the present time,
which, obviously, is not the case.
Additionally, it follows from Eq. (42) in Ref. [17]
that to avoid the problem of the fine structure constant
variation, the amplitude of the initial oscillations should
satisfy the condition: ψin . 10
−5MPl which, obviously,
completely agrees with our upper bound ψin . 10
−6GeV.
Summarizing we shown that LV effects can give addi-
tional restrictions on parameters of gravexcitons. First,
we found that gravexcitons should not be lighter than
10−13GeV. It is very close to the limit following from the
fifth-force experiment. Moreover, experiments for GRB
at frequencies ω > 1GeV can result in significant shift of
this lower limit making it much stronger than the fifth-
force estimates. Together with the non-overcloseness con-
dition, this estimate leads to the upper limit on the am-
plitude of the gravexciton initial oscillations. It should
not exceed ψin . 10
−6GeV. Thus, the bound on the ini-
tial amplitude obtained from the fine structure constant
variation is one magnitude weaker than our one even for
the limiting case of the gravexciton masses. Increasing
the mass of gravexcitons makes our limit stronger. Our
estimates for the present day amplitude of the gravexci-
ton oscillations, following from the obtained above lim-
itations, show that we cannot use the LV effect for the
direct detections of the gravexcitons. Nevertheless, the
obtained bounds can be useful for astrophysical and cos-
mological applications. For example, let us suppose that
gravexcitons with masses mψ > 10
−2GeV are produced
during late stages of the Universe expansion in some re-
gions and GRB photons travel to us through these re-
gions. Then, Eq. (A.3) is not valid for such gravexcitons
having astrophysical origin and the only upper limit on
the amplitude of their oscillations (in these regions) fol-
lows from Eq. (13). In the case of TeV masses we get
|κ0ψ| ∼ 10−16. If GRB photons have frequencies up to
1 TeV, ω ∼ 1TeV, then this estimate is increased by 6
orders of magnitude.
Acknowledgments
We thank G. Dvali, G. Gabadadze, A. Gruzinov, G.
Melikidze, B. Ratra, and A. Starobinsky for stimulating
discussions. T. K. and A. Zh. acknowledge hospital-
ity of Abdus Salam International Center for Theoreti-
cal Physics (ICTP) where this work has been started.
A.Zh. would like to thank the Theory Division of CERN
for their kind hospitality during the final stage of this
work. T.K. acknowledges partial support from INTAS
061000017-9258 and Georgian NSF ST06/4-096 grants.
A. Appendix: Dynamics of Light Gravexcitons
In this appendix we briefly summarize the main prop-
erties of the light gravexcitons necessary for our inves-
tigations. The more detail description can be found in
Refs. [17, 19].
The effective equation of motion for massive cosmolog-
ical gravexciton6 is
ψ + (3H + Γ)
ψ +m2ψψ = 0 , (A.1)
where H ∼ 1/t and Γ ∼ m3ψ/M2Pl are the Hubble pa-
rameter and decay rate (ψ → γγ) correspondingly. This
equation shows that at times when the Hubble parame-
ter is less than the gravexciton mass: H . mψ the scalar
field begins to oscillate (i.e. time tin ∼ H−1in ∼ 1/mψ
roughly indicates the beginning of the oscillations):
ψ ≈ CB(t) cos(mψt+ δ) . (A.2)
We consider cosmological gravexcitons with masses
10−12GeV ≤ mψ ≤ 10−2GeV. The lower bound fol-
lows both from the fifth-force experiments and Eq. (13).
The upper bound follows from the demand that the life-
time of these particles (with respect to decay ψ → γγ)
is larger than the age of the Universe: τ = 1/Γ ∼
(MPl/mψ)
tPl ≥ 1019sec > tuniv ∼ 4 × 1017 sec. Thus,
we can neglect the decay processes for these gravexci-
tons. Additionally, it can be easily seen that these par-
ticles start to oscillate before teq ∼ H−1eq when the en-
ergy densities of the matter and radiation become equal
to each other (matter/radiation equality). According to
the present WMAP data for the ΛCDM model it holds
Heq ≡ meq ∼ 10−56MPl ∼ 10−28eV. Thus, considered
particles have masses mψ >> meq and start to oscil-
late during the radiation dominated stage. They will not
overclose the observable Universe if the following condi-
tion is satisfied: mψ . meq(MPl/ψin)
4, where ψin is the
amplitude of the initial oscillations at the moment tin
(see Eq. (18) in Ref. [19]).
Prefactors C and B(t) in Eq. (A.2) for con-
sidered light gravexcitons respectively read: C ∼
(ψin/MPl) (MPl/mψ)
and B(t) ∼ MPl (MPlt)−3s/2.
Here, s = 1/2, 2/3 for oscillations during the radiation
6 We have seen that the interaction between gravexcitons and or-
dinary matter (in our case it is 4D-photons) is suppressed by the
Planck scale. Thus, gravexcitons are weakly interacting massive
particles (WIMPs).
dominated and matter dominated stages, correspond-
ingly. We are interested in the gravexciton oscillations
at the present time t = tuniv. In this case s = 2/3 and
for B(tuniv) we obtain: B(tuniv) ∼ t−1univ ≈ 10−61MPl.
Thus, the amplitude of the light gravexciton oscillations
at the present time reads:
|κ0ψ| ∼ 10−60
. (A.3)
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|
0704.0315 | The small deviations of many-dimensional diffusion processes and
rarefaction by boundaries | arXiv:0704.0315v1 [math.PR] 3 Apr 2007
THE SMALL DEVIATIONS OF MANY-DIMENSIONAL DIFFUSION
PROCESSES AND RAREFACTION BY BOUNDARIES
Vitalii A. Gasanenko
Abstract. We lead the algorithm of expansion of sojourn probability of many-dimensional
diffusion processes in small domain. The principal member of this expansion defines nor-
malizing coefficient for special limit theorems.
Introduction.
Let ξ(t) be a random process with measurable phase space (X,Σ(X)). Consider the
measurable connected domain D ∈ Σ(X) and small parameter ǫ. The investigations of
asymptotics of sojourn probability (small deviations)
P (ξ(t) ∈ ǫD, t ∈ [0, T ]) (1)
is jointed with many practice and theoretical problems [1-4]. In the literature, it was
researched both rough asymptotics of principal member of (1)(log from it)[5] and exact
asymptotics of diffusion processes of (1)[6-8]. In the works [9,10] was proved of algo-
rithms of expansions of exact asymptotics of small deviation for diffusion and piecewise
deterministic random processes for one-dimensional case.
The purpose this article is to present the algorithm of expansion of small deviation for
many-dimensional diffusion processes and to define all constants of principal member.
In Section 1 our main result is stated and proved. In section 2 we consider the limits
theorems about numbers of unabsorbed diffusion particles by boundaries of small domain.
I. The expansion.
We shall investigate of asymptote of following probability
P (ǫ, x) = P (ξ(t) ∈ ǫD, 0 ≤ t ≤ T ) , ǫ → 0,
where ξ(t) ∈ Rd is solution of the following stochastic differential equation
dξ(t) = a(t, ξ(t))dt+
bi(ξ(t))dwi(t), ξ(0) = x ∈ ǫD. (2)
where functions
1991 Mathematics Subject Classification. 60 J 65.
Key words and phrases. parabolic problem,small domain, algorithm of expansion, number of unab-
sorbed processes.
Typeset by AMS-TEX
http://arxiv.org/abs/0704.0315v1
2 VITALII A. GASANENKO
bi(x), a(t, x) : R
d → Rd and R+ ×R
d → Rd.
are differentiable.
Set σij(x) =
bik(x)b
k(x).
It is known that P (ǫ, x) = uǫ0(T, x). Here u
0(t, x) is solution of the following parabolic
boundary problem at 0 ≤ t ≤ T
∂uǫ0(t, x)
i,j=1
σij(x)
∂2uǫ0(t, x)
∂xi∂xj
ai(T − t, x)
∂uǫ0(t, x)
, x ∈ Dǫ;
u(t, x)|t=0 = 1; x ∈ Dǫ; u(t, x) = 0 x ∈ ∂Dǫ, 0 ≤ t ≤ T. (3)
where Dǫ = ǫD. It is assumed that D is a connected bounded domain from R
m; the
boundary ∂Q is the Lyapunov surface C(1,λ) and 0 ∈ D. We interest of the asymptotic
expansion ǫ → 0 of solution this problem uǫ0(t, x) at ǫ → 0.
We define the differential operator A : 1
1≤i,j≤d
σij(0)
∂xi∂xj
. Let σ be a matrix with
the following property
1≤i,j≤d
σij(0)zizj ≥ µ|~z|
Here µ, there is a fixed positive number, and ~z = (z1, · · · , zd) is an arbitrary real
vector.
This operator acts in the following space
HA = {u : u ∈ L2(D) ∩ Au ∈ L2(D) ∩ u(∂D) = 0}
with inner product (u, v)A = (Au, v). Here (, ) is inner product in L2(Q). The opera-
tor A is a positive operator[11]. It is known that the following eigenvalue problem
Au = −λu, u(∂D) = 0
has infinite set of real eigenvalues λi → ∞ and
0 < λ1 < λ2 < · · · < λs < · · · .
The corresponding eigenfunctions
f11, . . . , f1n1 , · · · , fs1, . . . , fsns , · · ·
form the complete system of functions both in HA and L
2(Q) := {u : u ∈ L2(Q) ∩
u(∂Q) = 0}. Here the number nk is equal to multiplicity of eigenvalue λk.
It is often convenient to present the system of eigenfunctions by one index: {fn(z)}.
The corresponding system of eigenvalues {λn} will be with recurrences. We shall use it
We introduce the spectral function
e(x, y, λ) =
fj(x)fj(y).
We shall need in the following theorem from the monograph [12].
THE SMALL DEVIATION OF MANY-DIMENSIONAL DIFFUSION PROCESSES 3
Theorem 1 ([12].Th.17.5.3). . There exists such constant Cα that
x,y∈D
|Dαx,ye(x, y, λ)| ≤ Cαλ
(n+|α|)/2
Here α is multi-index.
Theorem 2. . If the surface ∂D is Lyapunov surface and
(t,z)∈[0,T ]×D,1≤i,j≤d
∂ai(t, z)
∂bi(z)
∂ai(T − t, z)
then the following relation takes place at ǫ → 0
P (ǫ, zǫ) = exp
µ(t)dt
c1mf1m(z) (1 +O(ǫ)) , at z ∈ D,
where
µ(t) =
σij(0)ai(t, 0)aj(t, 0)− δijai(t, 0)aj(t, 0)
and c1m =
f1m(z)dz.
Proof. Make the change of variables and function
xi = ziǫ, u
1 = u
0 exp
ak(T − t, 0)zk
Now we obtain the following parabolic problem for function uǫ1
∂uǫ1(t, z)
i,j=1
σij(ǫz)
∂2uǫ1(t, z)
∂zi∂zj
ai(T − t, ǫz)−
σij(ǫz)aj(T − t, 0)
∂uǫ1(t, z)
σij(ǫz)ai(T − t, 0)aj(T − t, 0)− δijai(T − t, 0)aj(T − t, ǫz)− ǫ
∂ai(T − t, 0)
uǫ1, z ∈ D;
1(t, z)|t=0 = exp
ak(T, 0)zk
; z ∈ D; uǫ1(t, z) = 0 z ∈ ∂D, 0 ≤ t ≤ T.
We will construct the asymptotic expansion of solution for this initial - boundary
problem in the following form
uǫ1(t, z) =
vk(t, z)ǫ
k. (5)
Note that the famous expansion
4 VITALII A. GASANENKO
ak(T, 0)zk
= 1 + ǫ
ak(T, 0)zk +
ak(T, 0)zk
+ · · · ,
defines the initial conditions for vk, k ≥ 0:
v0(0, z) = 1, v1(0, z) =
ak(T, 0)zk, v2(0, z) =
ak(T, 0)zk
· · · .
Using the first fragment of Taylor series in zero point under conditions of theorem we
can obtain the following representations
σij(ǫz) = σij(0) + ǫσ
ij(z), ai(T − t, ǫz) = ai(T − t, 0) + ǫa
i(T − t, z), 1 ≤ i, j ≤ d (6)
where
z∈D,ǫ∈[0,1],1≤i,j≤d
|σǫij(z)| < ∞, sup
z∈D,t∈[0,T ],ǫ∈[0,1],1≤i≤d
|aǫi(T − t, z)| < ∞
Now, after substitution of (5),(6) to (4) we conclude that the v0 satisfies the problem
i,j=1
σij(0)
∂zi∂zj
v0 + µ(t)v0 (7)
v0|∂D = 0; v0(0, z) = 1, z ∈ D.
µ(t) =
σij(0)ai(T − t, 0)aj(T − t, 0)− δijai(T − t, 0)aj(T − t, 0)
Further, let us denote by Bǫ(t, z) the operator C
2(D) → C(D), for f ∈ C2(D) it’s
defined as follows:
ǫ(t, z)f =
i,j=1
σǫij(z)
∂zi∂zj
ai(T − t, ǫz)−
σij(ǫz)aj(T − t, 0)
i,j=1
σǫij(z)ai(T − t, 0)aj(T − t, 0)− δijai(T − t, 0)a
j(T − t, z)−
∂ai(T − t, 0)
i,j=1
σǫij(z)
∂zi∂zj
Aǫ1(t, z)f + ǫA
2(t, z).
THE SMALL DEVIATION OF MANY-DIMENSIONAL DIFFUSION PROCESSES 5
Now, formally the functions vk, k ≥ 1 are defined by the following recurrence system
problems
i,j=1
σij(0)
∂zi∂zj
vk +Bǫ(t, z)vk−1 (8)
v0|∂D = 0; vk(0, z) =
ak(T − t, 0)zk
, z ∈ D.
We shall solve the problems of (7),(8) by method of separation of variables. According
to this method the solutions are defined in the form
vk(t, z) =
qk,n(t)fn(z). (9)
For definition of principal number it suffices to construct of the v0. If we substitute
(9) at k = 0 to (7) then we obtain
−q̇0,n(t)−
q0,n(t) + µ(t)q0,n(t)
fn(z) = 0.
Set c0,n =
fn(z)dz (coefficients of expansion of indicator of set D). The initial
condition of v0 has the following stating
v0(0, z) =
q0,n(0)fn(z) =
c0,nfn(z) =
c0,lmflm(z), z ∈ D.
By definition of system of functions {fn(z)}, now we have the system of ordinary
differential equations
q̇0,n(t) +
− µ(t)
q0,n(t) = 0, q0,n(0) = c0,n.
From the latter one we have
q0,n(t) = c0,n exp
µ(s)ds
A0 = sup
ǫ≤1,z∈D;i,j
|σǫij(z)|, L0 =
l≥1,1≤m≤nl
(c0,ml)
A1 = sup
0≤ǫ≤1,z∈D,t∈[0,T ];i,j
ai(T − t, ǫz)−
σij(ǫz)aj(T − t, 0)
= sup
0≤ǫ≤1,z∈D,t∈[0,T ];i,j
ij(z)ai(T − t, 0)aj(T − t, 0)− δijai(T − t, 0)a
j(T − t, z)−
∂ai(T − t, 0)
We have the following relations for eigenvalues λl
6 VITALII A. GASANENKO
2/d ≤ λl ≤ k2l
2/d, max(k1, k2) < ∞
Applying Cauchy-Bunyakovskii inequality, Theorem 1 and the latter one, we get
aǫi,j(z)
∂zi∂zj
−λltǫ
µ(s)ds
c0,ml
aǫi,j(z)
∂2fml(z)
∂zi∂zj
≤ A0d
−λltǫ
µ(s)ds
(c0,ml)
∂2fml(z)
∂zi∂zj
≤ A0dC2,2L0
µ(s)ds
l ≤ exp
K0. (10)
Here K0 < ∞.
Reasoning similarly we convince ourselves that for other parts of Bǫ(t, z)v0 the fol-
lowing estimations take place
|Aǫ1(t, z)v0| ≤ A1dC1,1L0
µ(s)ds
l ≤ exp
K0,1; (11)
|Aǫ2(t, z)v0| ≤ A2dC0,0L0
µ(s)ds
l ≤ exp
K0,2, (12)
where max{K0,1,K0,2} < ∞.
Now let us estimate the coefficients βǫn(t) of expansion of B
ǫ(t, z)v0 by system {fn}n≥1.
Applying (10)-(12) and Cauchy-Bunyakovskii inequality, we get
|βǫn(t)| = |
Bǫ(t, z)v0(t, z)fn(z)dz| ≤
(Bǫ(t, z)v0)
n(z)dz
≤ exp(−λ1tǫ
K0 +K0,1
+ ǫK0,2
The latter one now gives
βǫn(s)ds| ≤ ǫγǫ(t), (13)
THE SMALL DEVIATION OF MANY-DIMENSIONAL DIFFUSION PROCESSES 7
where
0≤ǫ≤1,t∈[0,T ]
γǫ(t) < ∞.
Finally, let us estimate the difference rǫ(t, z) = uǫ1(t, z)−v0(t, z). By definition, r
ǫ(t, z)
is solution of the following problem
i,j=1
σij(0)
∂zi∂zj
rǫ +Bǫ(T − t, z)v0 z ∈ D; (14)
rǫ(t, z)|t=0 = exp
ak(T, 0)zk
− 1; z ∈ D; rǫ(t, z) = 0 z ∈ ∂D, 0 ≤ t ≤ T.
It is clear that rǫ(0, z) we can present as ǫrǫ1(0, z), where r
1(0, z) is uniform bounded
function of variables ǫ ∈ [0, 1] and z ∈ D. So, the coefficients of expansion this function
by system {fn(z)} have the following forms
rǫ(0, z)fn(z)dz = ǫµ
n, where sup
0≤ǫ≤1
(µǫn)
= M < ∞. (13)
Now we have the solution of (14) in the following form
rǫ(t, z) = ǫ
µǫn exp{−λntǫ
βǫn(s)ds}fn(z)
Applying latter one ,(13),(15), Theorem 1 and Cauchy-Bunyakovskii inequality we get
at t > 0
|rǫ(t, z)ǫ−1| ≤
(µǫn)
exp{−λntǫ
βǫn(s)ds}λ
n } ≤
≤ MC0,0 exp{−λ1tǫ
−2}K0,3, where K0,3 < ∞.
The proof of theorem is completed.
Remark 1. According to the above system of problems for definition of the functions
vk, k ≥ 1, we outline the construction of coefficients qk,n)(t) for the series (8):
q̇k,n(t) = +
+ µǫk−1,n(t)
qk,n(t),
qk,n(0) =
vk(0, z)fn(z)dz =
am(T, 0)zm
fn(z)dz
Here µǫk−1,n(t) =
fn(z)B
ǫ(t, z)vk−1(t, z)dz.
8 VITALII A. GASANENKO
Remark 2. Theorem 2 is coordinated with results of works [6-8] where the principal
member of small deviations in ball are investigated for more simple SDE.
II. The rarefaction of set of diffusion processes by boundaries of small
domains.
The following problem was investigated in works[13,14]. Let a set identical diffusion
random processes start at the initial time from the different points of domain D. These
processes are diffusion processes with absorbtion on the boundary ∂D. We are interested
in distribution of the number yet absorbed at the moment T . The initial number and
initial position of diffusion processes are defined either a random Poisson measure[14]
or deterministic measure [13]. The proved limits theorems described the situation when
T → ∞ and initial number of diffusion processes depended on T and it increased at the
rise of T . The role of normalizing function played principal member of asymptote of
solution of according parabolic problem at T → ∞.
Henceforth we shall assume that considered diffusion processes satisfy of the SDE (2)
with different initial points.
Now we consider the situation when initial number of absorbing diffusion processes
in small domain ǫD depends on ǫ → 0 and it increase under the condition of decrease
of ǫ. It is not hard to show, that now normalizing function is the principal member of
parabolic problem (3) at ǫ → 0.
The proofs of stated below theorems repeat the proofs of according theorems from
[13,14] almost word for word.
We will denote by η(ǫ, T ) the number of remaining processes in the region ǫD at the
moment T .
We will also assume that σ-additive measure ν is given on the Σν- algebra sets from
D, ν(D) < ∞. All eigenfunctions fij : D → R
1 are (Σν ,ΣY ) measurable. Here ΣY is
system of Borel sets from R1. Let ⇒ denote the weak convergence of random values or
measures.
At the beginning we assume that initial number and position of diffusion processes are
defined by deterministic measure N(ǫB, ǫ), B ∈ D. Thus, N(ǫB, ǫ) is equal to number of
starting points in the set ǫB.
Let us denote by νǫ(·) the measure
νǫ(ǫB) = exp
N(ǫB, ǫ).
where B ∈ Σν .
By definition of measure νǫ(·), we have
dνǫ(x) =
, if x = xk, k = 1, · · · , N(ǫD, ǫ)
0, otherwise.
Theorem 3. Under the assumptions of the Theorem 2 let the N(ǫ·, ǫ) satisfies the con-
dition
νǫ(ǫ ·) ⇒
ν(·).
THE SMALL DEVIATION OF MANY-DIMENSIONAL DIFFUSION PROCESSES 9
Then η(ǫ, T ) ⇒ η(T ) if ǫ → 0 where η(T ) has Poisson distribution function with
parameter
a(T ) = exp
µ(s)ds
F (z)dν(z),
where F (z) =
f1i(z)c1i, c1i =
f1i(z)dz
and µ(t) is the function from Theorem 2.
Now we consider the case when the initial number and positions of processes are
defined by the random Poisson measure µ(·, ǫ) in ǫD:
P (µ(ǫA, ǫ) = k) =
mk(ǫA, ǫ)
−m(ǫA,ǫ)
where m(ǫ ·, ǫ) is finitely additive positive measure on ǫD for fixed ǫ.
We assign
g(ǫ) = exp
Theorem 4. Under the assumptions of the Theorem 2 we suppose that m(ǫ·, ǫ) holds
the condition
m(ǫB, ǫ)g(ǫ) = ν(B), B ∈ Σν .
Then η(ǫ, T ) ⇒ η(T ) if ǫ → 0 where η(T ) has the Poisson distribution function with
the parameter a(T ) from Theorem 3.
References
1. GrahamR.,Path integral formulation of general diffusion processes, Z.Phys.(1979),B
26,pp.281-290.
2. Onsager L. andMachlup S. Fluctuation and irreversible processes, I,II, Phys.Rev.(1953)
91,pp.1505-1512,1512-1515.
3. Li W. V.,Shao Q.-M., Gaussian processes:inequalities, small ball probabilities and
applications, in : Stochastic Processes:Theory and Methods, in : Handbook of Statistics,
vol.19, 2001, pp. 533-597.
4. Lifshits M.A., Asymptotic behavior of small ball probabilities, in Probab. Theory
and Math.Statist., Proc. VII International Vilnius Conference (1998), pp. 453-468.
5. Lifshits M., Simon T., Small deviations for fractional stable processes, Ann. I. H.
Poincare - PR 41 (2005) pp. 725-752.
6. Mogulskii A.A, The method of Fourier for determination of asymptotics of small
deviations of Wiener process, Siberian Math. Journ. (1982),v.22,no.3,pp.161-174.
7. Fujita T. and Kotani S., The Onsager - Machlup Function for diffusion processes,
J.Math.Kyoto Uneversity.- 1982.-vol.22,no.22.pp.131-153.
8. Zeitoni O., On the Onsager-Machlup functional of diffusion processes around non
C2 curves, Ann. Probab.(1989),vol.17, no.3, pp.1037-1054.
10 VITALII A. GASANENKO
9. Gasanenko V.A., The total asymptotic expansion of sojourn probability of diffusion
process in thin domain with moving boundaries, Ukraine Math. Journ. (1999),v.51, no.
9, pp.1155-1164.
10. Gasanenko V.A., The jump like processes in thin domain, Analytic questions of
stochastic system, Kyiv:Institute of Mathematics (1992), pp. 4-9.
11. Mihlin S.G. Partial differential linear equations (1977), Vyshaij shkola, Moskow,
12.L.Hörmander, The analysis of Linear Partial differential Operators III (1985),
Spinger-Verlag.
13.Fedullo A., Gasanenko V.A., Limit theorems for rarefaction of set of diffusion
processes by boundaries, Theory of Stochastic Processes vol. 11(27), no.1-2,2005, pp.23-
14.Fedullo A., Gasanenko V.A.,Limit theorems for number of diffusion processes,
which did not absorb by boundaries, Central European Journal of Mathematics 4(4),
2006, pp.624-634.
Institute of Mathematics, National Academy of Science of Ukraine, Tereshchenkivska
3, 252601, Kiev, Ukraine
E-mail address: [email protected] or [email protected]
|
0704.0319 | Spin-orbit coupling effect on the persistent currents in mesoscopic ring
with an Anderson impurity | Spin-orbit coupling effect on the persistent currents in mesoscopic
ring with an Anderson impurity
Guo-Hui Ding and Bing Dong
Department of Physics, Shanghai Jiao Tong University, Shanghai, 200240, China
(Dated: November 4, 2018)
Abstract
Based on the finite U slave boson method, we have investigated the effect of Rashba spin-
orbit(SO) coupling on the persistent charge and spin currents in mesoscopic ring with an Anderson
impurity. It is shown that the Kondo effect will decrease the magnitude of the persistent charge
and spin currents in this side-coupled Anderson impurity case. In the presence of SO coupling,
the persistent currents change drastically and oscillate with the strength of SO coupling. The SO
coupling will suppress the Kondo effect and restore the abrupt jumps of the persistent currents. It
is also found that a persistent spin current circulating the ring can exist even without the charge
current in this system.
PACS numbers: 73.23.Ra, 71.70.Ej, 72.25.-b
http://arxiv.org/abs/0704.0319v1
I. INTRODUCTION
Recently the spin-orbit(SO) interaction in semiconductor mesoscopic system has attracted
a lot of interest[1]. Due to the coupling of electron orbital motion with the spin degree of
freedom, it is possible to manipulate and control the electron spin in SO coupling system by
applying an external electrical field or a gate voltage, and it is believed that the SO effect
will play an important role in the future spintronic application. Actually, various interesting
effects resulting from SO coupling have already been predicted, such as the Datta-Das spin
field-effect transistor based on Rashba SO interaction[2] and the intrinsic spin Hall effect[3].
In this paper we shall focus our attention on the persistent charge current and spin cur-
rent in mesoscopic semiconductor ring with SO interaction. The existence of a persistent
charge current in a mesoscopic ring threaded by a magnetic flux has been predicted decades
ago[4], and has been extensively studied in theory[5, 6, 7, 8, 9] and also observed in various
experiments[10, 11, 12]. The reason that a persistent charge current exists may be inter-
preted as that the magnetic flux enclosed by the ring introduces an asymmetry between
electrons with clockwise and anticlockwise momentum, thus leads to a thermodynamic state
with a charge current without dissipation. For a mesoscopic ring with a texture like inho-
mogeneous magnetic field, D. Loss et al.[13] predicted that besides the charge current there
are also a persistent spin current. The origin of the persistent spin current can be related
to the Berry phase acquired when the electron spin precesses during its orbital motion. The
persistent spin current has also been studied in semiconductor system with Rashba SO cou-
pling term[14, 15, 16]. Recently it is shown that a semiconductor ring with SO coupling can
sustain a persistent spin current even in the absence of external magnetic flux[17].
For the system of a mesoscopic ring with a magnetic impurity, the persistent charge
current has been investigated in the context of a mesoscopic ring coupled with a quantum
dot[18, 19, 20, 21, 22, 23, 24], where the quantum dot acts as an impurity level and will
introduce charge or spin fluctuations to the electrons in the ring. The Kondo effect arising
from a localized electron spin interacting with a band of electrons will be essential in the
charge transport in the ring. But to our knowledge in these systems the SO effect hasn’t
been considered. It might be expected that the interplay between the Kondo effect and
the SO coupling in the ring can give new features in the persistent currents. In this paper
we shall address this problem and investigate the SO effect on persistent charge and spin
currents in the ring system with an Anderson impurity. The Anderson impurity can act as
a magnetic impurity when the impurity level is in single electron occupied state and as well
as a barrier potential in empty occupied regime.
The outline of this paper is as follows. In section II we introduce the model Hamiltonian
of the system and also the method of calculation by finite-U slave boson approach[25, 26,
27, 28]. In section III the results of persistent charge current and spin current are presented
and discussed. In Section IV we give the summary.
II. MESOSCOPIC RING WITH AN ANDERSON IMPURITY
The electrons in a closed ring with SO coupling of Rashba term can be described by
following Hamiltonian in the polar coordinates[14, 29]
Hring = ∆(−i
[(σx cosϕ+ σy sinϕ)(−i
) + h.c.] , (1)
where ∆ = h̄2/(2mea
2), a is the radius of the ring. αR will characterize the strength of
Rashba SO interaction. Φ is the external magnetic flux enclosed by the ring, and Φ0 =
2πh̄c/e is the flux quantum.
We can write the above Hamiltonian in terms of creation and annihilation operators of
electrons in the momentum space,
Hring =
mσcmσ + 1/2
[tm(c
m+1↓cm↑ + c
m−1↑cm↓) + h.c.] , (2)
where ǫm = ∆(m+φ)
2, tm = αR(m+φ),(m = 0,±1, · · · ,±M) with φ = Φ/Φ0. One can see
that the SO interaction causes the m mode electrons coupled with m + 1 and m − 1 mode
electrons and spin-flip process. We consider the system with a side-coupled impurity which
can be described by the Anderson impurity model,
σdσ + Und↑nd↓ . (3)
The tunneling between the impurity level and the ring are given by
Hd−ring = tD
(d†σcmσ + h.c) . (4)
Then the total Hamiltonian for the system should be
H = Hring +Hd +Hd−ring . (5)
In order to treat the strong on-site Coulomb interaction in the impurity level. we adopt
the finite-U slave boson approach[25, 26]. A set of auxiliary bosons e, pσ, d is introduced for
the impurity level, which act as projection operators onto the empty, singly occupied(with
spin up and spin down), and doubly occupied electron states on the impurity, respectively.
Then the fermion operators dσ are replaced by dσ → fσzσ, with zσ = e
†pσ + p
σ̄d. In order
to eliminate un-physical states, the following constraint conditions are imposed :
σpσ +
e†e+ d†d = 1, and f †σfσ = p
σpσ + d
†d(σ =↑, ↓). Therefore, the Hamiltonian can be rewritten
as the following effective Hamiltonian in terms of the auxiliary boson e, pσ, d and the pesudo-
fermion operators fσ:
Heff =
mσcmσ + 1/2
[tm(c
m+1↓cm↑ + c
m−1↑cm↓) + h.c.]
σfσ + Ud
σcmσ + h.c.) + λ
p†σpσ + e
†e+ d†d− 1)
λ(2)σ (f
σfσ − p
σpσ − d
†d) , (6)
where the constraints are incorporated by the Lagrange multipliers λ(1) and λ(2)σ . The first
constraint can be interpreted as a completeness relation of the Hilbert space on the impurity
level, and the second one equates the two ways of counting the fermion occupancy for a given
spin. In the framework of the finite-U slave boson mean field theory[25, 26], the slave boson
operators e, pσ, d and the parameter zσ are replaced by real c numbers. Thus the effective
Hamiltonian is given as
HMFeff =
mσcmσ + 1/2
[tm(c
m+1↓cm↑ + c
m−1↑cm↓) + h.c.]
ǫ̃dσf
σfσ +
(t̃Dσf
σcmσ + h.c.) + Eg , (7)
where t̃Dσ = tDzσ represents the renormalized tunnel coupling between the impurity and
the mesoscopic ring. zσ can be regarded as the wave function renormalization factor. ǫ̃dσ =
σ is the renormalized impurity level and Eg = λ
2+d2−1)−
d2) + Ud2 is an energy constant.
In this mean field approximation the Hamiltonian is essentially that of a non-interacting
system, hence the single particle energy levels can be calculated by numerical diagonalization
of the Hamiltonian matrix. Then the ground state of this system |ψ0 > can be constructed
by adding electrons to the lowest unoccupied energy levels consecutively . By minimizing
the ground state energy with respect to the variational parameters a set of self-consistent
equations can be obtained as in Ref.[27,28], and they can be applied to determine the
variational parameters in the effective Hamiltonian.
III. THE PERSISTENT CHARGE CURRENT AND SPIN CURRENT
In this section we will present the results of our calculation of the persistent charge current
and spin current circulating the mesoscopic ring. Since there is still some controversial in the
literature for the definition of the spin current operator in the ring system with SO coupling
term[30]. We give both the formula of charge and spin currents used in this paper explicitly.
It is easy to obtain that the ϕ component of electron velocity operator in this SO coupled
ring is
[2∆(−i
+ φ) + αR(σx cosϕ+ σy sinϕ)] . (8)
Thereby the charge current operator is define as Î = −evϕ, and in terms of creation and
annihilation operator it can be written as
Î = −
c†mσcmσ(m+ φ) + αR
m+1↓cm↑ + c
m−1↑cm↓)] . (9)
At zero temperature, the persistent charge current is given by the expectation value of the
above charge current operator in the ground state, I = 1
< ψ0|Î|ψ0 >, and it can also be
calculated from the expression
I = −c
< ψ0|
|ψ0 > , (10)
where Egs is the ground state energy.
In Fig.1 the persistent charge current vs. the enclosed magnetic flux is plotted for a
set of values for the SO coupling strength. Here we have taken the model parameters
∆ = 0.01, tD = 0.3, U = 2.0 and the total number of electrons N is around 100. In this
case one can obtain the Fermi energy of the system EF = 6.25 and the level spacing δ = 0.5
around the Fermi surface. We consider the energy level of the Anderson impurity is well
below the Fermi energy( with ǫd − EF = −1.0), therefore the Anderson impurity is in the
Kondo regime. One can see in Fig.1 that the characteristic features of persistent charge
current depends on the parity of the total number of electrons(N), and can be distinguished
by two cases with N odd and N even. This is attributed to the different occupation patterns
of the highest occupied single particle energy level in the mean field effective Hamiltonian.
The persistent charge current for the system with N +2 electrons is different from that with
N electrons by a π phase shift IN+2(φ) = IN(φ+ π). In case (I) where the electron number
is odd(N = 4n− 1 and N = 4n + 1), one electron is almost localized on the impurity level
and forming a singlet with electron cloud in the conducting ring. This phenomena leads to
the well known Kondo effect. Fig.1 shows that the Kondo effect decreases the magnitude
of the persistent charge current, and also makes its curve shape resemble sinusoidal. In the
presence of finite SO coupling(αR < ∆), the spin-up and spin-down electrons are coupled and
it causes the splitting of the twofold degenerated energy levels in the effective Hamiltonian.
It turns out that the Kondo effect is suppressed and the abrupt jumps of the persistent
charge current with similarity to that of ideal ring case appears. It is explained in Ref.[14]
that the jumps of the persistent charge current in the case of odd number of electrons are
due to a crossing of levels with opposite spin. In case (II) where N is even (N = 4n and
N = 4n+2), The Kondo effect is manifested that the magnitude of persistent charge current
is significantly suppressed compared with ideal ring case and the rounding of the jumps of
persistent charge current due to the level crossing. In the presence of finite SO coupling, the
persistent charge current decreases with increasing the SO coupling strength when αR < ∆.
Fig.2 displays the persistent charge current as a function of the SO coupling strength
αR at different enclosed magnetic flux. The persistent charge current exhibits oscillations
with increasing the value of αR for both the systems with even or odd number of electrons.
Therefore by tuning the SO coupling strength, the magnetic response of this system can
change from paramagnetic to diamagnetic and vice versa. It indicates that SO coupling
can play a important role in electron transport in this mesoscopic ring. The curve of the
persistent charge current for odd number of electrons shows discontinuity in its derivation,
this can be attributed the level crossing in the energy spectrum by changing αR. It is also
noted that the position of this discontinuity for odd N also corresponds to the peak or valley
in even N case.
Since the electron has the spin degree of freedom as well as the charge, the electron
motion in the ring may give rise to a spin current besides the charge current. Now we turn
to study the persistent spin current in the ground state. The spin current operator is defined
by Ĵv = (v
ϕσv + σvv
ϕ)/2, which can be written explicitly as
Ĵv =
{2∆(−i
+ φ)σv +
[(σx cosϕ+ σy sinϕ)σv + h.c.]} , (11)
Therefore the three component of spin current operator in terms of creation and annihi-
lation operators are given by
Ĵz =
m↑cm↑ − c
m↓cm↓)(m+ φ)] , (12)
Ĵx =
m↑cm↓ + c
m↓cm↑)(m+ φ) +
m+1σ + c
m−1σ)cmσ] , (13)
Ĵy =
[−2i∆
m↑cm↓ − c
m↓cm↑)(m+ φ)− i
m+1σ − c
m−1σ)cmσ] , (14)
The expectation value of the spin current Jv =
< ψ0|Ĵv|ψ0 >.
In our calculation we find that only the z component of the spin current is nonzero in the
ground state. Fig.3 shows the persistent spin current Jz vs. magnetic flux at different SO
coupling strength. The persistent spin current is a periodic function of the magnetic flux
φ, which has the even parity symmetry Jz(−φ) = Jz(φ) and also an additional symmetry
Jz(φ) = Jz(π−φ). It is noted that the persistent spin current has quite different dependence
behaviors on magnetic flux compared with the persistent charge current in Fig.1. In the
presence of finite SO coupling, the persistent spin current is nonzero both for the systems
with odd N and even N at zero magnetic flux, it indicates that a persistent spin current can
be induced solely by SO interaction without accompany a charge current. This phenomena
is also shown in Ref.[17] where a SO coupling/normal hybrid ring was considered.
In Fig.4 the persistent spin current Jz as a function of SO coupling strength is plotted.
In the absence of SO coupling αR = 0, the persistent spin current is exactly zero for both
even and odd number electron system. In the presence of SO coupling, The persistent spin
current becomes nonzero and shows oscillations with increasing αR. It can change from
positive to negative values or vice versa by tuning the SO coupling strength. The sign of the
persistent spin current also shows dependence on the enclosed magnetic flux. For the system
with odd N , there is abrupt jumps in the curve of persistent spin current at certain value of
αR, the reason for the jump is the same as that in the charge current, and is due to the level
crossing in the energy spectrum. It is noted that the position of the jump coincides with
that of the persistent charge current. This kind of characteristic feature of the persistent
currents might provide a useful way to detect the SO coupling effects in semiconductor ring
system.
IV. CONCLUSIONS
In summary, we have investigated the Rashba SO coupling effect on the persistent charge
current and spin current in a mesoscopic ring with an Anderson impurity. The Anderson
impurity leads to the Kondo effect and decreases the amplitude of the persistent charge and
spin current in the ring. In the semiconducting ring with SO interaction, the persistent
charge current changes significantly by tuning the SO coupling strength, e.g. from the
paramagnetic to diamagnetic current. Besides the persistent charge current, there also
exists a persistent spin current, which also oscillates with the SO coupling strength. It is
shown that at zero magnetic flux a persistent spin current can exist even without the charge
current. Since the persistent spin current can generate an electric field[31], one might expect
that experiments on semiconductor ring with Rashba SO coupling can detect the persistent
spin current.
Acknowledgments
This project is supported by the National Natural Science Foundation of China, the
Shanghai Pujiang Program, and Program for New Century Excellent Talents in University
(NCET).
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http://arxiv.org/abs/cond-mat/0605748
0.0 0.2 0.4 0.6 0.8 1.0
0.0 0.2 0.4 0.6 0.8 1.0
0.0 0.2 0.4 0.6 0.8 1.0
(d)(b)
0.0 0.2 0.4 0.6 0.8 1.0
FIG. 1: The persistent charge current vs. magnetic flux for a set of values for the spin-orbit
coupling strength(αR/∆ = 0.0(solid line),0.5(dashed line), 0.7(dotted line),1.0(dash-dotted line)).
The total number of electrons N = 99 (a), 100(b), 101(c), 102(d). We take the other parameters
∆ = 0.01, td = 0.3, ǫd − EF = −1.0, U = 2.0 in the calculation. The persistent charge current is
measured in units of I0 = eN∆.
0 1 2 3 4
-0.15
-0.10
-0.05
0 1 2 3 4
-0.10
-0.05
(c)(a)
0 1 2 3 4
-0.10
-0.05
0 1 2 3 4
-0.10
-0.05
FIG. 2: The persistent charge current as a function of the spin-orbit coupling strength. The
magnetic flux (Φ/Φ0 = 0.125(solid line),0.25(dashed line), 0.375(dotted line)).
0.0 0.2 0.4 0.6 0.8 1.0
-0.15
-0.10
-0.05
0.0 0.2 0.4 0.6 0.8 1.0
-0.10
-0.05
0.0 0.2 0.4 0.6 0.8 1.0
-0.05
0.0 0.2 0.4 0.6 0.8 1.0
-0.05
FIG. 3: FIG.3: The persistent spin current Jz vs. magnetic flux for a set of values for the spin-
orbit coupling strength( with αR/∆ = 0.5(solid line),0.7(dashed line), 1.0(dotted line)). The panel
(a), (b), (c) and (d) corresponds the system with total number of electrons N = 99, 100, 101, 102,
respectively. The persistent spin current is measured in units of J0 = N∆, and we have taken the
other parameter values the same as that in Fig.1.
0 1 2 3 4
-0.10
-0.05
0 1 2 3 4
-0.10
-0.05
0 1 2 3 4
-0.10
-0.05
0 1 2 3 4
-0.10
-0.05
FIG. 4: FIG.4: The persistent spin current Jz as a function of the spin-orbit coupling strength.
The magnetic flux takes the value (Φ/Φ0 = 0.0(solid line),0.125(dashed line), 0.25(dotted line),
0.5(dash-dotted line)).
introduction
Mesoscopic ring with an Anderson impurity
the persistent charge current and spin current
conclusions
Acknowledgments
References
|
0704.0320 | Probability distributions generated by fractional diffusion equations | FRACALMO PRE-PRINT www.fracalmo.org
Probability distributions generated by
fractional diffusion equations1
Francesco MAINARDI(1), Paolo PARADISI(2) and Rudolf GORENFLO(3)
(1) Department of Physics, University of Bologna, and INFN,
Via Irnerio 46, I-40126 Bologna, Italy.
[email protected] [email protected]
(2) ISAC: Istituto per le Scienze dell’Atmosfera e del Clima del CNR,
Strada Provinciale Lecce-Monteroni Km 1.200, I-73100 Lecce, Italy.
[email protected]
(3) Department of Mathematics and Computer Science,
Freie Universität Berlin, Arnimallee 3, D-14195 Berlin, Germany.
[email protected]
Contents
Abstract . . . . . . . . . . . . . . . . . . . . . . . . p. 2
1. Introduction . . . . . . . . . . . . . . . . . . . . . . . p. 2
2. The Standard Diffusion Equation . . . . . . . . . . . . . . p. 4
3. The Time-Fractional Diffusion Equation . . . . . . . . . . . p. 8
4. The Cauchy Problem for the Time-Fractional Diffusion Equation p.10
5. The Signalling Problem for the Time-Fractional Diffusion Equation p.13
6. The Cauchy Problem for the Symmetric Space-Fractional
Diffusion Equation . . . . . . . . . . . . . . . . . . . . p.15
7. Conclusions . . . . . . . . . . . . . . . . . . . . . . . p.21
A. The Riemann-Liouville Fractional Calculus . . . . . . . . . p.22
B. The Stable Probability Distributions . . . . . . . . . . . . p.31
References . . . . . . . . . . . . . . . . . . . . . . . p.41
1This paper is based on an invited talk given by Francesco Mainardi at the International
Workshop on Econophysics held at Bolyai College, Eötvös University, Budapest, on
July 21-27, 1997. The paper was originally edited as a contribution for the book J.
Kertesz and I. Kondor (Editors), Econophysics: an Emerging Science, Kluwer
Academic Publishers, Dordrecht (NL) that should contain selected papers presented at
that Workshop and should have appeared in 1998 or 1999. Unfortunately the book was
not published. The present e-print is a revised version (with up-date annotations and
references) of that unpublished contribution, but essentially represents our knowledge of
that early time.
http://arXiv.org/abs/0704.0320v1
Abstract
Fractional calculus allows one to generalize the linear, one-dimensional,
diffusion equation by replacing either the first time derivative or the second
space derivative by a derivative of fractional order. The fundamental
solutions of these generalized diffusion equations are shown to provide
probability density functions, evolving on time or variable in space, which
are related to the peculiar class of stable distributions. This property is
a noteworthy generalization of what happens for the standard diffusion
equation and can be relevant in treating financial and economical problems
where the stable probability distributions are known to play a key role.
1 Introduction
Non-Gaussian probability distributions are becoming more common as data
models, especially in economics where large fluctuations are expected. In
fact, probability distributions with heavy tails are often met in economics
and finance, which suggests to enlarge the arsenal of possible stochastic
models by non-Gaussian processes. This conviction started in the early
sixties after the appearance of a series of papers by Mandelbrot and
his associates, who point out the importance of non-Gaussian probability
distributions, formerly introduced by Pareto and Lévy, and related scaling
properties, to analyse economical and financial variables, as reported in
the recent book by Mandelbrot (1997). Some examples of such variables
are common stock prices changes, changes in other speculative prices, and
interest rate changes. In this respect many works by different authors have
recently appeared, see e.g. the recent books by Bouchaud & Potter (1997),
Mantegna & Stanley (1998) and the references therein quoted.
It is well known that the fundamental solution (or Green function) of
the Cauchy problem for the standard linear diffusion equation provides at
any time the probability density function (pdf) in space of the Gauss (or
normal) law. This law exhibits all moments finite thanks to its exponential
decay at infinity. In particular, the space variance of the Green function
is proportional to the first power of time, a noteworthy property that
can be understood by means of an unbiased random walk model for the
Brownian motion, see e.g. Feller (1957). Less known is the property for
which the fundamental solution of the Signalling problem for the same
diffusion equation, provides at any position a unilateral pdf in time, known
as Lévy law, using the terminology of Feller (1966-1973). Because of its
algebraic decay at infinity as t−3/2 , this law has all moments of integer
order divergent, and consequently its expectation value and variance are
infinite.
Both the Gauss and Lévy laws belong to the general class of stable probability
distributions, which are characterized by an index α (0 < α ≤ 2), called
index of stability or characteristic exponent. In particular, the index of the
Gauss law is 2 , whereas that of the Lévy law is 1/2 .
In this paper we consider two different generalizations of the diffusion
equation by means of fractional calculus, which allows us to replace either the
first time derivative or the second space derivative by a suitable fractional
derivative. Correspondingly, the generalized equation will be referred to
as the time-fractional diffusion equation or the symmetric, space-fractional
diffusion equation. Here we show how the fundamental solutions of this
equation for the Cauchy and Signalling problems provide probability density
functions related to certain stable distributions, so providing a natural
generalization of what occurs for the standard diffusion equation.
The plan of the paper is as follows. First of all, for the sake of convenience
and completeness, we provide the essential notions of Riemann-Liouville
Fractional Calculus and Lévy Stable Probability Distributions in Appendix
A and B, respectively.
In Section 2, we recall the basic results for the standard diffusion
equation concerning the fundamental solutions of the Cauchy and Signalling
problems. In particular we provide the derivation of these solutions by the
Fourier and Laplace transforms and the interpretation in terms of Gauss
and Lévy stable pdf , respectively.
In Section 3, we consider the time-fractional diffusion equation and we
formulate for it the basic Cauchy and Signalling problems to be treated in the
subsequent two sections. Here we adopt the Riemann-Liouville approach to
Fractional Calculus, and the related definition for the Caputo time-fractional
derivative of a causal function of time.
In Section 4, we solve the Cauchy problem for the time-fractional diffusion
equation by using the technique of Fourier transform and we derive the
corresponding fundamental solution in terms of a special function of Wright
type in the similarity variable. In this case the solution can be interpreted
as a noteworthy symmetric pdf in space with all moments finite, evolving
in time. In particular, its space variance turns out to be proportional to a
power of time equal to the order of the time-fractional derivative.
In Section 5, we derive the fundamental solution for the Signalling problem
of the time-fractional diffusion equation by using the technique of Laplace
transform. In this case the solution, still expressed in terms of a special
function of Wright type, can be interpreted as a unilateral stable pdf in
time, depending on position, with index of stability given by half of the
order of the time-fractional derivative.
In Section 6, we consider the symmetric, space-fractional diffusion equation.
Here we adopt the Riesz approach to Fractional Calculus, and the related
definition for the symmetric space-fractional derivative of a function of a
single space variable. Here we treat the Cauchy problem by technique
of Fourier transform and we derive the series representation of the
corresponding Green function. In this case the fundamental solution is
interpreted in terms of a symmetric stable pdf in space, evolving in time,
with index of stability given by the order of the space-fractional derivative.
To approximate such evolution we propose a random walk model, discrete
in space and time, which is based on the Grünwald-Letnikov approximation
of the fractional derivative.
Finally, Section 7 is devoted to conclusions and remarks on related work.
2 The standard diffusion equation
For the standard diffusion equation we mean the linear partial differential
equation
u(x, t) = D
u(x, t) , u = u(x, t) , (2.1)
where D denotes a positive constant with the dimensions L2 T−1 , x and t
are the space-time variables, and u = u(x, t) is the field variable, which is
assumed to be a causal function of time, i.e. vanishing for t < 0 .
The typical physical phenomenon related to such an equation is the heat
conduction in a thin solid rod extended along x , so the field variable u is
the temperature.
In order to guarantee the existence and the uniqueness of the solution,
we must equip (1.1) with suitable data on the boundary of the space-time
domain. The basic boundary-value problems for diffusion are the so-called
Cauchy and Signalling problems. In the Cauchy problem, which concerns
the space-time domain −∞ < x < +∞ , t ≥ 0 , the data are assigned at
t = 0+ on the whole space axis (initial data). In the Signalling problem,
which concerns the space-time domain x ≥ 0 , t ≥ 0 , the data are assigned
both at t = 0+ on the semi-infinite space axis x > 0 (initial data) and at
x = 0+ on the semi-infinite time axis t > 0 (boundary data); here, as mostly
usual, the initial data are assumed to be vanishing.
Denoting by g(x) and h(t) two given, sufficiently well-behaved functions, the
basic problems are thus formulated as following:
a) Cauchy problem
u(x, 0+) = g(x) , −∞ < x < +∞ ; u(∓∞, t) = 0 , t > 0 ; (2.2a)
b) Signalling problem
u(x, 0+) = 0 , x > 0 ; u(0+, t) = h(t) , u(+∞, t) = 0 , t > 0 . (2.2b)
Hereafter, for both the problems, we derive the classical results which will be
properly generalized for the fractional diffusion equation in the subsequent
sections.
Let us begin with the Cauchy problem. It is well known that this initial value
problem can be easily solved making use of the Fourier transform and its
fundamental solution can be interpreted as a Gaussian pdf in x. Adopting
the notation g(x) ÷ ĝ(κ) with κ ∈ R and
ĝ(κ) = F [g(x)] =
e+iκx g(x) dx ,
g(x) = F−1 [ĝ(κ)] = 1
e−iκx ĝ(κ) dκ ,
the transformed solution satisfies the ordinary differential equation of the
first order (
+ κ2 D
û(κ, t) = 0 , û(κ, 0+) = ĝ(κ) , (2.3)
and consequently it turns out to be
û(κ, t) = ĝ(κ) e−κ
2 D t . (2.4)
Then, introducing
Gdc (x, t) ÷ Ĝdc (κ, t) = e−κ
2 D t , (2.5)
where the upper index d refers to (standard) diffusion, the required solution,
obtained by inversion of (2.4), can be expressed in terms of the space
convolution u(x, t) =
−∞ Gdc (ξ, t) g(x − ξ) dξ , where
Gdc (x, t) =
t−1/2 e−x
2/(4D t) . (2.6)
Here Gdc (x, t) represents the fundamental solution (or Green function) of
the Cauchy problem, since it corresponds to g(x) = δ(x) . It turns out
to be a function in x , even and normalized, i.e. Gdc (x, t) = Gdc (|x|, t) and∫ +∞
−∞ Gdc (x, t) dx = 1 . We also note the identity
|x| Gdc (|x|, t) =
Md(ζ) , (2.7)
where ζ = |x|/(
D t1/2) is the well-known similarity variable and
Md(ζ) =
2/4 . (2.8)
We note that Md(ζ) satisfies the normalization condition
d(ζ) dζ = 1 .
The interpretation of the Green function (2.6) in probability theory is
straightforward since we easily recognize
Gdc (x, t) = pG(x;σ) :=
2/(2σ2) , σ2 = 2D t , (2.9)
where pG(x;σ) denotes the well-known Gauss or normal pdf spread out
over all real x (the space variable), whose moment of the second order, the
variance, is σ2 . The associated cumulative distribution function (cdf) is
known to be
PG(x;σ) :=
′;σ) dx′ =
1 + erf
, (2.10)
where erf (z) := (2/
0 exp (−u2) du denotes the error function.
Furthermore, the moments of even order of the Gauss pdf turn out to be∫ +∞
2n pG(x;σ) dx = (2n − 1)!!σ2n , so
x2n Gdc (x, t) dx = (2n − 1)!! (2D t)n , n = 1, 2, . . . . (2.11)
Let us now consider the Signalling problem. This initial-boundary value
problem can be easily solved by making use of the Laplace transform.
Adopting the notation h(t) ÷ h̃(s) with s ∈ C and
h̃(s) = L [h(t)] =
e−st h(t) dt ,
h(t) = L−1
h̃(t)
est h̃(s) ds ,
where Br denotes the Bromwich path, the transformed solution of the
diffusion equation satisfies the ordinary differential equation of the second
order
ũ(x, s) = 0 , ũ(0+, s) = h̃(s) , ũ(+∞, s) = 0 . (2.12)
and consequently it turns out to be
ũ(x, s) = h̃(s) e−(x/
D) s1/2 . (2.13)
Then introducing
Gds (x, t) ÷ G̃ds (x, s) = e−(x/
D) s1/2 , (2.14)
the required solution, obtained by inversion of (2.13), can be expressed in
terms of the time convolution, u(x, t) =
0 Gds (x, τ)h(t − τ) dτ , where
Gds (x, t) =
t−3/2 e−x
2/(4D t) . (2.15)
Here Gds (x, t) represents the fundamental solution (or Green function) of the
Signalling problem, since it corresponds to h(t) = δ(t) . We note that
Gds (x, t) = pLS(t;µ) :=
2π t3/2
e−µ/(2t) , t ≥ 0 , µ = x
, (2.16)
where pLS(t;µ) denotes the one-sided Lévy-Smirnov pdf spread out over all
non negative t (the time variable). The associated cdf is, see e.g. Feller
(1966-1971) and Prüss (1993),
PL(t;µ) :=
′;µ) dt′ = erfc
= erfc
, (2.17)
where erfc (z) := 1 − erf (z) denotes the complenatary error function.
The Lévy-Smirnov pdf has all moments of integer order infinite, since it
decays at infinity as t−3/2 . However, we note that the absolute moments of
real order ν are finite only if 0 ≤ ν < 1/2 . In particular, for this pdf the mean
is infinite, for which we can take the median as expectation value. From
PLs(tmed;µ) = 1/2 , it turns out that tmed ≈ 2µ , since the complementary
error function gets the value 1/2 as its argument is approximatively 1/2.
We note that in the common domain x > 0 , t > 0 the Green functions of
the two basic problems satisfy the identity
xGdc (x, t) = tGds (x, t) , (2.18)
that we refer to as the reciprocity relation between the two fundamental
solutions of the diffusion equation. Furthermore, in view of (2.7) and (2.18)
we recognize the role of the function of the similarity variable, Md(ζ) ,
in providing the two fundamental solutions; we shall refer to it as to the
normalized auxiliary function of the diffusion equation for both the Cauchy
and Signalling problems.
3 The time-fractional diffusion equation
By the time-fractional diffusion equation we mean the linear evolution
equation obtained from the classical diffusion equation by replacing the first-
order time derivative by a fractional derivative (in the Caputo sense) of order
α with 0 < α ≤ 2. In our notation it reads
, u = u(x, t) , 0 < α ≤ 2 , (3.1)
where D denotes a positive constant with the dimensions L2 T−α . From
Appendix A we recall the definition of the Caputo fractional derivative of
order α > 0 for a (sufficiently well-behaved) causal function f(t) , see (A.9),
Dα∗ f(t) :=
Γ(m − α)
(t − τ)m−α f (m)(τ) dτ , (3.2)
where m = 1, 2, . . . , and 0 ≤ m − 1 < α ≤ m . According to (3.2) we thus
need to distinguish the cases 0 < α ≤ 1 and 1 < α ≤ 2 . In the the latter case
(3.1) may be seen as a sort of interpolation between the standard diffusion
equation and the standard wave equation. Introducing
Φλ(t) :=
tλ−1+
, λ > 0 , (3.3)
where the suffix + is just denoting that the function is vanishing for t < 0 ,
we easily recognize that the equation (3.1) assumes the explicit forms :
if 0 < α ≤ 1 ,
Φ1−α(t) ∗
Γ(1 − α)
(t − τ)−α
dτ = D
; (3.4)
if 1 < α ≤ 2 ,
Φ2−α(t) ∗
Γ(2 − α)
(t − τ)1−α
dτ = D ∂
. (3.5)
Extending the classical analysis for the standard diffusion equation (2.1) to
the above integro-differential equations (3.4-5), the Cauchy and Signalling
problems are thus formulated as in equations (2.2), i.e.
a) Cauchy problem
u(x, 0+) = g(x) , −∞ < x < +∞ ; u(∓∞, t) = 0 , t > 0 ; (3.6a)
b) Signalling problem
u(x, 0+) = 0 , x > 0 ; u(0+, t) = h(t) , u(+∞, t) = 0 , t > 0 . (3.6b)
However, if 1 < α ≤ 2 , the presence in (3.5) of the second order time
derivative of the field variable requires to specify the initial value of the first
order time derivative ut(x, 0
+) , since in this case two linearly independent
solutions are to be determined. To ensure the continuous dependence of our
solution on the parameter α also in the transition from α = 1− to α = 1+ ,
we agree to assume ut(x, 0
+) = 0 .
We recognize that our fractional diffusion equation (3.1), when subject to
the conditions (3.6), is equivalent to the integro-differential equation
u(x, t) = g(x) +
(t − τ)α−1
dτ , (3.7)
where 0 < α ≤ 2 . Such integro-differential equation has been investigated
by several authors, including Schneider & Wyss (1989), Fujita (1990), Prüss
(1993) and Engler (1997).
In view of our subsequent analysis we find it convenient to put
, 0 < ν < 1 . (3.8)
In fact the analysis of the time-fractional diffusion equation turns out to
be easier if we adopt as a key parameter the half of the order of the
time-fractional derivative. In future we shall provide the symbol α with
other relevant meanings, as the index of stability of a stable probability
distribution or the order of the space derivative in the space-fractional
diffusion equation.
Henceforth, we agree to insert the parameter ν in the field variable, i.e.
u = u(x, t; ν) . By denoting the Green functions of the Cauchy and Signalling
problems by Gc(x, t; ν) and Gs(x, t; ν) , respectively, the solutions of the two
basic problems are obtained by a space or time convolution, u(x, t; ν) =∫ +∞
−∞ Gc(ξ, t; ν) g(x−ξ) dξ , u(x, t; ν) =
0 Gs(x, τ ; ν)h(t−τ) dτ , respectively.
It should be noted that Gc(x, t; ν) = Gc(|x|, t; ν) , since the Green function
turns out to be an even function of x .
In the following two sections we shall compute the two fundamental solutions
with the same techniques (based on Fourier and Laplace transforms) used
for the standard diffusion equation and we shall provide their interpretation
in terms of probability distributions. Most of the presented results are based
on the papers by Mainardi (1994), (1995), (1996), (1997) and by Mainardi
& Tomirotti (1995), (1997).
4 The Cauchy problem for the time-fractional
diffusion equation
For the fractional diffusion equation (3.1) subject to (3.6a) the application
of the Fourier transform leads to the ordinary differential equation of order
α = 2ν ,
+ κ2 D
û(κ, t; ν) = 0 , û(κ, 0+; ν) = ĝ(κ) , (4.1)
Using the results of Appendix A, see (A.22-30), the transformed solution is
û(κ, t; ν) = ĝ(κ)E2ν
−κ2 D t2ν
, (4.2)
where E2ν(·) denotes the Mittag-Leffler function of order 2ν , and conse-
quently for the Green function we have
Gc(x, t; ν) = Gc(|x|, t; ν) ÷ Ĝc(k, t; ν) = E2ν
−κ2D t2ν
. (4.3)
Since the Green function is a real and even function of x, its (exponential)
Fourier transform can be expressed in terms of the cosine Fourier transform
and thus is related to its spatial Laplace transform as follows
Ĝc(k, t; ν) = 2
Gc(x, t; ν) cos κx dx =
G̃c(s, t; ν)
s=+ik
+ G̃c(s, t; ν)
s=−ik
(4.4)
Indeed, a split occurs also in (4.3) according to the duplication formula for
the Mittag-Leffler function, see (A.26),
Ĝc(k, t; ν) = E2ν(−κ2 D t2ν) =
[Eν(+iκ
D tν) + Eν(−iκ
D tν)]/2 .
(4.5)
When ν 6= 1/2 the inversion of the Fourier transform in (4.5) cannot be
obtained by using a standard table of Fourier transform pairs; however, for
any ν ∈ (0, 1) such inversion can be achieved by appealing to the Laplace
transform pair (A.37) with r = |x| , and s = ±iκ . In fact, taking into
account the scaling property of the Laplace transform, we obtain from (4.5)
and (A.37)
Gc(|x|, t; ν) =
( |x|√
, (4.6)
where M(ζ; ν) is the special function of Wright type, defined by (A.31-33),
, (4.7)
the similarity variable. We note the identity
|x| Gc(|x|, t; ν) =
M(ζ; ν) , (4.8)
which generalizes to the time-fractional diffusion equation the identity (2.7)
of the standard diffusion equation. Since
0 M(ζ; ν) dζ = 1 , see (A.40),
the function M(ζ; ν) is the normalized auxiliary function of the fractional
diffusion equation.
We note that for the time-fractional diffusion equation the fundamental
solution of the Cauchy problem is still a bilateral symmetric pdf in x (with
two branches, for x > 0 and x < 0 , obtained one from the other by
reflection), but is no longer of Gaussian type if ν 6= 1/2 . In fact, for large
|x| each branch exhibits an exponential decay in the ”stretched” variable
|x|1/(1−ν) as can be derived from the asymptotic representation (A.36) of the
auxiliary function M(·; ν) . In fact, by using (4.7-8) and (A.36), we obtain
Gc(x, t; ν) ∼ a∗(t) |x|(ν−1/2)/(1−ν) exp
−b∗(t)|x|1/(1−ν)
, (4.9)
as |x| → ∞ , where a∗(t) and b∗(t) are certain positive functions of time.
Furthermore, the exponential decay in x provided by (4.9) ensures that all
the absolute moments of positive order of Gc(x, t; ν) are finite. In particular,
using (4.8) and (A.39) it turns out that the moments (of even order) are
x2n Gc(x, t; ν) dx =
Γ(2n + 1)
Γ(2νn + 1)
(Dt2ν)n , n = 0 , 1 , 2 , . . . (4.10)
The formula (4.10) provides a generalization of the corresponding formula
(2.11) valid for the standard diffusion equation, ν = 1/2 . Furthermore, we
recognize that the variance associated to the pdf is now proportional to Dt2ν ,
which for ν 6= 1/2 implies a phenomenon of anomalous diffusion. According
to a usual terminology in statistical mechanics, the anomalous diffusion is
said to be slow if 0 < ν < 1/2 and fast if 1/2 < ν < 1 .
In Figure 1, as an example, we compare versus |x| , at fixed t , the
fundamental solutions of the Cauchy problem with different ν (ν =
1/4 , 1/2 , 3/4 ). We consider the range 0 ≤ |x| ≤ 4 and assume D = t = 1 .
0 1 2 3 4
Figure 1: The Cauchy problem for the time-fractional diffusion equation.
The fundamental solutions versus |x| with a) ν = 1/4 , b) ν = 1/2 , c)
ν = 3/4 .
We note the different behaviour of the pdf in the cases of slow diffusion (ν =
1/4 ) and fast diffusion (ν = 3/4 ) with respect to the Gaussian behaviour
of the standard diffusion (ν = 1/2). In the limiting cases ν = 0 and ν = 1
we have
Gc(x, t; 0) =
e−|x|
, Gc(x, t; 1) =
δ(x −
D t) + δ(x +
. (4.11)
We also recognize from the appendix B that for 1/2 ≤ ν < 1 any branch
of the fundamental solution is proportional to the corresponding positive
branch of an extremal stable pdf with index of stability α = 1/ν , which
exhibits an exponential decay at infinity. In fact, applying (B.29) with
α = 1/ν and y = ζ = |x|/(
Dtν) , from (4.7-8) we obtain
Gc(|x|, t; ν) =
|x|/(
D tν) ; − (2 − 1/ν)
· p1/ν (|x|; +1, 1, 0) , 1 < 1/ν ≤ 2 .
(4.12)
We also note that the stable distribution in (4.12) satisfies the condition
p1/ν (x; +1, 1, 0) dx = ν , 1 < 1/ν ≤ 2 . (4.13)
5 The Signalling problem for the time-fractional
diffusion equation
For the fractional diffusion equation (3.1) subject to (3.6b) the application
of the Laplace transform leads to the ordinary differential equation of order
ũ(x, s; ν) , ũ(0+, s; ν) = h̃(s) , ũ(+∞, s; ν) = 0 . (5.1)
Thus the transformed solution reads
ũ(x, s; ν) = h̃(s) e−(x/
D) sν , (5.2)
so for the Green function we have
Gs(x, t; ν) ÷ G̃s(x, s; ν) = e−(x/
D) sν . (5.3)
When ν 6= 1/2 the inversion of this Laplace transform cannot be obtained by
looking in a standard table of Laplace transform pairs. Also here we appeal
to a Laplace transform pair related to the Wright-type function M(ζ; ν). In
fact, using (A.40) with r = t , and taking into account the scaling property
of the Laplace transform, we obtain
Gs(x, t; ν) = ν
D t1+ν
. (5.4)
Introducing the similarity variable ζ = x/(
Dtν) , we recognize the identity
tGs(x, t; ν) = ν ζ M(ζ; ν) , (5.5)
which is the counterpart for the Signalling problem of the identity (4.8) valid
for the Cauchy problem.
Comparing (5.5) with (4.8) we obtain the reciprocity relation between the
two fundamental solutions of the time-fractional diffusion equation, in the
common domain x > 0 , t > 0 ,
2ν xGc(x, t; ν) = tGs(x, t; ν) . (5.6)
The interpretation of Gs(x, t; ν) as a one-sided stable pdf in time is
straightforward: in this respect we need to apply (B.28), with index of
stability α = ν and variable y = ζ−1/ν = t (
D/x)1/ν , in (5.5). We obtain
Gs(x, t; ν) =
; − ν
= pν (t; −1, 1, 0) . (5.7)
In Figure 2, as an example, we compare versus t , at fixed x , the fundamental
solutions of the Signalling problem with different ν (ν = 1/4 , 1/2 , 3/4 ). We
consider the range 0 ≤ t ≤ 3 and assume D = x = 1 .
We note the different behaviour of the pdf in the cases of slow diffusion
(ν = 1/4 ) and fast diffusion (ν = 3/4 ) with respect to the Lévy pdf for the
standard diffusion (ν = 1/2). In the limiting cases ν = 0 , 1 , we have
Gs(x, t; 0) = δ(t) , Gs(x, t; 1) = δ(t − x/
D) . (5.8)
0 1 2 3
Figure 2: The Signalling problem for the time-fractional diffusion equation.
The fundamental solutions versus t with a) ν = 1/4 , b) ν = 1/2 , c)
ν = 3/4 .
6 The Cauchy problem for the symmetric space-
fractional diffusion equation
The symmetric space-fractional diffusion equation is obtained from the
classical diffusion equation by replacing the second-order space derivative by
a symmetric space-fractional derivative (explained below) of order α with
0 < α ≤ 2 . In our notation we write this equation as
∂|x|α
, u = u(x, t;α) , x ∈ R , t ∈ R+0 , 0 < α ≤ 2 , (6.1)
where D is a positive coefficient with the dimensions Lα T−1 . The
fundamental solution for the Cauchy problem, Gc(x, t;α) is the solution of
(6.1), subject to the initial condition u(x, 0+;α) = δ(x) .
The symmetric space-fractional derivative of any order α > 0 of a sufficiently
well-behaved function φ(x) , x ∈ R , may be defined as the pseudo-
differential operator characterized in its Fourier representation by
d|x|α
φ(x) ÷ −|κ|α φ̂(κ) , x , k ∈ R , α > 0 . (6.2)
According to a usual terminology, −|κ|α is referred to as the symbol of our
pseudo-differential operator, the symmetric space-fractional derivative, of
order α . Here, we have adopted the notation introduced by Zaslavski, see
e.g. Saichev & Zaslavski (1997).
In order to properly introduce this kind of fractional derivative we need
to consider a peculiar approach to fractional calculus different from the
Riemann-Liouville one, already treated in Appendix A. This approach is
indeed based on the so-called Riesz potentials (or integrals), that we prefer
to consider later.
At first, let us see how things become highly transparent by using an
heuristic argument, originally due to Feller (1952). The idea is to start
from the positive definite differential operator
A := −
÷ κ2 = |κ|2 , (6.3)
whose symbol is |κ|2 , and form positive powers of this operator as pseudo-
differential operators by their action in the Fourier-image space, i.e.
Aα/2 :=
= |κ|α α > 0 . (6.4)
Thus the operator −Aα/2 can be interpreted as the required fractional
derivative, i.e.
Aα/2 ≡ − d
d|x|α
, α > 0 . (6.5)
We note that the operator just defined must not be confused with a power
of the first order differential operator d
for which the symbol is −iκ .
After the above considerations it is straightforward to obtain the Fourier
image of the Green function of the Cauchy problem for the space-fractional
diffusion equation. In fact, applying the Fourier transform to the equation
(6.1), subject to the initial condition u(x, 0+;α) = δ(x) , and accounting for
(6.2), we obtain
Gc(x, t;α) = Gc(|x|, t;α) ÷ Ĝc(k, t;α) = e−D t |κ|
, 0 < α ≤ 2 . (6.6)
We easily recognize that the Fourier transform of the Green function
corresponds to the canonic form of a symmetric stable distribution of index
of stability α and scaling factor γ = (Dt)1/α , see (B.8). Therefore we have
Gc(x, t;α) = pα(x; 0, γ, 0) , γ = (Dt)1/α . (6.7)
For α = 1 and α = 2 we easily obtain the explicit expressions of the
corresponding Green functions since in these cases they correspond to the
Cauchy and Gauss distributions,
Gc(x, t; 1) =
x2 + (D t)2
, (6.8)
see (B.5), and
Gc(x, t; 2)) =
2/(4D t) , (6.9)
in agreement with (2.6).
We easily recognize that
(D t)1/α
(6.10)
is the similarity variable for the space-fractional diffusion equation, in terms
of which we can express the Green function for any α ∈ (0, 2] . Indeed, we
recognize that
Gc(x, t;α) =
(D t)1/α
qα(η; 0) , (6.11)
where qα(η; 0) denotes the symmetric stable distribution of order α with
Feller-type characteristic function, see (B.14-15). Now we can express the
Green function using the Feller series expansions (B.21-22) with θ = 0 . We
obtain:
for 0 < α < 1 ,
qα(η; 0) = −
Γ(nα + 1)
, (6.12a)
for 1 < α ≤ 2 ,
qα(η; 0) =
(−1)m
Γ[(2m + 1)/α]
(2m)!
η2m . (6.12b)
In the limiting case α = 1 the above series reduce to geometrical series and
therefore are no longer convergent in all of C . In particular, they represent
the expansions of the function q1(η; 0) = 1/[π(1+η
2)] , convergent for η > 1
and 0 < η < 1 , respectively.
We also note that for any α ∈ (0, 2] the functions qα(η; 0) exhibit at the
origin the value qα(0; 0) = Γ(1/α)/(π α) , and at the queues, excluding the
Gaussian case α = 2 , the algebraic asymptotic behaviour, as η → ∞ ,
qα(η; 0) ∼
Γ(α + 1) sin
η−(α+1) , 0 < α < 2 . (6.13)
In Figure 3, as an example, we compare versus x , at fixed t , the fundamental
solutions of the Cauchy problem with different α (α = 1/2 , 1 , 3/2 , 2 ). We
consider the range −6 ≤ x ≤ +6 and assume D = t = 1 .
-6 -4 -2 0 2 4 6
-6 -4 -2 0 2 4 6
Figure 3: The Cauchy problem for the simmetric space-fractional diffusion
equation. The fundamental solutions versus x : plate a) α = 1/2
(continuous line), α = 1 (dashed line); plate b) α = 3/4 (continuous line),
α = 2 (dashed line).
Let us now express more properly our operator (6.4) (with symbol |κ|α)
as inverse of a suitable integral operator Iα whose symbol is |κ|−α . This
operator can be found in the approach by Marcel Riesz to Fractional
Calculus, see e.g. Samko, Kilbas & Marichev (1987-1993) and Rubin (1996).
We recall that for any α > 0 , α 6= 1 , 3 , 5 , . . . and for a sufficiently well-
behaved function φ(x) , x ∈ R , the Riesz integral or Riesz potential Iα and
its image in the Fourier domain read
Iα φ(x) :=
2Γ(α) cos(πα/2)
|x − ξ|α−1 φ(ξ) dξ ÷ φ̂(κ)
. (6.14)
On its turn, the Riesz potential can be written in terms of two Weyl integrals
Iα± according to
Iα φ(x) =
2 cos(πα/2)
Iα+φ(x) + I
−φ(x)
, (6.15)
where
Iα+ φ(x) :=
(x − ξ)α−1 φ(ξ) dξ ,
Iα− φ(x) :=
(ξ − x)α−1 φ(ξ) dξ .
(6.16)
Then, at least in a formal way, the space-fractional derivative (6.2) turns
out to be defined as the opposite of the (left) inverse of the Riesz fractional
integral, i.e.
d|x|α
φ(x) := −I−α φ(x) = −
2 cos(πα/2)
I−α+ φ(x) + I
− φ(x)
. (6.17)
Notice that (6.14) and (6.17) become meaningless when α is an integer odd
number. However, for our range of interest 0 < α ≤ 2 , the particular case
α = 1 can be singled out since the corresponding Green function is already
known, see (6.8). Thus, excluding the case α = 1 , our space-fractional
diffusion equation (6.1) can be re-written, x ∈ R , t ∈ R+0 , as
= −D I−α u , u = u(x, t;α) , 0 < α ≤ 2 , α 6= 1 , (6.18)
where the operator I−α is defined by (6.16-17).
Here, in order to evaluate the fundamental solution of the Cauchy problem,
interpreted as a probability density, we propose a numerical approach,
original as far as we know, based on a (symmetric) random walk model,
discrete in space and time, see also Gorenflo & Mainardi (1998a), Gorenflo
& Mainardi (1998b) and Gorenflo, De Fabritiis & Mainardi (1999). We shall
see how things become highly transparent, in that we properly generalize
the classical random-walk argument of the standard diffusion equation
to our space-fractional diffusion equation (6.18). So doing we are in
position to provide a numerical simulation of the related (symmetric) stable
distributions in a way analogous to the standard one for the Gaussian law.
The essential idea is to approximate the left inverse operators I−α± by the
Grünwald-Letnikov scheme, on which the reader can inform himself in the
treatises on fractional calculus, see e.g. Oldham & Spanier (1974), Samko,
Kilbas & Marichev (1987-1993), Miller & Ross (1993), or in the recent review
article by Gorenflo (1997). If h denotes a ”small” positive step-length, these
approximating operators read
± φ(x) :=
(−1)k
φ(x ∓ kh) . (6.19)
Assume, for simplicity, D = 1 , and introduce grid points xj = j h with
h > 0 , j ∈ Z , and time instances tn = n τ with τ > 0 , n ∈ N0 . Let there
be given probabilities pj,k ≥ 0 of jumping from point xj at instant tn to
point xk at instant tn+1 and define probabilities yj(tn) of the walker being
at point xj at instant tn. Then, by
yk(tn+1) =
pj,k uj(tn) ,
pj,k =
pj,k = 1 , (6.20)
with pj,k = pk,j , a symmetric random walk (more precisely a symmetric
random jump) model is described. With the approximation
yj(tn) ≈
∫ (xj+h/2)
(xj−h/2)
u(x, tn) dx ≈ hu(xj , tn) , (6.21)
and introducing the ”scaling parameter”
2 | cos(απ/2)|
, (6.22)
we have solved
yj(tn+1) − yj(tn)
= − hI−α yj(tn) , (6.23)
for yj(tn+1) . So we have proved to have a consistent (for h → 0) symmetric
random walk approximation to (6.18) by taking
i) for 0 < α < 1 , 0 < µ ≤ 1/2 ,
−α yj(tn) = µ
+ yj(tn) + hI
− yj(tn)
pj,j = 1 − 2µ , pj,j±k = µ
)∣∣ , k ≥ 1 ;
(6.24)
ii) for 1 < α ≤ 2 , 0 < µ ≤ 1/(2α) ,
−α yj(tn) = µ
+ yj+1(tn) + hI
− yj−1(tn)
pj,j = 1 − 2µ α , pj,j±1 = µ
pj,j±k = µ
)∣∣∣ , k ≥ 2 .
(6.25)
We note that our random walk model is not only symmetric, but also
homogeneous, the transition probabilities pj,j±k not depending on the index
In the special case α = 2 we recover from (6.25) the well-known three-point
approximation of the heat equation, because pj,j±k = 0 for k ≥ 2 . This
means that for approximation of common diffusion only jumps of one step
to the right or one to the left or jumps of width zero occur, whereas for
0 < α < 2 (α 6= 1) arbitrary large jumps occur with power-like decaying
probability, as it turns out from the asymptotic analysis for the transition
probabilities given in (6.24-25). In fact, as k → ∞ , one finds
pj,j+k ∼
(τ/hα)
Γ(α + 1) sin
k−(α+1) , 0 < α < 2 . (6.26)
This result thus provides the discrete counterpart of the asymptotic
behaviour of the long power-law tails of the symmetric stable distributions,
as foreseen by (6.13) when 0 < α < 2 .
7 Conclusions
We have treated two generalizations of the standard, one-dimensional,
diffusion equation, namely, the time-fractional diffusion equation and the
symmetric space-fractional diffusion equation. For these equations we have
derived the fundamental solutions using the transform methods of Fourier
and Laplace, and exhibited their connections to extremal and symmetric
stable probability densities, evolving on time or variable in space. For the
symmetric space-fractional diffusion equation we have presented a stationary
(in time), homogeneous (in space) symmetric random walk model, discrete
in space and time, the step-lengths of the spatial grid and the time lapses
between transitions properly scaled. In the limit of infinitesimally fine
discretization this model (based on the Grünwald-Letnikov approximation
to fractional derivatives) is consistent with the continuous diffusion process,
i.e. convergent if interpreted as a difference scheme in the sense of numerical
analysis2.
From the mathematical viewpoint the field of such ”fractional” general-
izations is fascinating as there several mathematical disciplines meet and
come to a fruitful interplay: e.g. probability theory and stochastic processes,
2Further generalizations have been considered by us and our collaborators in other
papers, in which we have given a derivation of discrete random walk models related to
more general space-time fractional diffusion equations. For a comprehensive analysis, see
Gorenflo et al. (2002). Readers interested to the fundamental solutions of these fractional
diffusion equations are referred to the paper by Mainardi et al. (2001) where analytical
expressions and numerical plots are found.
integro-differential equations, transform theory, special functions, numerical
analysis. As one may take from our References, one can observe that since
some decades there is an ever growing interest in using the concepts of
fractional calculus among physicists and economists. Among economists we
like to refer the reader to a collection of papers on the topic of ”Fractional
Differencing and Long Memory Processes”, edited by Baillie & King (1996).
Appendix A: The Riemann-Liouville Fractional
Calculus
Fractional calculus is the field of mathematical analysis which deals with the
investigation and applications of integrals and derivatives of arbitrary order.
The term fractional is a misnomer, but it is retained following the prevailing
use. This appendix is mostly based on the recent review by Gorenflo &
Mainardi (1997). For more details on the classical treatment of fractional
calculus the reader is referred to Erdélyi (1954), Oldham & Spanier (1974),
Samko et al. (1987-1993) and Miller & Ross (1993).
According to the Riemann-Liouville approach to fractional calculus, the
notion of fractional Integral of order α (α > 0) is a natural consequence
of the well known formula (usually attributed to Cauchy), that reduces the
calculation of the n−fold primitive of a function f(t) to a single integral of
convolution type. In our notation the Cauchy formula reads
Jnf(t) := fn(t) =
(n − 1)!
(t − τ)n−1 f(τ) dτ , t > 0 , n ∈ N , (A.1)
where N is the set of positive integers. From this definition we note that
fn(t) vanishes at t = 0 with its derivatives of order 1, 2, . . . , n − 1 . For
convention we require that f(t) and henceforth fn(t) be a causal function,
i.e. identically vanishing for t < 0. In a natural way one is led to extend
the above formula from positive integer values of the index to any positive
real values by using the Gamma function. Indeed, noting that (n − 1)! =
Γ(n) , and introducing the arbitrary positive real number α , one defines the
Fractional Integral of order α > 0 :
Jα f(t) :=
(t − τ)α−1 f(τ) dτ , t > 0 , α ∈ R+ , (A.2)
where R+ is the set of positive real numbers. For complementation we define
J0 := I (Identity operator), i.e. we mean J0 f(t) = f(t) . Furthermore, by
Jαf(0+) we mean the limit (if it exists) of Jαf(t) for t → 0+ ; this limit
may be infinite.
We note the semigroup property JαJβ = Jα+β , α , β ≥ 0 , which implies
the commutative property JβJα = JαJβ , and the effect of our operators Jα
on the power functions
Jαtγ =
Γ(γ + 1)
Γ(γ + 1 + α)
tγ+α , α ≥ 0 , γ > −1 , t > 0 . (A.3)
These properties are of course a natural generalization of those known when
the order is a positive integer.
Introducing the Laplace transform by the notation L {f(t)} :=∫∞
−st f(t) dt = f̃(s) , s ∈ C , and using the sign ÷ to denote a Laplace
transform pair, i.e. f(t) ÷ f̃(s) , we note the following rule for the Laplace
transform of the fractional integral,
Jα f(t) ÷ f̃(s)
, α ≥ 0 , (A.4)
which is the generalization of the case with an n-fold repeated integral.
After the notion of fractional integral, that of fractional derivative of order
α (α > 0) becomes a natural requirement and one is attempted to substitute
α with −α in the above formulas. However, this generalization needs some
care in order to guarantee the convergence of the integrals and preserve the
well known properties of the ordinary derivative of integer order.
Denoting by Dn with n ∈ N , the operator of the derivative of order n ,
we first note that Dn Jn = I , Jn Dn 6= I , n ∈ N , i.e. Dn is left-inverse
(and not right-inverse) to the corresponding integral operator Jn . In fact
we easily recognize from (A.1) that
Jn Dn f(t) = f(t) −
f (k)(0+)
, t > 0 . (A.5)
As a consequence we expect that Dα is defined as left-inverse to Jα. For
this purpose, introducing the positive integer m such that m − 1 < α ≤ m ,
one defines the Fractional Derivative of order α > 0 :
Dα f(t) := Dm Jm−α f(t) , m − 1 < α ≤ m , m ∈ N , (A.6)
namely
Dα f(t)=
Γ(m − α)
(t − τ)α+1−m
, m − 1 < α < m,
f(t) , α = m.
(A.6′)
Defining for complementation D0 = J0 = I , then we easily recognize that
Dα Jα = I , α ≥ 0 , and
Dα tγ =
Γ(γ + 1)
Γ(γ + 1 − α)
tγ−α , α ≥ 0 , γ > −1 , t > 0 . (A.7)
Of course, these properties are a natural generalization of those known when
the order is a positive integer.
Note the remarkable fact that the fractional derivative Dα f is not zero for
the constant function f(t) ≡ 1 if α 6∈ N . In fact, (A.7) with γ = 0 teaches
us that
Dα1 =
Γ(1 − α)
, α ≥ 0 , t > 0 . (A.8)
This, of course, is ≡ 0 for α ∈ N, due to the poles of the gamma function in
the points 0,−1,−2, . . .. We now observe that an alternative definition of
fractional derivative, originally introduced by Caputo (1967) (1969) in the
late sixties and adopted by Caputo and Mainardi (1971) in the framework
of the theory of Linear Viscoelasticity, is
Dα∗ f(t) := J
m−α Dm f(t) m − 1 < α ≤ m , m ∈ N , (A.9)
namely
D ∗α f(t) =
Γ(m − α)
f (m)(τ)
(t − τ)α+1−m
dτ , m − 1 < α < m,
f(t) , α = m.
(A.9′)
This definition is of course more restrictive than (A.6), in that requires
the absolute integrability of the derivative of order m. Whenever we use
the operator Dα∗ we (tacitly) assume that this condition is met. We easily
recognize that in general
Dα f(t) := Dm Jm−α f(t) 6= Jm−α Dm f(t) := Dα∗ f(t) , (A.10)
unless the function f(t) along with its first m − 1 derivatives vanishes at
t = 0+. In fact, assuming that the passage of the m-derivative under the
integral is legitimate, one recognizes that, for m − 1 < α < m and t > 0 ,
Dα f(t) = Dα∗ f(t) +
Γ(k − α + 1)
f (k)(0+) , (A.11)
and therefore, recalling the fractional derivative of the power functions (A.7),
f(t) −
f (k)(0+)
= Dα∗ f(t) . (A.12)
The alternative definition (A.9) for the fractional derivative thus incorpo-
rates the initial values of the function and of its integer derivatives of lower
order. The subtraction of the Taylor polynomial of degree m − 1 at t = 0+
from f(t) means a sort of regularization of the fractional derivative. In
particular, according to this definition, the relevant property for which the
fractional derivative of a constant is still zero can be easily recognized, i.e.
Dα∗ 1 ≡ 0 , α > 0 . (A.13)
We now explore the most relevant differences between the two fractional
derivatives (A.6) and (A.9). We agree to denote (A.9) as the Caputo
fractional derivative to distinguish it from the standard Riemann-Liouville
fractional derivative (A.6). We observe, again by looking at (A.7), that
Dαtα−1 ≡ 0 , α > 0 , t > 0 .
From above we thus recognize the following statements about functions
which for t > 0 admit the same fractional derivative of order α , with
m − 1 < α ≤ m , m ∈ N ,
Dα f(t) = Dα g(t) ⇐⇒ f(t) = g(t) +
α−j , (A.14)
Dα∗ f(t) = D
∗ g(t) ⇐⇒ f(t) = g(t) +
m−j . (A.15)
In these formulas the coefficients cj are arbitrary constants.
For the two definitions we also note a difference with respect to the formal
limit as α → (m − 1)+ ; from (A.6) and (A.9) we obtain respectively,
Dα f(t) → Dm J f(t) = Dm−1 f(t) ; (A.16)
Dα∗ f(t) → J Dm f(t) = Dm−1 f(t) − f (m−1)(0+) . (A.17)
We now consider the Laplace transform of the two fractional derivatives.
For the standard fractional derivative Dα the Laplace transform, assumed to
exist, requires the knowledge of the (bounded) initial values of the fractional
integral Jm−α and of its integer derivatives of order k = 1, 2, . . . ,m−1 . The
corresponding rule reads, in our notation,
Dα f(t) ÷ sα f̃(s) −
Dk J (m−α) f(0+) sm−1−k , (A.18)
where m − 1 < α ≤ m .
The Caputo fractional derivative appears more suitable to be treated by
the Laplace transform technique in that it requires the knowledge of the
(bounded) initial values of the function and of its integer derivatives of
order k = 1, 2, . . . ,m− 1 , in analogy with the case when α = m . In fact, by
using (A.4) and noting that
Jα Dα∗ f(t) = J
α Jm−α Dm f(t) = Jm Dm f(t) = f(t) −
f (k)(0+)
(A.19)
we easily prove the following rule for the Laplace transform,
Dα∗ f(t) ÷ sα f̃(s) −
f (k)(0+) sα−1−k , m − 1 < α ≤ m . (A.20)
Indeed, the result (A.20), first stated by Caputo (1969) by using the
Fubini-Tonelli theorem, appears as the most ”natural” generalization of the
corresponding result well known for α = m .
Gorenflo and Mainardi (1997) have pointed out the major utility of the
Caputo fractional derivative in the treatment of differential equations of
fractional order for physical applications. In fact, in physical problems,
the initial conditions are usually expressed in terms of a given number
of bounded values assumed by the field variable and its derivatives of
integer order, no matter if the governing evolution equation may be a
generic integro-differential equation and therefore, in particular, a fractional
differential equation3.
We now analyze the most simple differential equations of fractional order,
including those which, by means of fractional derivatives, generalize the well-
known ordinary differential equations related to relaxation and oscillation
3We note that the Caputo fractional derivative was so named after the book by
Podlubny (1999). It coincides with that introduced, independently and a few later,
by Dzherbashyan and Nersesyan (1968) as a regularization of the Riemann-Liouville
fractional derivative. Nowadays, some Authors refer to it as the Caputo-Dzherbashyan
fractional derivative. The prominent role of this fractional derivative in treating initial
value problems was recognized in interesting papers by Kochubei (1989), (1990).
phenomena. Generally speaking, we consider the following differential
equation of fractional order α > 0 ,
Dα∗ u(t) = D
u(t) −
u(k)(0+)
= −u(t) + q(t) , t > 0 , (A.21)
where u = u(t) is the field variable and q(t) is a given function. Here m is
a positive integer uniquely defined by m − 1 < α ≤ m , which provides the
number of the prescribed initial values u(k)(0+) = ck , k = 0, 1, 2, . . . ,m−1 .
Implicit in the form of (A.21) is our desire to obtain solutions u(t) for which
the u(k)(t) are continuous. In particular, the cases of fractional relaxation
and fractional oscillation are obtained for 0 < α < 1 and 1 < α < 2 ,
respectively
The application of the Laplace transform through the Caputo formula (A.20)
yields
ũ(s) =
sα−k−1
sα + 1
sα + 1
q̃(s) . (A.22)
Now, in order to obtain the Laplace inversion of (A.22), we need to recall
the Mittag-Leffler function of order α > 0 , Eα(z) . This function, so named
from the great Swedish mathematician who introduced it at the beginning
of this century, is defined by the following series and integral representation,
valid in the whole complex plane,
Eα(z) =
Γ(αn + 1)
σα−1 e σ
σα − z
dσ , α > 0 . (A.23)
Here Ha denotes the Hankel path, i.e. a loop which starts and ends at −∞
and encircles the circular disk |σ| ≤ |z|1/α in the positive sense. It turns out
that Eα(z) is an entire function of order ρ = 1/α and type 1 .
The Mittag-Leffler function provides a simple generalization of the expo-
nential function, to which it reduces for α = 1 . Particular cases from which
elementary functions are recovered, are
= cosh z , E2
= cos z , z ∈ C , (A.24)
E1/2(±z1/2) = ez
1 + erf (±z1/2)
= ez erfc (∓z1/2) , z ∈ C , (A.25)
where erf (erfc) denotes the (complementary) error function. defined as
erf (z) :=
du , erfc (z) := 1 − erf (z) , z ∈ C .
A noteworthy property of the Mittag-Leffler function is based on the
following duplication formula
Eα(z) =
Eα/2(+z
1/2) + Eα/2(−z1/2)
. (A.26)
In (A.25-26) we agree to denote by z1/2 the main branch of the complex
root of z .
The Mittag-Leffler function is connected to the Laplace integral through the
equation ∫ ∞
e−u Eα (u
α z) du =
1 − z
α > 0 . (A.27)
The integral at the L.H.S. was evaluated by Mittag-Leffler who showed that
the region of its convergence contains the unit circle and is bounded by the
line Re z1/α = 1 . The above integral is fundamental in the evaluation of the
Laplace transform of Eα (−λ tα) with α > 0 and λ ∈ C . In fact, putting in
(A.27) u = st and uα z = −λ tα with t ≥ 0 and λ ∈ C , we get the Laplace
transform pair
Eα (−λ tα) ÷
sα + λ
, Re s > |λ|1/α . (A.28)
Then, using (A.28), we put for k = 0, 1, . . . ,m − 1 ,
uk(t) := J
keα(t) ÷
sα−k−1
sα + 1
, eα(t) := Eα(−tα) , (A.29)
and, from inversion of the Laplace transforms in (A.22), we find
u(t) =
ck uk(t) −
q(t − τ)u′0(τ) dτ . (A.30)
In particular, the formula (A.30) encompasses the solutions for α = 1 , 2 ,
since e1(t) = exp(−t) , e2(t) = cos t . When α is not integer, namely for
m − 1 < α < m , we note that m − 1 represents the integer part of α
(usually denoted by [α]) and m the number of initial conditions necessary
and sufficient to ensure the uniqueness of the solution u(t). Thus the m
functions uk(t) = J
keα(t) with k = 0, 1, . . . ,m−1 represent those particular
solutions of the homogeneous equation which satisfy the initial conditions
+) = δk h , h, k = 0, 1, . . . ,m − 1 , and therefore they represent the
fundamental solutions of the fractional equation (A.21), in analogy with the
case α = m . Furthermore, the function uδ(t) = −u′0(t) = −e′α(t) represents
the impulse-response solution.
The Mittag-Leffler function of order less than one turns out to be related
through the Laplace integral to another special function of Wright type,
denoted by M(z, ν) with 0 < ν < 1 , following the notation introduced
by Mainardi (1994, 1995). Since this function turns out to be relevant in
the general framework of fractional calculus with special regard to stable
probability distributions, we are going to summarize its basing properties.
For more details on this function, see Mainardi (1997), Appendix A.
Let us first recall the more general Wright function Wλ,µ(z) , z ∈ C , with
λ > −1 and µ > 0 . This function, so named from the British mathematician
who introduced it between 1933 and 1941, is defined by the following series
and integral representation, valid in the whole complex plane,
Wλ,µ(z) =
n! Γ(λn + µ)
eσ + zσ
−λ dσ
, (A.31)
where Ha denotes the Hankel path. It is possible to prove that the Wright
function is entire of order 1/(1+λ) , hence of exponential type if λ ≥ 0 . The
case λ = 0 is trivial since W0,µ(z) = e
z/Γ(µ) . The case λ = −ν , µ = 1 − ν
with 0 < ν < 1 provides the function M(z, ν) of special interest for us.
Specifically, we have
M(z; ν) := W−ν,1−ν(−z) =
W−ν,0(−z) , 0 < ν < 1 , (A.32)
and therefore from (A.31-32)
M(z; ν) =
(−z)n−1
(n − 1)!
Γ(ν n) sin (ν n π)
eσ − zσ
, 0 < ν < 1 .
(A.33)
In the series representation we have used the reflection formula for the
Gamma function, Γ(x) Γ(1−x) = π/ sin πx . Explicit expressions of M(z; ν)
in terms of simpler known functions are expected in particular cases when
ν is a rational number. Relevant cases are ν = 1/2 , 1/3 for which
M(z; 1/2) =
− z2/4
, (A.34)
M(z; 1/3) = 32/3 Ai
z/31/3
, (A.35)
where Ai denotes the Airy function.
When the argument is real and positive, i.e. z = r > 0 , the existence of
the Laplace transform of M(r; ν) is ensured by the asymptotic behaviour,
as derived by Mainardi & Tomirotti (1995), as r → +∞ ,
M(r/ν; ν) ∼ a(ν) r(ν − 1/2)/(1 − ν) exp
−b(ν) r1/(1 − ν)
, (A.36)
where a(ν) = 1/
2π (1 − ν) , b(ν) = (1 − ν)/ν .
It is an instructive exercise to derive the Laplace transform by interchanging
the Laplace integral with the Hankel integral in (A.33) and recalling the
integral representation (A.23) of the Mittag-Leffler function. We obtain the
Laplace transform pair
M(r; ν) ÷ Eν(−s) , 0 < ν < 1 . (A.37)
For ν = 1/2 , (A.37) with (A.25) and (A.34) provides the result, see e.g.
Doetsch (1974),
M(r; 1/2) :=
− r2/4
÷ E1/2(−s) := exp
erfc (s) . (A.38)
It would be noted that, since M(r, ν) is not of exponential order,
transforming term-by-term the Taylor series of M(r; ν) yields a series of
negative powers of s , which represents the asymptotic expansion of Eν(−s)
as s → ∞ in a certain sector around the real axis.
We also note that (A.37) with (A.23) allows us to compute the moments of
any real order δ ≥ 0 of M(r; ν) in the positive real axis. We obtain
r δ M(r; ν) dr =
Γ(δ + 1)
Γ(νδ + 1)
, δ ≥ 0 . (A.39)
When δ is integer we note that the moments are provided by the derivatives
of the Mittag-Leffler function in the origin, i.e.
rn M(r; ν) dr = lim
(−1)n
Eν(−s) =
Γ(n + 1)
Γ(νn + 1)
, (A.40)
where n = 0, 1, 2, . . . . The normalization condition
0 M(r; ν) dr =
Eν(0) = 1 is recovered for n = 0 . The relation with the Mittag-Leffler
function stated in (A.40) can be extended to the moments of non integer
order if we replace the ordinary derivative, of order n, with the corresponding
fractional derivative, of order δ 6= n, in the Caputo sense.
Another exercise on the function M concerns the inversion of the Laplace
transform exp(−sν) , either by the complex integral formula or by the formal
series method. We obtain the Laplace transform pair
M (1/rν ; ν) ÷ exp (−sν) , 0 < ν < 1 . (A.41)
For ν = 1/2 , (A.41) with (A.34) provides the known result, see e.g. Doetsch
(1974),
2 r3/2
M(1/r1/2; 1/2) :=
π r3/2
exp [− 1/(4r)] ÷ exp
−s1/2
. (A.42)
We recall that a rigorous proof of (A.41) was formerly given by Pollard
(1946), based on a formal result by Humbert (1945). The Laplace transform
pair was also obtained by Mikusiński (1959) and, albeit unaware of the
previous results, by Buchen & Mainardi (1975) in a formal way.
Appendix B: The Stable Probability Distributions
The stable distributions are a fascinating and fruitful area of research in
probability theory; furthermore, nowadays, they provide valuable models in
physics, astronomy, economics, and communication theory.
The general class of stable distributions was introduced and given this name
by the French mathematician Paul Lévy in the early 1920’s, see Lévy (1924,
1925). The inspiration for Lévy was the desire to generalize the celebrated
Central Limit Theorem, according to which any probability distribution
with finite variance belongs to the domain of attraction of the Gaussian
distribution.
Formerly, the topic attracted only moderate attention from the leading
experts, though there were also enthusiasts, of whom the Russian
mathematician Alexander Yakovlevich Khintchine should be mentioned first
of all. The concept of stable distributions took full shape in 1937 with the
appearance of Lévy’s monograph, see Lévy (1937-1954), soon followed by
Khintchine’s monograph, see Khintchine (1938).
The theory and properties of stable distributions are discussed in some
classical books on probability theory including Gnedenko & Kolmogorov
(1949-1954), Lukacs (1960-1970), Feller (1966-1971), Breiman (1968-1992),
Chung (1968-1974) and Laha & Rohatgi (1979). Also treatises on fractals
devote particular attention to stable distributions in view of their properties
of scale invariance, see e.g. Mandelbrot (1982) and Takayasu (1990). Sets of
tables and graphs have been provided by Mandelbrot & Zarnfaller (1959),
Fama & Roll (1968), Bo’lshev & Al. (1968) and Holt & Crow (1973).
Only recently, monographs devoted solely to stable distributions and related
stochastic processes have been appeared, i.e. Zolotarev (1983-1986), Janicki
& Weron (1994), Samorodnitsky & Taqqu (1994), Uchaikin & Zolotarev
(1999). We now can cite the paper by Mainardi, Luchko & Pagnini (2001)
where the reader can find (convergent and asymptotic) representations and
plots of the symmetric and non-symmetric stable densities generated by
fractional diffusion equations.
Stable distributions have three exclusive properties, which can be briefly
summarized stating that they 1) are invariant under addition, 2) possess
their own domain of attraction, and 3) admit a canonic characteristic
function.
Let us now illustrate the above properties which, providing necessary and
sufficient conditions, can be assumed as equivalent definitions for a stable
distribution. We recall the basic results without proof.
A random variable X is said to have a stable distribution P (x) = Prob {X ≤
x} if for any n ≥ 2 , there is a positive number cn and a real number dn such
X1 + X2 + . . . + Xn
= cn X + dn , (B.1)
where X1,X2, . . . Xn denote mutually independent random variables with
common distribution P (x) with X . Here the notation
= denotes equality
in distribution, i.e. means that the random variables on both sides have the
same probability distribution.
When mutually independent random variables have a common distribution
[shared with a given random variable X], we also refer to them as
independent, identically distributed (i.i.d) random variables [independent
copies of X]. In general, the sum of i.i.d. random variables becomes
a random variable with a distribution of different form. However, for
independent random variables with a common stable distribution, the sum
obeys to a distribution of the same type, which differs from the original
one only for a scaling (cn) and possibly for a shift (dn). When in (B.1) the
dn = 0 the distribution is called strictly stable.
It is known, see Feller (1966-1971), that the norming constants in (B.1) are
of the form
cn = n
1/α with 0 < α ≤ 2 . (B.2)
The parameter α is called the characteristic exponent or the index of stability
of the stable distribution.
We agree to use the notation X ∼ Pα(x) to denote that the random variable
X has a stable probability distribution with characteristic exponent α . We
simply refer to P (x) , p(x) := dP/dx (probability density function = pdf)
and X as α-stable distribution, density, random variable, respectively.
The definition (B.1) with the theorem (B.2) can be stated in an alternative
version that needs only two i.i.d. random variables. see also Lukacs (1960-
1970). A random variable X is said to have a stable distribution if for any
positive numbers A and B, there is a positive number C and a real number
D such that
AX1 + B X2
= C X + D , (B.3)
where X1 and X2 are independent copies of X . Then there is a number
α ∈ (0, 2] such that the number C in (B.3) satisfies Cα = Aα + Bα .
For a strictly stable distribution (B.3) holds with D = 0 . This implies that
all linear combinations of i.i.d. random variables obeying to a strictly stable
distribution is a random variable with the same type of distribution.
A stable distribution is called symmetric if the random variable −X has the
same distribution. Of course, a symmetric stable distribution is necessarily
strictly stable.
Noteworthy examples of stable distributions are provided by the Gaussian
(or normal) law (with α = 2) and by the Cauchy-Lorentz law (α = 1). The
corresponding pdf are known to be
pG(x;σ, µ) :=
e−(x − µ)
2/(2σ2) , x ∈ R , (B.4)
where σ2 denotes the variance and µ the mean, and
pC(x; γ, δ) :=
(x − δ)2 + γ2
, x ∈ R , (B.5)
where γ denotes the semi-interquartile range and δ the ”shift”.
Another (equivalent) definition states that stable distributions are the only
distributions that can be obtained as limits of normalized sums of i.i.d.
random variables. A random variable X is said to have a domain of
attraction,i.e. if there is a sequence of i.i.d. random variables Y1, Y2, . . .
and sequences of positive numbers {γn} and real numbers {δn}, such that
Y1 + Y2 + . . . Yn
d⇒X . (B.6)
The notation
d⇒ denotes convergence in distribution.
It is clear that the previous definition (B.1) yields (B.6), e.g. , by taking the
Yis to be independent and distributed like X . The converse is easy to show,
see Gnedenko & Kolmogorov (1949-1954). Therefore we can alternatively
state that a random variable X is said to have a stable distribution if it has
a domain of attraction.
When X is Gaussian and the Yis are i.i.d. with finite variance, then (B.6)
is the statement of the ordinary Central Limit Theorem. The domain
of attraction of X is said normal when γn = n
1/α ; in general, γn =
n1/α h(n) where h(x) , x > 0 , is a slow varying function at infinity, that
is, lim
h(ux)/h(x) = 1 for all u > 0 , see Feller (1971). The function
h(x) = log x , for example, is slowly varying at infinity.
Another definition specifies the canonic form that the characteristic function
(cf) of a stable distribution of index α must have. Recalling that the cf is
the Fourier transform of the pdf , we use the notation p̂α(κ) := 〈exp (iκX)〉 ÷
pα(x) . We first note that a stable distribution is also infinitely divisible, i.e.
for every positive integer n its cf can be expressed as the nth power of
some cf . In fact, using the characteristic function, the relation (B.1) is
transformed into
[p̂α(κ)]
n = p̂α(cn κ) e
idnκ . (B.7)
The functional equation (B.7) can be solved completely and the solution is
known to be
p̂α(κ;β, γ, δ) = exp {iδκ − γα |κ|α [1 + i (sign κ)β ω(|κ|, α)]} , (B.8)
where
ω(|κ|, α) =
tan (α π/2) , if α 6= 1 ,
−(2/π) log |κ| , if α = 1 . (B.9)
Consequently a random variable X is said to have a stable distribution if
there are four real parameters α, β, γ, δ with 0 < α ≤ 2 , −1 ≤ β ≤ +1 ,
γ > 0 , such that its characteristic function has the canonic form (B.8-9).
Then we write pα(x;β, γ, δ)÷ p̂α(κ;β, γ, δ) and X ∼ Pα(x;β, γ, δ) , so partly
following the notation of Holt & Crow (1973) and Samorodnitsky & Taqqu
(1994).
We note in (B.8-9) that β appears with different signs for α 6= 1 and α = 1 .
This minor point has been the source of great confusion in the literature, see
Hall (1980) for a discussion. The presence of the logarithm for α = 1 is the
source of many difficulties, so this case has often to be treated separately.
The cf (B.8-9) turns out to be a useful tool for studying α-stable distri-
butions and for providing an interpretation of the additional parameters,
β (skewness parameter), γ (scale parameter) and δ (shift parameter), see
Samorodnitsky & Taqqu (1994). When α = 2 the cf refers to the Gaussian
distribution with variance σ2 = 2 γ2 and mean µ = δ ; in this case the value
of the skewness parameter β is not specified because tan π = 0 , and one
conventionally takes β = 0 .
One easily recognizes that a stable distribution is symmetric if and only if
β = δ = 0 and is symmetric about δ if and only if β = 0 . Stable distributions
with extremal values of the skewness parameter are called extremal. One
can prove that all the extremal stable distributions with 0 < α < 1 are
one-sided, the support being R+0 if β = −1 , and R
0 if β = +1 .
For the stable distributions Pα(x;β, γ, δ) we now consider the asymptotic
behaviour of the tail probabilities, T+(λ) := Prob {X > λ} and T−(λ) :=
Prob {X < −λ} , as λ → ∞ . For the Gaussian case α = 2 the result is well
known, see e.g. Feller (1957),
α = 2 : T±(λ) ∼ 1
2/(4γ2)
, λ → ∞ . (B.10)
Because of the above exponential decay all the moments of the corresponding
pdf turn out to be finite, which is an exclusive property of this stable
distribution. For all the other stable distributions the singularity of the
characteristic function in the origin is responsible for the algebraic decay of
the tail probabilities as indicated below, see e.g. Samorodnitsky & Taqqu
(1994),
0 < α < 2 : lim
λα T±(λ) = Cα γ
α (1 ∓ β)/2 , (B.11)
where
x−α sin x dx
1 − α
Γ(2 − a) cos (απ/2)
, if α 6= 1 ,
2/π , if α = 1 .
(B.12)
We note that for extremal distributions (β = ±1) the above algebraic decay
holds true only for one tail, the left one if β = +1 , the right one if β = −1 .
The other tail is either identically zero if 0 < α < 1 (the distribution is
one-sided !), or exhibits an exponential decay if 1 ≤ α < 2 . Because of the
algebraic decay we recognize that
0 < α < 2 :
|x|>λ
pα(x;β, γ, δ) dx = O(λ
−α) , (B.13)
so the absolute moments of a stable non-Gaussian pdf turn out to be finite
if their order ν is 0 ≤ ν < α and infinite if ν ≥ α . We are now convinced
that the Gaussian distribution is the unique stable distribution with finite
variance. Furthermore, when α ≤ 1 , the first absolute moment 〈|X|〉 is
infinite as well, so we need to use the median to characterize the expected
value.
There is however a fundamental property shared by all the stable
distributions that we like to point out: for any α the stable pdf are unimodal
and indeed bell-shaped, i.e. their n-th derivative has exactly n zeros, see
Gawronski (1964).
We now come back to the cf of a stable distribution, in order to provide for
α 6= 1 and δ = 0 a simpler canonic form which allow us to derive convergent
and asymptotic power series for the corresponding pdf . We first note that
the two parameters γ and δ in (B.8), being related to a scale transformation
and a translation, are not so essential since they do not change the shape
of distributions. If we take γ = 1 and δ = 0 , we obtain the so-called
standardized form of the stable distribution and X ∼ Pα(x;β, 1, 0) is referred
to as the α-stable standardized random variable. Furthermore, we can choose
the scale parameter γ in such a way to get from (B.8-9) the simplified canonic
form used by Feller (1952, 1966-1971) and Takayasu (1990) for strictly stable
distributions (δ = 0) with α 6= 1 , which reads in an ad hoc notation,
q̂α(κ; θ) :=
eiκ y pα(y; θ) dy = exp
−|κ|α e±i θ π/2
, (B.14)
where the symbol ± takes the sign of κ . This canonic form, that we refer to
as the Feller canonic form, is derived from (B.8-9) if in addition to α 6= 1
and δ = 0 we require
γα = cos
, tan
= β tan
. (B.15)
Here θ is the skewness parameter instead of β and its domain is restricted
in the following region (depending on α)
|θ| ≤
α , if 0 < α < 1 ,
2 − α , if 1 < α < 2 . (B.16)
Thus, when we use the Feller canonic form for strictly stable distributions
with index α 6= 1 and skewness θ , we implicitly select the scale parameter
γ (0 < γ ≤ 1), which is related to α , β and θ by (B.15). Specifically, the
random variable Y ∼ Qα(y; θ) turns out to be related to the standardized
random variable X ∼ Pα(x;β, 1, 0) by the following relations
Y = X/γ , pα(x;β, 1, 0) = γ qα(y = γx; θ) , (B.17)
with
γ = [cos (θπ/2)]1/α ,
θ = (2/π) arctan [β tan (απ/2)] ,
tan (θπ/2)
tan (απ/2)
(B.18)
We recognize that qα(y, θ) = qα(−y,−θ) , so the symmetric stable
distributions are obtained if and only if θ = 0 . We note that for the
symmetric stable distributions we get the identity between the standardized
and the Lévy canonic forms, since in (B.18) β = θ = 0 implies γ = 1 .
A particular but noteworthy case is provided by p2(x; 0, 1, 0) = q2(y; 0) ,
corresponding to the Gaussian distribution with variance σ2 = 2 .
The extremal stable distributions, corresponding to β = ±1 , are now
obtained for θ = ±α if 0 < α < 1 , and for θ = ∓(2 − α) if 1 < α < 2 ; for
them the scaling parameter turns out to be γ = [cos (|α|π/2)]1/α . It may be
an instructive exercise to carry out the inversion of the Fourier transform
when α = 1/2 and θ = −1/2 . In this case we obtain the analytical expression
for the corresponding extremal stable pdf , known as the (one-sided) Lévy-
Smirnov density,
q1/2(y;−1/2) =
y−3/2 e−1/(4y) , y ≥ 0 . (B.19)
The standardized form for this distribution can be easily obtained from
(B.19) using (B.17-18) with α = 1/2 and θ = −1/2 . We get γ =
[cos (−π/4)]2 = 1/2 , β = −1 , so
p1/2(x;−1, 1, 0) =
q1/2(x/2;−1/2) =
x−3/2 e−1/(2x) , (B.20)
where x ≥ 0 , in agreement with Holt & Crow (1973) [§2.13, p. 147].
Feller (1952) has obtained from (B.14) the following representations by
convergent power series for the stable distributions valid for y > 0 , with
0 < α < 1 (negative powers),
qα(y; θ) =
(−y−α)n Γ(nα + 1)
(θ − α)
, (B.21)
1 < α ≤ 2 (positive powers),
qα(y; θ) =
(−y)n
Γ(n/α + 1)
(θ − α)
. (B.22)
The values for y < 0 can be obtained from (B.21-22) using the identity
qα(−y; θ) = qα(y;−θ) , y > 0 . As a consequence of the convergence in all of
C of the series in (B.21-22) we recognize that the restrictions of the functions
y qα(y; θ) on the two real semi-axis turn out to be equal to certain entire
functions of argument 1/|y|α for 0 < α < 1 and argument |y| for 1 < α ≤ 2 .
It has be shown, see e.g. Bergström (1952), Chao Chung-Jeh (1953), that the
two series in (B.21-22) provide also the asymptotic (divergent) expansions to
the stable pdf with the ranges of α interchanged from those of convergence.
From (B.21-22) a relation between stable pdf with index α and 1/α can be
derived as noted in Feller (1966-1971). Assuming 1/2 < α < 1 and y > 0 ,
we obtain
q1/α(y
−α; θ) = qα(y; θ
∗) , θ∗ = α(θ + 1) − 1 . (B.23)
A quick check shows that θ∗ falls within the prescribed range, |θ∗| ≤ α ,
provided that |θ| ≤ 2 − 1/α .
We now consider two particular cases of the Feller series (B.21-22), of
particular interest for us, which turn out to be related to the entire function
of Wright type, M(z; ν) with 0 < ν < 1 , reported in Appendix A. These
cases correspond to the following extremal distributions
Φ1(y) := qα(y;−α) , y > 0 , 0 < α < 1 , (B.24)
Φ2(y) := qα(y;α − 2) , y > 0 , 1 < α ≤ 2 , (B.25)
for which the Feller series (B.21-22) reduce to
Φ1(y) =
(−1)n−1 y−αn−1 Γ(nα + 1)
sin (nπα) , y > 0 , (B.26)
Φ2(y) =
(−1)n−1 yn−1 Γ(n/α + 1)
, y > 0 . (B.27)
In fact, recalling the series representation of the general Wright function,
Wλ,µ(z) with λ > −1 , µ > 0 , see (A.31), and the definition of the function
M(z; ν) with 0 < ν < 1 , see (A.32-33), we recognize that
Φ1(y) =
W−α,0(−y−α) =
M(y−α;α) , y > 0 , (B.28)
Φ2(y) =
W−1/α,0(−y) =
M(y; 1/α) , y > 0 . (B.29)
We would like to remark that the above relations with the Wright functions
have been noted also by Engler (1997).
It is worth to point out that, whereas Φ1(y) totally represents the one-
sided stable pdf qα(y;−α) , 0 < α < 1 , with support in R+0 , Φ2(y) is the
restriction on the positive axis of qα(y;α− 2) , 1 < α ≤ 2 , whose support is
all of R . Since the function M(z; ν) turns out to be normalized in R+0 , see
(A.39-40), we also note
Φ1(y) dy = 1 ;
Φ2(y) dy = 1/α . (B.30)
Using the results (A.41) and (A.37) we can easily evaluate the Laplace
transforms of Φ1(y) and Φ2(y) , respectively. We obtain
L[Φ1(y)] = Φ̃1(s) = exp (−sα) , 0 < α < 1 , (B.31)
L[Φ2(y)] = Φ̃2(s) =
E1/α (−s) , 1 < α ≤ 2 , (B.32)
where E1/α(·) denotes the Mittag-Leffler function of order 1/α , see (A.23).
It is an instructive exercise to derive the asymptotic behaviours of Φ1(y) and
Φ2(y) as y → 0+ and y → +∞ . By using the expressions (B.28−29) in terms
of the function M and recalling the series and asymptotic representations of
this function, see (A.33) and (A.36), we obtain
Φ1(y) =
y−(2−α)/[2(1−α)] e−c1 y
−α/(1−α)
, as y → 0+ ,
Γ(1 − α)
y−α−1 [1 + O (y−α)] , as y → +∞ ,
(B.33)
Φ2(y) =
Γ(1 − 1/α)
[1 + O (y)] , as y → 0+ ,
y(2−α)/[2(α−1)] e−c2 y
α/(α−1)
, as y → +∞ ,
(B.34)
where c1 , c2 are positive constants depending on α . We note that the
exponential decay is found for Φ1(y) as y → 0+ but as y → +∞ for Φ2(y) .
Explicit expressions for stable pdf can be derived form those for the function
M(z; ν) when ν = 1/2 and ν = 1/3 , given in Appendix A, see (A.34-
35). Of course the ν = 1/2 expression can be used to recover the well-
known (symmetric) Gaussian distribution q2(y; 0) accounting for (B.29), and
the (one-sided) Lévy distribution q1/2(y;−1/2), see (B.19), accounting for
(B.28). The ν = 1/3 expression provides, accounting for (B.28),
q1/3(y;−1/3) = 3−1/3 y−4/3 Ai
(3y)−1/3
y−3/2 K1/3
(B.35)
where Ai denotes the Airy function and K1/3 the modified Bessel function of
the second kind of order 1/3 . The equivalence between the two expressions
in (B.35) can be proved in view of the relation, see Abramowitz & Stegun
(1965-1972) [(10.4.14)],
Ai (z) =
. (B.36)
The case α = 1/3 has also been discussed by Zolotarev (1983-1986), who
has quoted the corresponding expression of the pdf in terms of K1/3 .
A general representation of all stable distributions (thus including the
extremal distributions above considered) in terms of special functions has
been only recently achieved by Schneider (1986). In his remarkable
(but almost ignored) article, Schneider has established that all the stable
distributions can be characterized in terms of a general class of special
functions, the so-called Fox H functions, so named after Charles Fox (1961).
For details on Fox H functions, see e.g. the books Mathai & Saxena (1978),
Srivastava & Al. (1982) and the most recent paper by Kilbas and Saigo
(1999). These functions are expressed in terms of special integrals in the
complex-plane, the Mellin-Barnes integrals4.
4The names refer to the two authors, who in the first 1910’s developed the theory of
these integrals using them for a complete integration of the hypergeometric differential
equation. However, as pointed out in the the Bateman Project Handbook on High
Transcendental Functions, see Erdelyi (1953), these integrals were first used by S. Pincherle
in 1888. For a revisited analysis of the pioneering work of Pincherle (1853-1936, Professor
of Mathematics at the University of Bologna from 1880 to 1928) we refer to the paper by
Mainardi and Pagnini (2003).
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Introduction
The standard diffusion equation
The time-fractional diffusion equation
The Cauchy problem for the time-fractional diffusion equation
The Signalling problem for the time-fractional diffusion equation
The Cauchy problem for the symmetric space-fractional diffusion equation
Conclusions
|
0704.0321 | Fabrication of half metallicity in a ferromagnetic metal | Fabrication of half metallicity in a ferromagnetic metal
Kalobaran Maiti∗
Department of Condensed Matter Physics and Materials’ Science,
Tata Institute of Fundamental Research, Homi Bhabha Road, Colaba, Mumbai - 400 005, INDIA
(Dated: August 15, 2021)
We investigate the growth of half metallic phase in a ferromagnetic material using state-of-the-art
full potential linearized augmented plane wave method. To address the issue, we have substituted Ti
at the Ru-sites in SrRuO3, where SrRuO3 is a ferromagnetic material. Calculated results establish
Ti4+ valence states (similar to SrTiO3), which was predicted experimentally. Thus, Ti substitution
dilutes the Ru-O-Ru connectivity, which is manifested in the calculated results in the form of
significant band narrowing leading to finite gap between t2g and eg bands. At 75% substitution, a
large gap (> 2eV) appears at the Fermi level, ǫF in the up spin density of states, while the down spin
states contributes at ǫF characterizing the system a half-metallic ferromagnet. The t2g − eg gap can
be tailored judiciously by tuning Ti concentrations to minimize thermal effects, which is often the
major bottleneck to achieve high spin polarization at elevated temperatures in other materials. This
study, thus, provides a novel but simple way to fabricate half-metallicity in ferromagnetic materials,
which are potential candidates for spin-based technology.
PACS numbers: 85.70.Ay, 75.30.-m, 71.70.Ch, 71.15.Ap
The search of half metallic ferromagnetic materials has
seen an explosive growth in the recent times due to its
potential technological applications. In these materials,
the electronic density of states (DOS) at the Fermi level,
ǫF corresponds to only one kind of spin, while the other
spin density of states exhibit an energy gap at ǫF . Thus,
in the polarized condition, electronic conduction strongly
depends on the spin of the charge carriers; the material
is insulating for one kind of spin and metallic for the
other. This unique property makes them ideal candidates
for the development of spin-based electronics. Various
theoretical studies predicted half metallicity in Heusler
alloys [1], double perovskites [2], manganates [3], CrO2
[4], graphene nanoribbons [5] etc. However, experimen-
tal studies on very few materials such as manganates [3]
and CrO2 [4], etc. exhibit half metallicity at low temper-
atures. Thermal fluctuations often lead to a reduction in
spin polarization at elevated temperatures [6] making it
difficult for technological applications.
In this study, we investigate the evolution of the elec-
tronic density of states in SrRu1−xTixO3 as a function
of x. SrRuO3 is a ferromagnetic metal with Curie tem-
perature of 165 K. Spin polarization at ǫF is found to be
negative in the ferromagnetic ground state [7, 8]. SrTiO3,
on the other hand, is a band insulator. Various experi-
mental studies [9, 10] suggest (4+) valence state of Ti in
the intermediate compositions (similar to SrTiO3), which
corresponds to 3d0 electronic configuration. Thus, in ad-
dition to disorder effect, Ti substitution leads to a dilu-
tion of Ru-O-Ru connectivity. Transport measurements
in SrRu1−xTixO3 exhibit a range of novel phase tran-
sitions involving disorder induced correlated metal, An-
derson insulator, correlated insulator and band insulators
[11] for different values of x.
Using ab initio calculations, we find that Ti substitu-
tion at Ru-sites in ferromagnetic SrRuO3 leads to half
FIG. 1: (color online) Crystal structure of SrRu0.25Ti0.75O3.
In order to obtain the structure of SrRuTiO3, we replaced Ti2
by Ru, and all the Ti and Ru sites are made equivalent.
metallicity. Here, reduced Ru-O-Ru connectivity due to
Ti-substitution leads to significant narrowing of Ru 4d
band and thus, the up spin band moves below ǫF . In-
terestingly, the energy gap between t2g and eg bands can
be tuned by Ti-concentration. 75% substituted sample
exhibits gap as high as 2 eV. Experimental realization of
such method on different systems would provide a new
direction in the search of HMFs for spin-based technol-
The electronic density of states of SrRu1−xTixO3 for
x = 0.0, 0.5, 0.75 and 1.0 were calculated using state-
of-the-art full potential linearized augmented plane wave
method (FLAPW) within the local spin density approxi-
mations (LSDA) using WIEN2K software [12]. The crystal
structure of SrTiO3 is cubic with the lattice constant, a =
3.905 Å. SrRuO3 possesses close to cubic structure with
small orthorhombic distortion. This is manifested clearly
by the similar density of states (DOS) of SrRuO3 in real
structure vis-a-vis in the equivalent cubic structure [7].
Ti-substitution in SrRuO3 leads the system towards cu-
http://arxiv.org/abs/0704.0321v1
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FIG. 2: (color online) (a) TDOS, (b) Ti 3d PDOS, (c) Ru 4d
PDOS, (d) O 2p PDOS and (e) Sr 4d PDOS of SrRu1−xTixO3.
Thin and thick solid lines represent DOS corresponding to x
= 0.5 and 0.75, respectively.
bic structure. Thus, we have considered cubic structure
for all the calculations in this study. A typical unit cell
for SrRu0.25Ti0.75O3 is shown in Fig. 1. There are 8 for-
mula units in the unit cell constructed by doubling the
lattice constant of SrTiO3. In order to preserve cubic
symmetry, three types of Ti are considered occupying
corners (Ti1), edge centers (Ti2) and face centered posi-
tions (Ti3). The body centered position is occupied by
Ru. There are three non-equivalent oxygens; O1 forms
the octahedra around Ti1-sites, O2 forms the octahedra
around Ru-sites and the rest of the oxygen positions are
occupied by O3. Thus, the connectivity between Ru-sites
occurs via Ru-O2 bondings. The muffin-tin radii (RMT )
for Sr, Ru, Ti and O were set to 1.16 Å 0.95 Å 0.95 Å and
0.74 Å respectively. The convergence for different calcu-
lations were achieved considering 512 k points within the
first Brillouin zone. The error bar for the energy conver-
gence was set to < 0.25 meV per formula unit. In every
case, the charge convergence was achieved to be less than
10−3 electronic charge.
In Fig. 2, We show the total DOS calculated for
SrRu1−xTixO3 (x = 0.5 and 0.75) and the partial DOS
obtained by projecting the eigenstates onto the Ti 3d,
Ru 4d, O 2p and Sr 4d states. The figure exhibits 5
distinctly separable features. The energy region -1.5 eV
to -5 eV is primarily contributed by O 2p partial DOS
with negligible contributions from other electronic states.
Thus, these contributions are characterized due to the
non-bonding O 2p states. Sr 4d partial DOS shown in
Fig. 2(e) appear above 5 eV. The peak appears to shift
towards higher energy with increasing x. This can be
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89:;<
FIG. 3: (color online) (a) TDOS, (b) Ti 3d and Ru 4d PDOS,
(c) O 2p PDOS and (d) Sr 4d PDOS of SrTiO3 and SrRuO3.
Dashed line represent Sr 4d PDOS rescaled by 20 times.
understood by comparing the same in the end members,
SrTiO3 and SrRuO3 as demonstrated in Fig. 3. Sr 4d
states appear at much higher energies in SrTiO3 com-
pared to that in SrRuO3. One reason for such a large
shift may be related to the shift of the Fermi level to
the top of the O 2p band in SrTiO3. However, the shift
of Sr 4d band in the intermediate compositions, where
the Fermi level is pinned by the occupancy of the Ru 4d
band, indicates that the Madelung potential at Sr-sites
increases with the increase in Ti concentrations.
Ti 3d partial DOS appears 2 eV above the Fermi level.
This clearly demonstrates that the occupancy of Ti 3d
states is essentially zero and hence correspond to Ti4+
valency. Such valence states was predicted in the x-ray
photoemission spectra [9]. This study provides evidence
of such effect theoretically within the effective single par-
ticle approach itself. The width of the Ti 3d t2g band is
significantly small in x = 0.5 sample (∼ 0.65 eV), which
increases to 1.5 eV in x = 0.75 sample and 2.5 eV at x
= 1.0 (see Fig. 3).
Ru 4d partial DOS exhibit three regions. The narrow
and intense feature between the energy range -1.6 to 0.5
eV correspond to the electronic states having t2g sym-
metry. The electronic states above 1.8 eV appears due
to Ru 4d states having eg symmetry. Notably, the O
2p states also contribute in all the three energy regions.
Thus, DOS appearing below -5 eV can be attributed to
the Ru 4d - O 2p bonding states having a large O 2p
character, and the energy region above -1.5 eV are the
anti-bonding states having primarily Ru 4d character.
Most interestingly, both the compounds exhibit metallic
ground state. However, the t2g bandwidth, W reduces
significantly with the increase in x. While W is close to
2.6 eV in SrRuO3, it is about 1.7 eV for x = 0.5 and 0.54
eV for x = 0.75. Such reduction in W is understand-
able as Ti-substitution leads to a significant reduction
in the hopping interaction strength due to the reduced
degree of Ru-O-Ru connectivity. This is clearly evident
in Fig. 1; if we assume homogeneous distribution of Ru
and Ti atoms in the solid, all the RuO6 octahedra are
separated by TiO6 octahedra at x = 0.5. At x = 0.75,
the number of Ru-[O-Ti-O]-Ru connectivity reduces to
half of that at x = 0.5. Subsequently, U/W (U = local
Coulomb interactions strength) will increase significantly
and presumably play a role in the transport properties
in these compositions [11].
In order to understand the bonding of Ru 4d electronic
states with various O 2p states, we compare the Ru 4d t2g
and eg bands with the 2p bands corresponding to O1, O2
and O3 for x = 0.75 and 0.5 sample in Fig. 4(a) and 4(b),
respectively. All the oxygens are equivalent in the x = 0.5
sample. The energy distribution of O2 2p partial DOS is
almost identical in Fig. 4(a) to that observed in Ru 4d
partial DOS. This is expected as the RuO6 octahedra is
formed by O2 atoms only. The width of the O2 2p band
is significantly larger than that of O1 and O3. The most
interesting observation is that the t2g and eg bands are
separated by a distinct energy gap. This gap is already
visible in Ru 4d partial DOS of x = 0.5 sample in Fig.
4(b) and is absent in SrRuO3 as shown in Fig. 3 and in
the literature as well [7, 13].
We calculate the crystal field splitting of the Ru 4d
band by measuring the separation of the center of gravity
of the Ru 4d t2g and eg bands as shown in Fig. 4 by closed
circles in both the compositions. It is evident that crystal
field splitting, ∆ remains almost the same (∼ 2.1 eV) in
both the compositions and is very close to 2 eV found
in SrRuO3. Thus, the large energy gap between the t2g
and eg bands appears purely due to the band narrowing.
Such effect has strong implication in the magnetic phase
as described below.
It is already well established that the magnetic ground
state can be exactly described by these band structure
calculations [7, 14, 15, 16]. Thus, we have calculated the
ground state energies for ferromagnetic arrangement of
moments of the constituents using local spin-density ap-
proximations. Interestingly, the eigen energy for the fer-
romagnetic ground state in x = 0.5 sample is 5.67 meV/fu
lower than the lowest eigen energy for the non-magnetic
solution. This is higher than 1.2 meV/fu observed in
SrRuO3 in real structure and significantly smaller than
30.4 meV/fu observed in the equivalent cubic struc-
ture of SrRuO3. This energy difference between the
non-magnetic and magnetic solutions increases to 33.95
meV/fu in x = 0.75. All these results suggest that the
stability of the ferromagnetic ground state increases with
the decrease in the degree of charge delocalization of the
CD EF GH I J K L M
TU VW XY Z [ \ ] ^
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ij kl
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qrs tu vwxyz{|}
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·¸¹º»¼ ½¾¿À
FIG. 4: (color online) Ru 4d partial DOS with t2g and eg
symmetry are compared with the O 2p partial DOS in (a)
SrRu0.25Ti0.75O3 and (b) SrRu0.5Ti0.5O3.
ÇÈ É Ê Ë Ì
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Energy (eV)
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FIG. 5: (color online) Up and down spin density of stated
corresponding to (a) Ru 4d in SrRu0.5Ti0.5O3, (b) O 2p in
SrRu0.5Ti0.5O3, (c) Ru 4d in SrRu0.25Ti0.75O3, and (d) O
2p in SrRu0.25Ti0.75O3. This figure demonstrates that band
narrowing in Ru 4d band leads to a gap in the up spin channel
leading to half metallicity.
valence electrons.
The spin magnetic moment centered at Ru-sites is
found to be about 0.6 µB in x = 0.5 sample. Inter-
estingly, magnetic moment at the interstitial electronic
states is significantly large (∼ 0.36 µB). The moment
at the O sites is about 0.05 µB. The Ti sites also ex-
hibit very small moment (∼ -0.03 µB). Thus the total
magnetic moment of the solid becomes 1.24 µB per Ru-
atom. This is very similar to that observed (1.2 µB) in
SrRuO3. The magnetic moments increase significantly
with the increase in x. The moments at Ru site becomes
0.88 µB in x = 0.75 sample. The moments of the intersti-
tial states and 2p states at O2 sites also enhance to 0.66
µB and 0.066 µB, respectively. Thus, the total moment
turns out to be 1.99 µB, which is very close to the spin
only value of 2 µB corresponding to Ru 4t
2g electronic
configuration. It is to note here that although the local
moment of the highly extended 4d states is significantly
smaller than the spin only value as opposed to the case
in 3d transition metal oxides [15], Ru 4d moment induces
a large degree of polarization in the interstitial and O 2p
electrons. These results evidently suggest applicability
of Stoner description to capture magnetic properties of
these systems.
In order to investigate the exchange splitting and the
character of density of states in the vicinity of ǫF , we plot
the spin-resolved DOS corresponding to Ru 4d and O 2p
partial DOS in Fig. 5. In the x = 0.5 sample, both the up
and down spin states contribute at ǫF and the exchange
splitting is found to be about 0.47 eV. This is again very
similar to the case in SrRuO3 [7]. The exchange splitting
increases to 0.65 eV in x = 0.75 sample as shown in the
figure. Interestingly, the up spin band moves significantly
below ǫF and the contributions at ǫF appears only due
to the down spin states indicating a half-metallic behav-
ior. No contribution of the up spin states observed in
the total density of states (not shown here). Considering
the paucity of half-metallic materials for various tech-
nological applications, achieving half metallicity in the
ferromagnetic SrRuO3 by Ti-substitution is remarkable.
It is believed that the half metallicity can be achieved
via strong d − d hybridization in Heusler alloys involv-
ing two transition metal elements in the compound [17].
In transition metal oxides, often doping of large amount
of electrons or holes leads to a shift of the Fermi level
towards the energy gap of one spin channel leading to
half metallicity [3]. The primary difficulty to use these
systems in technological applications is the loss of half
metallicity at elevated temperatures, where thermal ex-
citations leads to significant mixing of various spin chan-
nels due to small energy gap at ǫF [6]. In the present
case, mechanism to achieve half metallicity is simple and
easily achievable experimentally. The most important
aspect is that the energy gap between t2g and eg bands
can be tailored judiciously by tuning the composition to
minimize thermal effects.
In summary, we investigate the possibility of fabricat-
ing half metallicity by Ti-substitution at the Ru-sites in
a ferromagnetic material, SrRuO3. The calculated re-
sults using FLAPW method within the local spin density
approximations reveal tetravalency of Ti in all the com-
positions consistent with the experimental predictions.
The Ru 4d band exhibit significant narrowing with the
increase in Ti-substitution; the crystal field splitting re-
mains almost the same across the whole series. Thus,
an energy gap develops between the t2g and eg bands,
which gradually grows with the increase in x. Conse-
quently, the up spin density of states exhibit an energy
gap at the Fermi level, while the down spin states still
contribute leading to half metallicity. Most interestingly,
the t2g − eg gap can be engineered by tuning x and thus
spin mixing effects due to thermal excitations can be min-
imized. This study thus provide a novel but simple way
to fabricate half metallicity in ferromagnetic materials,
which are potential candidates for spin based technol-
ogy. Experimental realization of this method would help
both chemists and physicists to cultivate new materials.
In addition, this study demonstrates that effective sin-
gle particle approaches provide a remarkable description
of the electronic properties of these systems, which are
predicted experimentally.
∗ Electronic mail: [email protected]
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K.H.J. Buschow, Phys. Rev. Lett. 50, 2024-2027 (1983),
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|
0704.0322 | Emergence of spatiotemporal chaos driven by far-field breakup of spiral
waves in the plankton ecological systems | Emergence of spatiotemporal chaos driven by far-field breakup of spiral waves in the
plankton ecological systems
Quan-Xing Liu,1 Gui-Quan Sun,1 Bai-Lian Li,2 and Zhen Jin1, ∗
Department of Mathematics, North University of China,
Taiyuan, Shan’xi 030051, People’s Republic of China
Ecological Complexity and Modeling Laboratory, Department of Botany and Plant Sciences,
University of California, Riverside, CA 92521-0124, USA
(Dated: October 25, 2018)
Alexander B. Medvinsky et al [A. B. Medvinsky, I. A. Tikhonova, R. R. Aliev, B.-L. Li, Z.-S. Lin,
and H. Malchow, Phys. Rev. E 64, 021915 (2001)] and Marcus R. Garvie et al [M. R. Garvie and C.
Trenchea, SIAM J. Control. Optim. 46, 775-791 (2007)] shown that the minimal spatially extended
reaction-diffusion model of phytoplankton-zooplankton can exhibit both regular, chaotic behavior,
and spatiotemporal patterns in a patchy environment. Based on that, the spatial plankton model
is furtherly investigated by means of computer simulations and theoretical analysis in the present
paper when its parameters would be expected in the case of mixed Turing-Hopf bifurcation region.
Our results show that the spiral waves exist in that region and the spatiotemporal chaos emerge,
which arise from the far-field breakup of the spiral waves over large ranges of diffusion coefficients
of phytoplankton and zooplankton. Moreover, the spatiotemporal chaos arising from the far-field
breakup of spiral waves does not gradually involve the whole space within that region. Our results
are confirmed by means of computation spectra and nonlinear bifurcation of wave trains. Finally,
we give some explanations about the spatially structured patterns from the community level.
PACS numbers: 87.23.Cc, 82.40.Ck, 82.40.Bj, 92.20.jm
Keywords: Spiral waves; Spatio-temporal pattern; Plankton dynamics; Reaction-diffusion system
I. INTRODUCTION
There is a growing interest in the spatial pattern dy-
namics of ecological systems [1, 2, 3, 4, 5, 6, 7, 8, 9, 10,
11, 12, 13]. However, many mechanisms of the spatio-
temporal variability of natural plankton populations are
not known yet. Pronounced physical patterns like ther-
moclines, upwelling, fronts and eddies often set the frame
for the biological process. Measurements of the underwa-
ter light field are made with state-of-the-art instruments
and used to calculate concentrations of phytoplankton
biomass (as chlorophyll) as well as other forms of organic
matter. Very high diffusion of the marine environment
would prevent the formation of any stable patch spatial
distribution with much longer life-time than the typical
time of biodynamics. Meanwhile, in addition to very
changeable transient spatial patterns, there also exist
other spatial patterns in marine environment, much more
stable spatial structure associated with ocean fronts, spa-
tiotemporal chaos [10, 11, 14], cyclonic rings, and so
called meddies [15]. In fact, it is significant to create
the biological basis for understanding spatial patterns
of plankton [16]. For instance, the impact of space on
the persistence of enriched ecological systems was proved
in laboratory experiments [17]. Recently, it has been
shown both in laboratory experiments [18] and theoreti-
cally [14, 19, 20, 21] that the existence of a spatial struc-
ture makes a predator-prey system less prone to extinc-
∗Corresponding author; Electronic address: [email protected]
tion. This is due to the temporal variations of the density
of different sub-populations can become asynchronous
and the events of local extinction can be compensated
due to re-colonization from other sites in the space [22].
During a long period of time, all the spiral waves have
been widely observed in diverse physical, chemical, and
biological systems [23, 24, 25, 26]. However, a quite lim-
ited number of documents [11, 12, 27, 28, 29] concern
the spiral wave pattern and its breakup in the ecological
systems.
The investigation of transition from regular patterns
to spatiotemporally chaotic dynamics in spatially ex-
tended systems remains a challenge in nonlinear sci-
ence [14, 23, 30, 31]. In a nonlinear ecology system, the
two most commonly seen patterns are spiral waves and
turbulence (spatio-temporal chaos) for the level of the
community [32]. It has been recently shown that sponta-
neous spatiatemoporal pattern formation is an instrinsic
property of a predator-prey system [11, 14, 33, 34, 35, 36]
and spatiotemporal structures play an important role in
ecological systems. For example, spatially induced speci-
ation prevents the extinction of the predator-prey mod-
els [11, 12, 37]. So far, plankton patchiness has been ob-
served on a wide range of spatial temporal scales [38, 39].
There exist various, often heuristic explanations of the
spatial patterns phenomenon for these systems. It should
be noted that, although conclusive evidence of ecological
chaos is still to be found, there is a growing number of
indications of chaos in real ecosystems [40, 41, 42, 43].
Recently developed models show that spatial self-
structuring in multispecies systems can meet both cri-
teria and provide a rich substrate for community-level
http://arxiv.org/abs/0704.0322v3
mailto:[email protected]
section and a major transition in evolution. In present
paper, the scenario in the spatially extended plankton
ecological system is observed by means of the numeri-
cal simulation. The system has been demonstrated to
exhibit regular or chaostic, depending on the initial con-
ditions and the parameter values [10, 29]. We find that
the far-field breakup of the spiral wave leads to complex
spatiotemporal chaos (or a turbulentlike state) in the spa-
tially extended plankton model (1). Our results show
that regular spiral wave pattern shifts into spatiotempo-
ral chaos pattern by modulating the diffusion coefficients
of the species.
II. MODEL
In this paper we study the spatially extended nutrient-
phytoplankton-zooplankton-fish reaction-diffusion sys-
tem. Following Scheffer’s minimal approach [44], which
was originally formulated as a system of ordinary diff-
ential equation (ODEs) and later developed models [10,
11, 29, 45, 46], as a further investigation, we study a
two-variable phytoplankton and zooplankton model on
the level of the community to describe pattern formation
with the diffusion. The dimensionless model is written
= rp(1 − p)−
1 + bp
h+ dp∇
2p, (1a)
1 + bp
h−mh− f
n2 + h2
+ dh∇
2h, (1b)
where the parameters are r, a, b, m, n, dp, dh, and
f which refer to work in Refs. [10, 11]. The explana-
tion of model (1) relates to the nutrient-phytoplankton-
zooplankton-fish ecological system [see Refs. [10, 29, 44]
for details]. The local dynamics are given by
g1(p, h) = rp(1− p)−
1 + bp
h, (2a)
g2(p, h) =
1 + bp
h−mh− f
n2 + h2
. (2b)
From the earlier results [45] about non-spatial system
of model (1) by means of numerical bifurcation analysis
show that the bifurcation and bistability can be found in
the system (1) when the parameters are varied within a
realistic range. For the fixed parameters (see the caption
of Fig. 1 and 2), we can see that the f controls the dis-
tance from Hopf bifurcation. For larger f , there exists
only one stable steady state. As f is decreased further,
the homogeneous steady state undergoes a saddle node
bifurcation (SN), that is fSN = 0.658. In this case, a
stable and an unstable steady state become existence.
Moreover, the bistability will emerge when the parame-
ter f lies the interval fSN > f > fc = 0.445 (this value is
more than the Hopf onset, fH = 0.3397). There are three
steady states: with these kinetics A and C are linearly
stable while B is unstable. Outside this interval, the sys-
tem (1) has unique nontrivial equilibrium. Recent stud-
ies [11, 29] shown that the systems (1) can well-develop
the spiral waves in the oscillation regime, but where the
authors only consider the special case, i.e., dp = dh. A
few important issue have not yet been properly addressed
such as the spatial pattern if dp 6= dh.
Here we report the result that emergence of spatiotem-
poral chaos due to breakup in the system under the
dh 6= dp case. We may now use the f and diffusion
ratio, ν = dh/dp, as control parameters to evaluate
the region for the spiral wave. Turing instability in
reaction-diffusion can be recast in terms of matrix sta-
bility [47, 48]. Such with the help of Maple software
assistance algebra computing, we obtain the parameters
space (f, ν) bifurcation diagrams of the spiral waves as
showing Fig. 2, in which two lines are plotted, Hopf line
(solid) and Turing lines (dotted) respectively. In domain
I, located above all three bifurcation lines, the homo-
geneous steady states is the only stable solution of the
system. Domain II are regions of homogeneous oscilla-
tion in two dimensional spaces [49]. In domain III, both
Hopf and Turing instabilities occur, (i.e., mixed Turing-
Hopf modes arise), in which the system generally pro-
duces the phase waves. Our results show that the system
has spiral wave in this regions. One can see that a Hopf
bifurcation can occur at the steady when the parameter
f passes through a critical values fH while the diffusion
coefficients dp = dh = 0 and the bifurcation periodic so-
lutions are stable. From our analysis (see Fig. 2), one
could also see that the diffusion can induce Turing type
instability for the spatial homogeneous stable periodic
solutions and the spatially extended model (1) exhibit
spatio-temporal chaos patterns. These spatial pattern
formation arise from interaction between Hopf and Tur-
ing modes, and their subharmonics near hte codimension-
two Hopf-Turing bifucation point. Special, it is interest-
ing that spiral wave and travelling wave will appear when
the parameters correspond to the Turing-Hopf bifurca-
tion region III in the spatially extended model (1), i.e.,
the Turing instability and Hopf bifurcation occur simul-
taneously.
III. NUMERICAL RESULTS
The simulation is done in a two-dimensional (2D)
Cartesian coordinate system with a grid size of 600×600.
The fourth order Runger-Kutta integrating method is
applied with a time step ∆t = 0.005 time unit and a
space step ∆x = ∆y = 0.20 length unit. The results
remain the same when the reaction-diffusion equations
were solved numerically in one and two spatial dimen-
sions using a finite-difference approximation for the spa-
tial derivatives and an explicit Euler method for the time
integration. Neumann (zero-flux) boundary conditions
FIG. 1: The sketch map for the bistability and the Hopf bi-
furcation in the system (2) with r = 5.0, a = 5.0, b = 5.0,
m = 0.6, and n = 0.4. The black curve is the g1(p, h). The
colored curves are g2(p, h) with different values of f . The red
curve: f = 0.3; the blue: f = 0.445; the green: f = 0.5; and
the cyan: f = 0.658.
5 10 15
Turing instability
FIG. 2: The sketch map of parameter space (f, ν) bifurcation
diagrams for the spatially extended system (1) with r = 5.0,
a = 5.0, b = 5.0, m = 0.6, dp = 0.05, and n = 0.4.
were emmployed in our simulation. The diffusion terms
in Eqs. (1a) and (1b) often describe the spatial mixing
of species due to self-motion of the organism. The typi-
cal diffusion coefficient of plankton patterns dp is about
0.05, based on the parameters estimatie of Refs [50, 51]
using the relationship between turbulent diffusion and
the scale of the space in the sea. In the previous stud-
ies [10, 11, 29, 45, 46], the authors provided a valueable
insight into the role of spatial pattern for the system (1)
if dp = dh. From the biological meaning, the diffusion
coefficients should satisfy dh ≥ dp. However, in nature
waters it is turbulent diffusion that is supposed to domi-
nate plankton mixing [52], when dh < dp is allowed. The
other reason for choosing such parameter is that it is well-
known new patterns, such as Turing patterns, can emerge
in reaction-diffusion systems in which there is an imbal-
ance between the diffusion coefficients dp and dh [23, 53].
Therefore, we set ν = dh/dp, and investigated whether a
spiral wave would break up into complex spatiotemporal
chaos when the diffusion ratio was varied. Throughout
this paper, we fix dp = 0.05 and dh is a control parameter.
In the following, we will show that the dynamic behav-
ior of the spiral wave qualitatively change as the control
parameter dh increases from zero, i.e., the diffusion ra-
tio ν increases from zero, to more than one. For large
ν (ν > 1), the outwardly rotating spiral wave is com-
pletely stable everywhere, and fills in the space when the
proper parameters are chosen, as shown in Fig. 3(A). Fig-
ure 3(A) shows a series of snapshots of a well-developed
single spiral wave formed spontaneously for the variable
p in system (1). The spiral is initiated on a 600×600 grid
by the cross-field protocol (the initial distribution chosen
in the form of allocated “constant-gradient” perturbation
of the co-existence steady state) and zero boundary con-
ditions are employed for simulations in the two dimen-
sions. From Fig. 3(A) we can see that the well-developed
spiral waves are formed firstly by the evolution. Inside
the domain, new waves emerge, but are evolved by the
spiral wave growing from the center. The spiral wave
can steadily grow and finally prevail over the whole do-
main (a movie illustrating the dynamical evolution for
this case [54] [partly movie−1, movie−2, and movie−3
for dh = 0.2]). Fig. 3(B) shows that the spiral wave
first break up far away from the core center and even-
tually relatively large spiral fragments are surrounded
by a ‘turbulent’ bath remain. The size of the surviv-
ing part of the spiral does not shrink when dh is further
decreasing until finally dh equals to 0, which is different
from phenomenon that is observed previous in the two-
dimensional space Belousov-Zhabotinsky and FitzHugn-
Nagumo oscillatory system [30, 31, 55, 56, 57], in which
the breakup gradually invaded the stable region near the
core center, and finally the spiral wave broke up in the
whole medium. Figure 3(C) is the time sequences (ar-
bitrary units) of the variables p and h at an arbitrary
spatial point within the spiral wave region, from which
we can see that the spiral waves are caused by the ac-
cepted as “phase waves” with substantially group veloc-
ity, phase velocity and sinusoidal oscillation rather than
the relaxational oscillation with large amplitude. This
breakup scenario is similar to the breakup of rotating
spiral waves observed in numerical simulation in chemi-
cal systems [30, 31, 55, 56, 57], and experiments in BZ
systems [58, 59], which shows that spiral wave breakup
in these systems was related to the Eckhaus instability
and more important, the absolute instability.
The corresponding trajectories of the spiral core and
the spiral arm (far away from the core center) at y = 300
are shown in Fig. 4, respectively. From Fig. 4, we can
see that the spiral core is not completely fixed, but oscil-
lates with a large amplitude. However, as dh decreases
to a critical value, an unstable modulation develops in
200 220 240 260 280 300
(D) t (arb. units)
FIG. 3: Well developed spiral waves and some properties of
them. The figures show simulations of the system (1) with
r = 5, a = 5, b = 5, m = 0.6, n = 0.4, dp = 0.05, and
f = 0.3. (A)Well developed spiral waves shown at subsequent
snapshot in time, dh = 0.2. (B) Far-field breakup of the spiral
waves shown at subsequent snapshot in time, dh = 0.002.
The white (black) areas correspond to maximum (minimum)
values of p [Additional movie format available from Ref. [54]].
(C) Oscillations of the variable p and h at an arbitrary spatial
point within the regular spiral wave region for both scenarios.
Each figure is ran the long time until it spatial patterns are
unchange.
regions which is far away from the spiral core (cf. the
middle column of the Fig. 4). These oscillations eventu-
ally grow large enough to cause the spiral arm far away
from the core to breakup into complex multiple spiral
waves, while the core region remains stable (the corre-
sponding movie can be viewed in the online supplemen-
tal in Ref. [54] [partly movie−1 and movie−2, and for
dh = 0.02]). Figures 3(B) and 4(B) show the dynamic
behavior for dh = 0.02, i.e., ν = 0.4. The regular tra-
jectories far away from the core are now the same as the
region of the spatial chaos (cf. the middle column of the
Fig. 4). It is shown that an decrease in the diffusion ra-
tio ν which leads to population oscillations of increasing
amplitude (cf. the left column of the Fig. 4). In the
tradition explain that the minimum value of the popula-
tion density decreases and population extinction becomes
more probable due to stochastic environmental perturba-
tions. However, from the spatial evolution of system (1)
(see Fig. 3), the temporal variations of the density of
different sub-population can become asynchronous and
the events of local extinction can be compensated due to
re-colonization (or diffusion) from other sites.
FIG. 4: The corresponding trajectories (from left to right)
for locations (300, 300), (250, 300), and (50, 300) respectively.
The parameters in (A), and (B) were the same as these in
Fig. 3(A) and (B), respectively.
Furthermore, it is well known that the basic arguments
in spiral stability analysis can be carried out by reducing
the system to one dimensional space [30, 31, 55, 56, 57].
Here we show some essential properties of the spiral
breakup resulting from the numerical simulation. In the
next section we will give the theoretical computation by
using the eigenvalue spectra. In this model, it is worth
noting that we do not neglect the oscillation of the dy-
namics in the core as shown in Fig. 4 due to the system
exhibiting spatial periodic wave trains when the model
is simulated in one-dimensional space. Breakup occurs
first far away from the core (the source of waves). The
spiral wave breaks towards the core until it gets to some
constant distance and then the surviving part of the spi-
ral wave stays stable. These minimal stable wavelengths
are called λmin. So the one-parameter family may be
described by a dispersion curve λ(dh) (see Fig. 5). The
minimal stable wavelength λmin of the spiral wave are
shown in Fig. 5 coming from the simulation in two di-
mensional space. The results of Fig. 5 can be interpreted
as follows: the minimal stable wavelengths decrease with
respect to the decrease of dh but eventually stay at a
relative constant value, which is that the stable spiral
waves are always existing for a larger region values of dh.
Space-time plots at different times are shown in Fig. 6
for two different dh, i.e., different ν, which display the
time evolution of the spiral wave along the cross section
in the two-dimensional images of Fig. 3(A) and (B). As
shown in Fig. 6(A) and (B) for dh = 0.2 and dh = 0.02
respectively, the waves far away from the core display
unstable modulated perturbation due to convective in-
stability [30, 31, 55, 56, 57], but this perturbation is
gradually advected to the left and right sides, and finally
disappears. The instability manifests itself to produce
the wave train breakup several waves from the far-field,
as shown in Figs. 6(B).
FIG. 5: Dependence of the wavelength λmin on the parameter
dh for the system (1) with r = 5.0, a = 5.0, b = 5.0, m = 0.6,
dp = 0.05, and n = 0.4. Note the log scale for dh.
IV. SPECTRA AND NONLINEAR
BIFURCATION OF THE SPIRAL WAVE
In this section, we concentrate on the linear stabil-
ity analysis of spiral wave by using the spectrum the-
ory [56, 60, 61, 62, 63]. From the results in Refs. [56, 62]
we know that the absolute spectrum must be computed
numerically for any given reaction-diffusion systems. In
practice, such computations only require discretization
in one-dimensional space and compare with computing
eigenvalues of the full stability problem on a large do-
main due to the spiral wave exhibitting traveling waves
in the plane (see Fig. 6 about the space-time graphes).
For spiral waves on the unbounded plane, the essential
FIG. 6: Space-time plots of variable p for different time and
dh. The parameters in (A), and (B) are the same as those in
Fig. 3(A) and (B), respectively.
spectrum is also required to compute, since it determined
only by the far-field wave trains of the spiral. The lin-
ear stability spectrum consists of point eigenvalues and
the essential spectrum that is a continuous spectrum for
spiral waves.
For sake of simplicity, the Eqs. (1a) and (1b) can been
written as following
= dp∇
2p+ g1(p, h), (3a)
= dh∇
2h+ g2(p, h). (3b)
Suppose that (p∗, h∗) are a solutions and refer to them
as steady spirals of Eq. (3) that rotate rigidly with a
constant angular velocity ω, and that are asymptotically
periodic along rays in the plane. In a coratating coordi-
nate frame, using the standardized analysis method for
the spiral waves [62, 63], the Eq. (3) is given by
= dp∇
ρ,θp+ ω
+ g1(p
∗, h∗), (4a)
= dh∇
ρ,θh+ ω
+ g2(p
∗, h∗), (4b)
where (ρ, θ) denote polar coordinates, spirals waves are
relative equilibria, then the statianry solutions p∗(ρ, θ)
and h∗(ρ, θ) both are 2π-periodic functions with θ = ϕ−
ωt. In Eqs. (4a) and (4b) the operator∇2ρ,θ denotes ∂ρρ+
A. Computation of spiral spectra
Next, we commpute the leading part of its linear stabil-
ity spectrum for the system (4). Consider the linearized
evolution equation in the rotating frame, the eigenvalue
problem of Eqs. (4a) and (4b) associated with the planar
spiral solutions p∗(ρ, θ) and h∗(ρ, θ) are given by
ρ,θp+ ω
∗, h∗)p+ gh1 (p
∗, h∗)h = λp, (5a)
ρ,θh+ ω
∗, h∗)p+ gh2 (p
∗, h∗)h = λh, (5b)
where g
1 , · · · , g
2 denote the derivatives of the nonlin-
ear functions and g
1(p, h) = r(1 − p) − rp −
(1+bp)2
, gh1 (p, h) = −
2(p, h) =
− abph
(1+bp)2
, and
gh2 (p, h) =
−m− 2fnh
n2+h2
+ 2fnh
(n2+h2)2
. We shall ignore
isolated eigenvalues that belong to the point spectrum,
instabilities caused by point eigenvalues lead to mean-
deringor drifting waves, or to an unstable tip motionin
in excitable media and oscillation media [56, 64, 65, 66].
This phenomenon is not shown in the present paper. In-
stead, we focus on the continuous spectrum that is re-
sponsible for the spiral wave breakup in the far field (see
Fig. 3(b)). By the results in Ref. [62], it turns out that
the boundary of the continuous spectrum depends only
on the limiting equation for ρ → ∞. Thus, we have that
λ is the boundary of the continuous spectrum if, and only
if the limiting equation
ρ,ρp+ ω
∗, h∗)p+ gh1 (p
∗, h∗)h = λp, (6a)
ρ,ρh+ ω
∗, h∗)p+ gh2 (p
∗, h∗)h = λh, (6b)
have solutions p(ρ, θ) and h(ρ, θ) for (ρ, θ) ∈ R+× [0, 2π],
which are bounded but does not decay as ρ → ∞. Since
spiral waves are rotating waves in the plane, the wave
train solutions have the form as u(t, x, y) = u(ρ, ϕ− ωt)
for an appropriate wave numbers k and temporal fre-
quency ω, where we assume that u is 2π-periodic in
its argument so that u(ξ) = u(ξ + 2π) for all ξ and
u = (p, h)T. Spiral waves converge to wave trains
u(ρ, ϕ − ωt) → uwt(kρ + ϕ − ωt) for ρ → ∞, which
are corresponding to asymptotically Archimedean in the
two-dimensional space. Assume that k 6= 0 and ω 6= 0,
and in this case, we can pass from the theoretical frame ρ
to the comoving frame ξ = kρ+ϕ−ωt (ξ ∈ R) in which
the eigenvalue equation (6) becomes
2∇2ξ,ξp+ωpξ+g
1(uwt(ξ))p+g
1 (uwt(ξ))h = λp, (7a)
2∇2ξ,ξh+ ωhξ + g
2(uwt(ξ))p + g
2 (uwt(ξ))h = λh.(7b)
Indeed, any nontrivial solution u(ξ) = (p(ξ), h(ξ))T cor-
responding to the linearization eigenvalue problem (7)
give a solution U(ρ, ·) of the eigenvalue problem for the
temporal period map of (3) in the corotating frame via
U(ρ, ·) = eλtu(kρ− ωt), U(ρ, T ) = eλTu(kρ− 2π).
We write the equations (7) as the first-order systems
= p1,
= h1,
= k−2d−1p
µp− ωp1 − g
1(uwt(ξ))p− g
1 (uwt(ξ))h
= k−2d−1
µh− ωh1 − g
2(uwt(ξ))p− g
2 (uwt(ξ))h
in the radial variable ρ. Then the spatial eigenvalues or
spatial Floquet exponents are deternined as the roots of
the Wronskian
A(λ, k) :=
0 0 1 0
0 0 0 1
(λ− g
1(uwt(ξ))) −
gh1 (uwt(ξ)) −
2(uwt(ξ))
(λ− gh2 (uwt(ξ))) 0 −
where k ∈ R. The function U(ρ, ·) = eλteikρu0(kρ− ωt)
satisfies the equation (3) when the spatial and temporal
exponents ik and λ satisfy the complex dispersion rela-
tion det(A(λ, k) − ik) = 0 for λ ∈ C. We call the ik
in spectrum of A(λ, k) as spatial eigenvalues or spatial
Floquet exponents.
The stability of the spiral waves state (p∗, h∗) on the
plane is determined by the essential spectrum given by
Σess = {λ ∈ C; det(A(λ, k) − ik) = 0 for some k ∈ R}.
Now, we compute the continuous spectrum with the
equation (9) that are parameterized by the wave num-
ber k. For each λ, there are infinitely many stable and
unstable spatial eigenvalues. We plot λ in the complex
plane associated spatial spectrum, see Fig. 7. By the ex-
plaination of Sandstede et al [60], one would know that
if the real part of the essentail spectra is positive, then
the associated eigenmodes grow exponentially toward the
boundary, i.e., they correspond to a far-field instability.
Note that we find the essentail spectra are not sensitive
to temporal frequency, ω.
Re(λ)
K30 K20 K10 0
Im(λ)
Re(λ)
K0.8 K0.6 K0.4 K0.2 0 0.2
Im(λ)
FIG. 7: The essentail spectra of wave trains are obtained by
using the algorithms outlined in Refs. [60, 61]. The param-
eters of (A) and (B) are corresponding to the values used in
the simulations of Fig. 3(A) and (B).
B. Existence and properties of wave trains
Suppose that a reaction-diffusion system on the one-
dimensional space such that the variables equal to a
homogeneous stationary solution. If the homogeneous
steady-state destabilizes, then its linearization accommo-
dates waves of the form ei(kx−ωt) for certain values k and
ω. Typically, near the transition to instability, small spa-
tially periodic travelling waves arise for any wave number
close to kc, which is the critical wavenumber. Their wave
speed is approximately equal to ωc
, where ωc is corre-
sponding to kc. In present paper, we focus exclusively on
the situation where ωc = 0 and kc 6= 0. The bifurcation
with ωc = 0 and kc 6= 0 is known as the Turing bifur-
cation, and the bifurcating spatially periodic steady pat-
terns are often referred to as Turing patterns. Another
class of moved patterns will appear when the instabilities
modulated by Hopf-Turing bifurcation, which is resem-
ble a travelling waves. Moreover, the common feature
of the spiral waves in one-dimensional space mentioned
above is the presence of wave trains which are spatially
periodic travelling waves of the form pwt(kx−ωt; k) and
hwt(kx − ωt; k), where pwt(φ; k) and hwt(φ; k) are 2π-
periodic about φ. Typically, the spatial wavenumber k
and the temporal frequency ω are related via the non-
linear dispersion relation ω = ω(k) so that the phase
velocity is given by
. (12)
A second quantity related to the nonlinear dispersion
relation is the group velocity, cg =
, of the wave
train which also play a central role in the spiral waves.
The group velocity cg gives the speed of propagation of
small localized wave-package perturbations of the wave
train [67]. Here, we are only concerned the existence of
travelling wave solution. In fact, the spiral waves move
at a constant speed outward from the core (see Fig. 6),
so that they have the mathematical form p(x, t) = P (z),
and h(x, t) = H(z) where z = x−cpt. Substituting these
solution forms into Eq. (3) gives the ODEs
+ g1(P,H) = 0, (13a)
+ g2(P,H) = 0. (13b)
Here, we investigate numerically the existence, speed
and wavelength of travelling wave patterns. Our ap-
proach is to use the bifurcation package Matcont 2.4 [68]
to study the pattern ODEs (13). To do this, the most
natural bifurcation parameters are the wave speed cp and
f , but they give no information about the stability of
travelling wave as solutions of the model PDEs (3).
Our starting point is the homogeneous steady state of
Eq. (13) with in the domain III of Fig. 2. The typical bi-
furcation diagrams are illustrated in Fig. 8, which shows
that steady spatially peroidic travelling waves exist for
the larger values of the speed cp, but it is unstable for
small values of cp. The changes in stability occur via
Hopf bifurcation, from which a branch of periodic orbits
emanate. Note that here we use the terms “stable” and
“unstable” as referring to the ODEs system (13) rather
than the model PDEs. Fig. 8(B) illustrates the max-
imun stable wavelength against the bifucation parame-
ter, speed cp, and the small amplitudes have very long
wavelength. It is known that cp =
, hence the tavelling
wave solution exist when the cp 6= 0, i.e., k 6= 0, ω 6= 0.
Using Matcont 2.4 package, it is possible to track the lo-
cus of the Hopf bifurcation points and the Limit point
(fold) bifurcation in a parameter plane, and a typical ex-
ample of this for the cp-f and cp-dh plane are illustrated
in Fig. 9. The travelling wave solutions exist for values
of cp and f lying in left of Hopf bifurcation locus (see
Fig. 9(A)). The same structure about the cp-dh plane
is shown in Fig. 9(B). These reuslts confirm our previ-
ous analysis coming from the algebra computation (see
Fig. 2) and the numerical results (see Fig. 6).
V. CONCLUSIONS AND DISCUSSION
We have investigated a spatially extended plank-
ton ecological system within two-dimensional space and
0 0.1 0.2 0.3 0.4 0.5 0.6 0.7 0.8
Speed, c
Hopf bifurcation point
FIG. 8: Typical bifurcation diagrams for the pattern
ODEs (13). (A) The spatially periodic travelling waves of
system (3) is existence. The changes in stability occur via
Hopf bifurcation, from which a branch of periodic orbits em-
anate. Thus unstable travelling waves appear. (B) Maxi-
mum stable wavelength along the bifurcation parametercp,
i.e., k 6= 0, ω 6= 0. The parameter values in (A) and (B) are
the same as Fig. 3(A).
found that its spatial patterns exhibit spiral waves dy-
namics and spatial chaos patterns. Specially, the sce-
nario of the spatiotemporal chaos patterns arising from
the far-field breakup is observed. Our research is based
on numerical analysis of a kinematic mimicking the dif-
fusion in the dynamics of marine organisms, coupled to a
two component plankton model on the level of the com-
munity. By increasing (decreasing) the diffusion ratio of
the two variables, the spiral arm first broke up into a
turbulence-like state far away from the core center, but
which do not invade the whole space. From the previous
studies in the Belousov-Zhabotinsky reaction, we know
the reason causing this phenomenon can be illuminated
theoretically by the M. Bär and L. Brusch [30, 31], as
well as by using the spectrum theory that poses by B.
Sandstede, A. Scheel et al [56, 60, 61, 69]. The far-field
breakup can be verified in field observation and is useful
to understand the population dynamics of oceanic ecolog-
ical systems. Such as that under certain conditions the
interplay between wake (or ocean) structures and bio-
logical growth leads to plankton blooms inside mesoscale
hydrodynamic vortices that act as incubators of primary
production. From Fig. 3 and corresponding the movies,
we see that spatial peridic bloom appear in the phyto-
plankton populations, and the details of spatial evolution
of the distribution of the phytoplankton population dur-
ing one bloom cycle, respectively.
In Ref. [70], the authors study the optimal control of
the model (1) from the spatiotemporal chaos to spiral
waves by the parameters for fish predation treated as a
multiplicative control variable. Spatial order emerges in a
range of spatial models of multispecies interactions. Un-
surprisingly, spatial models of multispecies systems often
0 0.5 1 1.5 2 2.5
Speed, c
Locus of Hopf bifurcation points
0.5 1 1.5 2
Speed, c
Locus of Hopf bifurcation points
FIG. 9: An illustration of the variations in parameter space of
the pattern ODEs (13). We plot the loci of Hopf bifurcation
points. (A) f − cp planes; (B) dh − cp planes. The parameter
values in (A) and (B) are the same as Fig. 3(A).
manifests very different behaviors from their mean-field
counterparts. Two important general features of spatial
models of multispecies systems are that they allow the
possibility of global persistence in spite of local extinc-
tions and so are usually more stable than their mean-field
equivalents, and have a tendency to self-organzie spa-
tially or regular spatiotemporal patterns [70, 71]. The
spatial structures produces nonrandom spatial patterns
such as spiral waves and spatiotemporal chaos at scales
much larger than the scale of interaction among individ-
uals level. These structures are not explicitly coded but
emerge from local interaction among individuals and lo-
cal diffusion.
As we know that plankton plays an important role in
the marine ecosystem and the climate, because of their
participation in the global carbon and nitrogen cycle at
the base of the food chain [72]. From the review [73], a
recently developed ecosystem model incorporates differ-
ent phytoplankton functional groups and their competi-
tion for light and multiple nutrients. Simulations of these
models at specific sites to explore future scenarios sug-
gest that global environmental change, including global-
warming-induced changes, will alter phytoplankton com-
munity structure and hence alter global biogeochemical
cycles [74]. The coupling of spatial ecosystem model to
global climate raises again a series of open questions on
the complexity of model and relevant spatial scales. So
the study of spatial model with large-scale is more impor-
tant in the ecological system. Basing on numerical simu-
lation on the spatial model, we can draft that the oceanic
ecological systems show permanent spiral waves and spa-
tiotemporal chaos in large-scale over a range of parame-
ter values dh, which indicates that periodically sustained
plankton blooms in the local area. As with all areas of
evolutionary biology, theoretical development advances
more quickly than does empiraical evidence. The most
powerful empirical approach is to conduct experiments
in which the spatial pattern can be measured directly,
but this is difficulties in the design. However, we can in-
directly measured these phenomenona by the simulation
and compared with the satellite pictures. For example,
the spatiotemporal chaos patterns agree with the per-
spective observation of the Fig. 3 in Ref. [73]. Also, some
satellite imageries [http://oceancolor.gsfc.nasa.gov] have
displayed spiral patterns that represent the phytoplank-
ton [the chlorophyll] biomass and thus demonstrated that
plankton patterns in the ocean occur on much broader
scales and therefore mechanisms thought diffusion should
be considered.
Acknowledgments
This work is supported by the National Natural Sci-
ence Foundation of China under Grant No. 10471040
and the Natural Science Foundation of Shan’xi Province
Grant No. 2006011009.
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|
0704.0323 | General sequential quantum cloning | General Sequential Quantum Cloning
Gui-Fang Dang and Heng Fan
Institute of Physics, Chinese Academy of Sciences, Beijing 100080, China.
(Dated: November 4, 2018)
Some multipartite quantum states can be generated in a sequential manner which may be im-
plemented by various physical setups like microwave and optical cavity QED, trapped ions, and
quantum dots etc. We analyze the general N to M (N ≤ M) qubits Universal Quantum Cloning
Machine (UQCM) within a sequential generation scheme. We show that the N to M sequential
UQCM is available. The case of d-level quantum states sequential cloning is also presented.
PACS numbers: 03.67.Mn, 03.65.Ud, 52.50.Dv
Quantum entanglement plays a key role in quantum
computation and quantum information [1]. Multipartite
entangled states arise as a resource for quantum infor-
mation processing tasks such as the well known quantum
teleportation[2], quantum communication [3, 4], clock
synchronization [5] etc. In general it is extremely dif-
ficult to generate experimentally multipartite entangled
states through single global unitary operations. In this
sense, the sequential generation of the entangled states
appears to be promising. Actually most of the quantum
computation networks are designed to implement quan-
tum logic gates through a sequential procedure [6]. Re-
cently sequential implementing of quantum information
processing tasks has been attracting much attention. It
is pointed out that photonic multiqubit states can be
generated by letting a source emit photonic qubits in a
sequential manner [7]. The general sequential generation
of entangled multiqubit states in the realm of cavity QED
was systematically studied in Refs.[8, 9]. It is also shown
that the class of sequentially generated states is identical
to the matrix-product-state (MPS) which is very useful
in study of spin chains of condensed matter physics [10].
On the other hand, much progress has already been
made in the past years in studying quantum cloning ma-
chines, for reviews see, for example, Refs.[11, 12, 13].
And various quantum cloning machines have been im-
plemented experimently by polarization of photons [14,
15, 16, 17, 18],nuclear spins in Nuclear Magnetic Reso-
nance [19, 20], etc. However, these experiments are for
1 to 2 (one qubit input and two-qubit output) or 1 to 3
cloning machines. The more general case will be much
difficult. There are some schemes proposed for the gen-
eral quantum cloning machines which are not in a sequen-
tial manner, see for example, [21, 22]. Recently a 1 to M
sequential universal quantum cloning is proposed [23] by
using the cloning transformation presented in Ref.[24].
Since it is in a sequential procedure, potentially it re-
duces the difficult in implementing this quantum cloning
machine. However, as is well known the collective quan-
tum cloning machine (the N identical input states are
cloned collectively to M copies) is better than the quan-
tum cloning machine which can only deal with the in-
dividual input(only one input is copied to several copies
each time). We know that the general N to M cloning
transformation is also available in Refs.[24, 25]. Then a
natural question arise is that whether the general N to
M sequential cloning machine is possible. In this Letter,
we will present the general sequential universal quantum
cloning machine.
The 1 to M cloning transformations used in Ref.[23]
was proposed by Gisin and Massar in Ref.[24]. And the
N toM UQCM was also presented in Ref.[24]. However,
to use the method proposed in Refs.[8, 23] to find the se-
quential cloning machine, the input state |Φ〉⊗N should
be expanded in computational basis {|0〉, |1〉}. The ex-
plicit quantum cloning transformations with this kind of
input were proposed by Fan et al in Ref.[25]. In this Let-
ter, based on the result of Ref.[25], the general sequential
UQCM will be presented.
As presented in Refs.[8, 23], the sequential generation
of a multiqubit state is like the following. Let HA be
a D-dimensional Hilbert space which acts as the ancil-
lary system, and a single qubit (e.g., a time-bin qubit)
is in a two-dimensional Hilbert space HB. In every step
of the sequential generation of a multiqubit state, a uni-
tary time evolution will be acting on the joint system
HA ⊗HB. We assume that each qubit is initially in the
state |0〉 which is like a blank or an empty state and
will not be written out in the formulas. So the unitary
time evolution is written in the form of an isometry V :
HA → HA⊗HB, where V =
i,α,β V
α,β |α, i〉〈β|, each V i
is a D×D matrix, and the isometry condition takes the
i=0 V
i†V i = 1. By applying successively n oper-
ations of V (not necessarily the same) on an initial ancil-
lary state |φI〉 ∈ HA, we obtain |Ψ〉 = V [n]...V [2]V [1]|φI〉.
The generated n qubits are in general an entangled state,
but the last step qubit-ancilla interaction can be chosen
so as to decouple the final multiqubit entangled state
from the auxiliary system, so the sequentially generated
state is
|ψ〉 =
i1...in=0
〈φF |V [n]in ...V [1]i1 |φI〉|in, ..., i1〉, (1)
where |φF 〉 is the final state of the ancilla. This is the
MPS. It was proven that any MPS can be sequentially
generated [8].
http://arxiv.org/abs/0704.0323v2
Suppose there are N identical pure quantum states
|Φ〉⊗N = (x0|0〉+x1|1〉)⊗N need to be cloned toM copies,
where |x0|2 + |x1|2 = 1. We know that the input state
can be represented by a basis in symmetric subspace.
|Φ〉⊗N =
xN−m0 x
CmN |(N −m)0,m1〉, (2)
where |(N − m)0,m1〉 denotes the symmetric and nor-
malized state with (N −m) qubits in the state |0〉 and m
qubits in the state |1〉, and we have CmN = N !/(N−m)!m!
in standard notation. So if we find the quantum cloning
transformations for all states in symmetric subspace, we
can clone N pure states to M copies. The UQCM with
input in symmetric subspace can be written as [25],
|(N −m)0,m1〉 → |ΦmM 〉, (3)
where
|ΦmM 〉 =
βmj |(M −m− j)0, (m+ j)1〉 ⊗Rj ,(4)
βmj =
M−N−j
M−m−jC
(m+j)
/CN+1M+1, (5)
where Rj are the ancillary states of the cloning machine
and are orthogonal with each other for different j. For
a sequential quantum cloning machine in this Letter, we
choose a realization Rj ≡ |(M −N − j)1, j0〉 for the an-
cilla states. This UQCM is optimal in the sense that
the fidelity between single qubit output state reduced
density operator ρoutreduced and the single input |Φ〉 is op-
timal. The optimal fidelity is F = 〈Φ|ρoutreduced|Φ〉 =
(MN +M + N)/M(N + 2), see Refs.[11, 12, 13] for re-
views and the references therein. A realization of this
UQCM with photon stimulated emission can be found in
Ref.[22] which is not in a sequential manner. We next
show that this general N to M UQCM can be generated
through a sequential procedure.
The basic idea is to show that the final state of the
cloning, |ΦmM 〉 in (4), can be expressed in its MPS form.
As shown in Ref.[8], any MPS can be sequentially gen-
erated. We shall follow the method, for example, as in
Refs.[23, 26]. By Schmidt decomposition, we first ex-
press the quantum state |ΦmM 〉 as a bi-partite state across
1 : 2... cut,
|ΦmM 〉 = λ
1 |0〉|φ
[2...(2M−N)]
1 〉+ λ
2 |1〉|φ
[2...(2M−N)]
Γ[1]i1α1 λ
|i1〉|φ[2...(2M−N)]α1 〉, (6)
where Γ
α1 = δα1,1,Γ
α1 = δα1,2, and λ
α1 are eigen-
values of the first qubit reduced density operator, and
we find λ
∑M−m−1
k=−m β
M−1/C
M , λ
∑M−m−1
k=−m β
mk+1C
M−1/C
m+k+1
M . To correspond with
the MPS in (1), we can define V
[1]i1
α1 = Γ
[1]i1
α1 . Suc-
cessively by Schmidt decomposition, the quantum state
|ΦmM 〉 in (4) is divided into a bi-partite state with the first
n qubits as one part, and the rest as another part, where
1 < n ≤M − 1. We find
|ΦmM 〉 =
j+1|(n− j)0, j1〉|φ
[(n+1)...(2M−N)]
j+1 〉, (7)
when 1 < n ≤M−N+m,n′ = n; whenM−N+m< n ≤
M − 1, n′ =M −N +m, λ[n]j+1 are eigenvalues of the first
n qubits reduced density operator of |ΦmM 〉. According to
the results in Eqs.(4,5), we can obtain,
j+1 =
M−m−n
m(j+k)
Cm+kM−n
m+j+k
. (8)
And we also have
|φ[(n+1)...(2M−N)]j+1 〉 =
M−m−n
β2m(j+k) ×
(m+k)
m+j+k
|(M − n−m− k)0, (m+ k)1〉 ⊗Rj+k.
By induction and a concise formula, we have
|Φn...(2M−N)]j+1 〉
αn,in
[n]in
(j+1)αn
λ[n]αn |in〉|φ
[(n+1)...(2M−N)]
[n−1]
|0〉|φ[(n+1)...(2M−N)]j+1 〉
+|1〉|φ[(n+1)...(2M−N)]j+2 〉
, (9)
where we denote
(j+1)αn
= δ(j+1)αn
n−1/(λ
[n−1]
n), (10)
(j+1)αn
= δ(j+2)αn
n−1/(λ
[n−1]
n ). (11)
Still we define that
V [n]inαnαn−1 = Γ
[n]in
αn−1αn
λ[n]αn . (12)
It is thus in the MPS representation. We can further con-
sider other cases including the ancilla state of the cloning
machine represented as Rj (Note it is not the ancilla state
in the MPS representation). We can find that the out-
put state of the general UQCM can be expressed as MPS
as in form (1). So it can be created sequentially. The
explicit results are summarized in the appendix.
We have shown that the output states of the general
UQCM in (4,5) are MPS’s and thus can be generated
sequentially. The sequential matrices V [n] of course de-
pend on the input |(N−m)0,m1〉 which are W-like states
and are generally multiqubit entangled. For later con-
venience, we denote V (m) to express that it depends
on input state for different m. By a straightforward
method, the sequential cloning operation, i.e., the iso-
metrices, depending on different input may take the form
m |(N − m)0,m〉〈(N − m)0,m1| ⊗ V (m). However,
this operation may need a single global unitary opera-
tor which involves N -qubit entangled states except for
m = 0,m = N . This contradicts with our aim that each
operation should be divided into sequential unitary oper-
ators in a quDit (quantum state in D-dimensional space)
times qubit system. Here we can use a scheme like the
following: the ancillary state interacts with each qubit
according to the (N + 1) × D-dimensional isometrices
CmN |0〉〈0|⊗N−m⊗|1〉〈1|⊗m⊗V (m) sequentially,
here a whole normalization factor is omitted. We know
that the operation |0〉〈0|⊗N−m ⊗ |1〉〈1|⊗m acts on each
qubit individually. Thus this scheme reduces the com-
plexity of the operation. This finishes our general se-
quential UQCM for the case of qubit. In case N = 1,
we recover the result of Ref.[23] for 1 to M cloning.
We should remark that similar as the case of sequen-
tial 1 to M UQCM in Ref.[23], for the general sequential
UQCM, the minimal dimension D of the ancillary state
grows linearly at most with M −N/2 + 1 for even N or
M − (N − 1)/2 for odd N .
Next we will consider a more general case that the se-
quential cloning machine is about the quantum state in d-
dimensional Hilbert space. We will use the d-dimensional
UQCM proposed by Fan et al in Ref. [25]. This UQCM
is a generalization of the cloning machine proposed in
Ref.[24] and we can use this UQCM to study its sequen-
tial form for d-dimensional case.
An arbitrary d-dimensional pure state takes the form
|Φ〉 =
i=0 xi|i〉 with
i=0 |xi|2 = 1. N identical pure
states can be expanded in terms of state in symmet-
ric subspace |Φ〉⊗N =
m1!...md!
xm10 ...x
d−1|~m〉,
where |~m〉 ≡ |m1, ...,md〉 is a symmetric state with mi
states of |i − 1〉, and also mi should satisfy a relation
i=1mi = N . The cloning transformations with states
in symmetric subspace can be written as
|~m〉 → |Φ~mM 〉 =
|~m+~j〉 ⊗ |~j〉, (13)
i=1 C
mi+ji
CM−NM+d−1
where ~j should satisfy
i ji = M − N . This cloning
machine is optimal and the corresponding fidelity of a
single quantum state between input and output is F =
(N(d+M) +M −N) /(d+N)M .
As for qubit system, we next show that the output
states for all symmetric states input can be expressed
as the sequential form. We consider the case 1 < n ≤
M − 1, and the state |Φ~mM 〉 is a bipartite state across
1...n : (n+ 1)... cut,
|Φ~mM 〉 =
|~j〉|φ[(n+1)...(M+1)]
〉 (15)
where
~m(~j−~m+~k)
i=1 C
ji+ki
, (16)
|φ[(n+1)...(M+1)]
~m(~j−~m+~k)
i=1 C
ji+ki
|~k〉|~j − ~m+ ~k〉/λ[n]
. (17)
By the same procedure as that of qubit case, we can
obtain the following
|φ[n...(M+1)]
[n]in
λ[n]αn |in〉|φ
(n+1)...(M+1)]
〉. (18)
Then we have
[n]in
= δαn(~j+~ein+1)
jin+1 + 1
[n−1]
. (19)
Still we can define V
[n]in
αnαn−1 = Γ
[n]in
αn−1αnλ
αn , and thus we
can find that each state |Φ~mM 〉 is a MPS and thus can be
sequentially generated. The detailed result of this part
will be presented elsewhere [27].
In conclusion, we show that the generalN toM univer-
sal quantum cloning machine can be implemented by a se-
quential manner. Since the sequential generation of mul-
tipartite state can be implemented in various physical se-
tups such as microwave and optical cavity QED, trapped
ions and quantum dots etc. This general sequential quan-
tum cloning machine may be implemented much easier
than the single global implementation scheme. This re-
duces dramatically the complexity in implementing the
general UQCM. We also show that for d-dimensional
quantum state, the sequential UQCM is also available.
Besides the universal cloning machine, the 1 toM phase-
covariant quantum cloning machine can also be sequen-
tially implemented. It will be interesting to consider sim-
ilarly the generalN toM phase-covariant cloning and the
economic phase-covariant cloning. The sequential asym-
metric quantum cloning machine may also be an inter-
esting topic.
Acknowledgements: HF was supported by ”Bairen”
program, NSFC and ”973” program (2006CB921107).
Appendix.–The explicit form of matrices V are pre-
sented as:
V [n]0αnαn−1 = δαnαn−1 ×
∑M−m−n
k=−m X
m+αn−1−1+k
∑M−m−n+1
k=−m X
M−n+1
m+αn−1−1+k
V [n]0αnαn−1 = δαnαn−1+1 ×
∑M−m−n
k=−m X
m+αn−1+k
∑M−m−n+1
k=−m X
M−n+1
m+αn−1−1+k
where notations X = β2m(αn−1−1+k), X
′ = β2m(αn−1+k)
are used. For case 1 < n ≤ M − N + m,αn−1 =
1, ..., n;αn = 1, ..., (n+1), and for caseM−N+m < n ≤
M − 1, αn−1, αn = 1, ..., (M −N +m+1). We can check
that the above defined V satisfies the isometry condition
V [n]in
V [n]in = 1. Similarly we have
V [M ]0αMαM−1 = δαMαM−1 ×
M−1−1−m)
M−1−1−m)
M−1−1
M−1−m)
V [M ]1αMαM−1 = δαM (αM−1+1) ×
M−1−m)
M−1−1−m)
M−1−1
M−1−m)
where 0 ≤ m ≤ N −m,αM−1, αM = 1, 2, ..., (M − N +
m+ 1).
For case concerning about ancilla state of the UQCM,
assume 1 ≤ l ≤M −N , we have
V [M+l]0αM+lαM+l−1 = δαM+l(αM+l−1−1) ×
αM+l−1 −m− 1
M −N − l + 1
V [M+l]1αM+lαM+l−1 = δαM+lαM+l−1 ×
M −N − l − αM+l−1 +m+ 1
M −N − l + 1
(1) For (m+ 1) ≤ αM+l ≤ (M −N +m− l+ 1),
(m+ 2) ≤ αM+l−1 ≤ (M −N +m− l+ 2),
[M+l]0
αM+lαM+l−1 = δαM+l(αM+l−1−1)
αM+l−1−m−1
M−N−l+1 .
For αM+l = (M −N +m− l + 2), 1 ≤ αM+l−1 ≤
(M −N +m+ 1), V [M+l]0αM+lαM+l−1 = 0. Otherwise
[M+l]0
αM+lαM+l−1 = δαM+lαM+l−1
(2) For (m+ 1) ≤ αM+l, αM+l−1 ≤
(M −N +m− l + 1), V [M+l]1αM+lαM+l−1 =
δαM+lαM+l−1
M−N−l−αM+l−1+m+2
M−N−l+1 . For αM+l =
(M −N +m− l + 2), 1 ≤ αM+l−1 ≤ (M −N +m+ 1),
[M+l]0
αM+lαM+l−1 = 0. Otherwise V
[M+l]0
αM+lαM+l−1 =
δαM+lαM+l−1
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|
0704.0324 | On the pseudospectrum of elliptic quadratic differential operators | ON THE PSEUDOSPECTRUM OF ELLIPTIC QUADRATIC
DIFFERENTIAL OPERATORS
Karel Pravda-Starov
University of California, Berkeley
Abstract. We study the pseudospectrum of a class of non-selfadjoint differential
operators. Our work consists in a detailed study of the microlocal properties, which
rule the spectral stability or instability phenomena appearing under small pertur-
bations for elliptic quadratic differential operators. The class of elliptic quadratic
differential operators stands for the class of operators defined in the Weyl quantiza-
tion by complex-valued elliptic quadratic symbols. We establish in this paper a simple
necessary and sufficient condition on the Weyl symbol of these operators, which en-
sures the stability of their spectra. When this condition is violated, we prove that
it occurs some strong spectral instabilities for the high energies of these operators,
in some regions which can be far away from their spectra. We give a precise geo-
metrical description of them, which explains the results obtained for these operators
in some numerical simulations giving the computation of “false eigenvalues” far from
their spectra by algorithms for eigenvalues computing.
Key words. Spectral instability, pseudospectrum, semiclassical quasimodes, non-
selfadjoint operators, non-normal operators, condition (Ψ), subellipticity.
2000 AMS Subject Classification. 35P05, 35S05.
1. Introduction
1.1. Miscellaneous facts about pseudospectrum. In recent years, there has been
a lot of interest in studying the pseudospectrum of non-selfadjoint operators. The
study of this notion has been initiated by noticing that for certain problems of sci-
ence and engineering involving non-selfadjoint operators, the predictions suggested by
spectral analysis do not match with the numerical simulations. This fact lets thinking
that in some cases the only knowledge of the spectrum of an operator is not enough to
understand sufficiently its action. To supplement this lack of information contained
in the spectrum, some new subsets of the complex plane called pseudospectra have
been defined. The main idea about the definition of these new subsets is that it is
interesting to study not only the points where the resolvent of an operator is not de-
fined, i.e. its spectrum, but also where this resolvent is large in norm. This explains
the following definition of the ε-pseudospectrum σε(A) of a matrix or an operator A,
σε(A) =
z ∈ C, ‖(zI −A)−1‖ ≥ 1
for any ε > 0, if we write by convention that ‖(zI − A)−1‖ = +∞ for every point z
belonging to the spectrum σ(A) of the operator.
Let us mention that there exists an abundant literature about this notion of pseu-
dospectrum. We refer here for the definition and some general properties of pseu-
dospectra to the paper [15] of L.N. Trefethen. Let us also point out the more recently
published book [16], which draws up a wide all-round view of this topic and gives a
lot of illustrations.
According to the previous definition, studying the pseudospectra of an operator is
exactly studying the level lines of the norm of its resolvent. What is interesting in
studying such level lines is that it gives some information about the spectral stability
http://arxiv.org/abs/0704.0324v1
of the operator. Indeed, pseudospectra can be defined in an equivalent way in term of
spectra of perturbations of the operator. For instance, we have for any A ∈ Mn(C),
σε(A) = {z ∈ C, z ∈ σ(A +B) for some B ∈ Mn(C) with ‖B‖ ≤ ε}.
It follows that a complex number z belongs to the ε-pseudospectrum of a matrix A if
and only if it belongs to the spectrum of one of its perturbations A+B with ‖B‖ ≤ ε.
More generally, if A is a closed unbounded linear operator with a dense domain on a
complex Hilbert space H , the result of Roch and Silbermann in [13] gives that
σε(A) =
B∈L(H), ‖B‖L(H)≤ε
σ(A +B),
where L(H) stands for the set of bounded linear operators on H . From this second
description, we understand the interest in studying such subsets if we want for example
to compute numerically some eigenvalues of an operator. Indeed, we start to do
it by discretizing this operator. This discretization and inevitable round-off errors
will generate some perturbations of the initial operator. Eventually, algorithms for
eigenvalues computing will determine the eigenvalues of a perturbation of the initial
operator, i.e. a value in a ε-pseudospectrum of the initial operator but not necessarily
a spectral one. This explains why it is important in such numerical computations to
understand if the ε-pseudospectra of studied operators contain more or less deeply
their spectra.
Let us first notice that this study is a priori non-trivial only for non-selfadjoint
operators, or more precisely for non-normal operators. Indeed, we have for a normal
operator A an exact expression of the norm of its resolvent given by the following
classical formula (see for example (V.3.31) in [8]),
(1.1.1) ∀z 6∈ σ(A), ‖(zI −A)−1‖ = 1
z, σ(A)
where d
z, σ(A)
stands for the distance between z and the spectrum of the operator,
when A is a closed unbounded linear operator with a dense domain on a complex
Hilbert space. This formula proves that the resolvent of a normal operator cannot
blow up far from its spectrum. It ensures the stability of its spectrum under small
perturbations because the ε-pseudospectrum is exactly equal in this case to the ε-
neighbourhood of the spectrum
(1.1.2) σε(A) =
z ∈ C : d
z, σ(A)
Nevertheless it is well-known that this formula (1.1.1) is no more true for non-normal
operators. For such operators, it can occur that their resolvents are very large in
norm far from their spectra. This induces that the spectra of these operators can be
very unstable under small perturbations. To illustrate this fact, let us consider the
case of the rotated harmonic oscillator and the following numerical computation of its
spectrum. The rotated harmonic oscillator is a simple example of elliptic quadratic
differential operator
Hc = D
x + cx
2, Dx = i
−1∂x,
with c = eiπ/4. The numerical computation is performed on the matrix discretization
(HcΨi,Ψj)L2(R)
1≤i,j≤N
where N is an integer taken equal to 100 and (Ψj)j∈N∗ stands for the basis of L
composed by Hermite functions. The black dots appearing on this computation stand
for the numerically computed eigenvalues. We can notice on this numerical simulation
that the computed low energies are very close to theoretical ones since the spectrum
Figure 1. Computation of some level lines of the norm of the resol-
vent ‖(Hc− z)−1‖ = ε−1 for the rotated harmonic oscillator Hc with
c = eiπ/4. The right column gives the corresponding values of log10 ε.
0 20 40 60 80 100 120 140 160
dim = 100
of the rotated harmonic oscillator is only composed of eigenvalues regularly spaced
out on the half-line eiπ/8R∗+,
σ(Hc) = {eiπ/8(2n+ 1) : n ∈ N}.
However we notice that it is no more true for the high energies. It occurs for them
some strong spectral instabilities, which lead to the computation of “false eigenvalues”
far from the half-line eiπ/8R∗+. Let us mention that some comparable computations
can be found in [3]. In this paper, we are interested in studying when and how this
kind of phenomena occurs in the class of elliptic quadratic differential operators.
1.2. Elliptic quadratic differential operators. We study here the class of elliptic
quadratic differential operators. It is the class of pseudodifferential operators defined
in the Weyl quantization
(1.2.1) q(x, ξ)wu(x) =
(2π)n
ei(x−y).ξq
(x+ y
u(y)dydξ,
by some symbols q(x, ξ), where (x, ξ) ∈ Rn×Rn and n ∈ N∗, which are some complex-
valued elliptic quadratic forms i.e. complex-valued quadratic forms verifying
(1.2.2) (x, ξ) ∈ Rn × Rn, q(x, ξ) = 0 ⇒ (x, ξ) = (0, 0).
Let us first notice that since the symbols of these operators are some quadratic forms,
these are only some differential operators, which are a priori non-selfadjoint because
their Weyl symbols are complex-valued. As mentioned before, the rotated harmonic
oscillator is an example of such an operator since we have
D2x + e
iθx2 = (ξ2 + eiθx2)w, 0 < θ < π,
if Dx = i
−1∂x. This operator is a very simple example of non-selfadjoint operator for
which we have noticed on the previous numerical simulation that it occurs some strong
spectral instabilities under small perturbations for its high energies. These phenomena
have been studied in several recent works. We can mention in particular the works of
L.S. Boulton [1], E.B. Davies [3], K. Pravda-Starov [10] and M. Zworski [18], which
have given a good understanding of these phenomena.
A question, which has been at the origin of this work, has been to study if these
phenomena peculiar to the rotated harmonic oscillator are representative, or not, of
what occurs more generally in the class of elliptic quadratic differential operators in
every dimension. We have tried to answer to the following questions:
- Does it always occur some strong spectral instabilities under small perturba-
tions for the high energies of these operators ?
- If it is not the case, is it possible to give a necessary and sufficient condition on
the Weyl symbols of these operators, which ensures their spectral stability ?
- Can we precisely describe the geometry, which separates the regions of the
resolvent sets where the resolvents of these operators blow up in norm from
the ones where one keeps a control on their sizes ?
To understand these spectral stability or instability phenomena, we need to study
the microlocal properties, which rule these phenomena in the class of elliptic quadratic
differential operators. Let us mention that it is M. Zworski who first underlined in [18]
the close link between these questions of spectral instabilities and some results of
microlocal analysis about the solvability of pseudodifferential operators.
1.3. Semiclassical pseudospectrum. To answer to these previous questions, it is
interesting to use a semiclassical setting and to study a notion of pseudospectrum
in this new setting. We define for a semiclassical family (Ph)0<h≤1 of operators on
L2(Rn), with a domain D, the following notions of semiclassical pseudospectra.
Definition 1.3.1. For all µ ≥ 0, the set
Λscµ (Ph) =
z ∈ C : ∀C > 0, ∀h0 > 0, ∃ 0 < h < h0, ‖(Ph − z)−1‖ ≥ Ch−µ
is called semiclassical pseudospectrum of index µ of the semiclassical family (Ph)0<h≤1.
The semiclassical pseudospectrum of infinite index is defined by
Λsc∞(Ph) =
Λscµ (Ph).
With this definition, the points in the complement of the semiclassical pseudospectrum
of index µ are the points of the complex plane where we have the following control of
the resolvent’s norm for sufficiently small values of the semiclassical parameter h,
(1.3.1) ∃C > 0, ∃h0 > 0, ∀ 0 < h < h0, ‖(Ph − z)−1‖ < Ch−µ.
To prove the existence of semiclassical pseudospectrum of index µ, we will study the
question of existence of semiclassical quasimodes
(1.3.2) ∀C > 0, ∀h0 > 0, ∃ 0 < h < h0, ∃uh ∈ D,
‖uh‖L2(Rn) = 1 and ‖Phuh − zuh‖L2(Rn) ≤ Chµ,
in some points z of the resolvent set, which can be considered as some “almost eigen-
values” in O(hµ) in the semiclassical limit. Let us notice that the definition chosen
here for the notions of semiclassical pseudospectra differ from the one given in [5] for
a semiclassical pseudodifferential operator. In fact, we have chosen a definition for
semiclassical pseudospectra inspired by the remark made p.388 in [5], because this
definition only depends on the properties of the semiclassical operator rather than on
its symbol.
The interest of working in a semiclassical setting is a matter of geometry. We can
explain this choice by the fact that it is easier for an elliptic quadratic differential oper-
ator q(x, ξ)w to describe the geometry of semiclassical pseudospectra of its associated
semiclassical operator (q(x, hξ)w)0<h≤1, than to describe directly the geometry of its
ε-pseudospectra. The semiclassical setting is particularly well-adapted for the study
of elliptic quadratic differential operators because there exists a simple link between
this semiclassical setting and the quantum one. Indeed, using that the symbols of
these operators are some quadratic forms q, we obtain from the change of variables,
y = h1/2x with h > 0, the following identity between the quantum operator q(x, ξ)w
and its associated semiclassical operator (q(x, hξ)w)0<h≤1,
(1.3.3) q(x, ξ)w − z
q(y, hη)w − z
if z ∈ C. This identity allows to get some information about the resolvent’s norm
behaviour of the quantum operator
q(x, ξ)w − z
if we have some information about semiclassical pseudospectra for its associated semi-
classical operator. Let us mention for example that if a non-zero complex number z
belongs to the semiclassical pseudospectrum of infinite index of the operator
(q(x, hξ)w)0<h≤1,
the identity (1.3.3) induces that the resolvent’s norm of the quantum operator blows
up along the half-line zR+ with a rate faster than any polynomials
(1.3.4) ∀N ∈ N, ∀C > 0, ∀η0 ≥ 1, ∃η ≥ η0, ‖
q(x, ξ)w − zη
)−1‖ ≥ CηN ,
and this, even if this half-line zR+ does not intersect the spectrum of the opera-
tor q(x, ξ)w. Conversely, in the case where z 6∈ Λscµ
q(y, hη)w
, z 6= 0 and 0 ≤ µ ≤ 1,
the identity (1.3.3) shows that we can find some positive constants C1 and C2 such
that the resolvent of the operator q(x, ξ)w remains bounded in norm in some regions
of the resolvent set of the shape
(1.3.5)
u ∈ C : |u| ≥ C1, d(∆, u) ≤ C2|proj∆u|1−µ
∩ C \ σ
q(x, ξ)w
where ∆ = zR+ and proj∆u stands for the orthogonal projection of u on the closed
half-line ∆. Indeed, we obtain from (1.3.1) and (1.3.3) that
∃C > 0, ∃η0 ≥ 1, ∀η ≥ η0,
q(x, ξ)w − ηeiargz
∥ < Cηµ−1,
which induces that for all v ∈ D
q(x, ξ)w
and η ≥ η0,
q(x, ξ)w − ηeiargz
L2(Rn)
≥ C−1η1−µ‖v‖L2(Rn),
q(x, ξ)w
stands for the domain of the operator q(x, ξ)w . Then, we can find a
constant η̃0 ≥ 1 such that if z̃ belongs to
u ∈ C : |u| ≥ η̃0, d(eiargzR+, u) ≤ 2−1C−1|projeiargzR+u|
∩ C\σ
q(x, ξ)w
|projeiargzR+ z̃| ≥ η0.
This induces using the previous estimates and the triangular inequality that if z̃
belongs to
u ∈ C : |u| ≥ η̃0, d(eiargzR+, u) ≤ 2−1C−1|projeiargzR+u|
∩ C\σ
q(x, ξ)w
we have for all v ∈ D
q(x, ξ)w
q(x, ξ)w − z̃
q(x, ξ)w − projeiargzR+ z̃
eiargzR+, z̃
‖v‖L2
≥ 2−1C−1|projeiargzR+ z̃|
1−µ‖v‖L2
≥ 2−1C−1η1−µ0 ‖v‖L2,
because µ ≤ 1. This last estimate shows that the resolvent of the operator q(x, ξ)w is
bounded in norm by 2Cη
0 on the set
u ∈ C : |u| ≥ η̃0, d(eiargzR+, u) ≤ 2−1C−1|projeiargzR+u|
∩ C\σ
q(x, ξ)w
We notice that depending directly on the value of the index µ, 0 ≤ µ < 1, the previous
set contains more or less deeply in its interior the half-line
{u ∈ C : |u| ≥ η̃0, u ∈ zR+}.
This fact explains why in the following we will precise carefully the index of the
semiclassical pseudospectrum to which a point does not belong when there is no
semiclassical pseudospectrum of infinite index in that point.
2. Statement of the results
2.1. Some notations and some preliminary facts about elliptic quadratic
differential operators. Let us begin by giving some notations and recalling known
results about elliptic quadratic differential operators. Let q be a complex-valued
elliptic quadratic form
q : Rnx × Rnξ → C
(x, ξ) 7→ q(x, ξ),
with n ∈ N∗, i.e. a complex-valued quadratic form verifying (1.2.2). The numerical
range Σ(q) of q is defined by the subset in the complex plane of all values taken by
this symbol
(2.1.1) Σ(q) = q(Rnx × Rnξ ),
and the Hamilton map F ∈ M2n(C) associated to the quadratic form q is uniquely
defined by the identity
(2.1.2) q
(x, ξ); (y, η)
(x, ξ), F (y, η)
, (x, ξ) ∈ R2n, (y, η) ∈ R2n,
where q
stands for the polar form associated to the quadratic form q and σ is the
symplectic form on R2n,
(2.1.3) σ
(x, ξ), (y, η)
= ξ.y − x.η, (x, ξ) ∈ R2n, (y, η) ∈ R2n.
Let us first notice that this Hamilton map F is skew-symmetric with respect to σ.
This is just a consequence of the properties of skew-symmetry of the symplectic form
and symmetry of the polar form
(2.1.4) ∀X,Y ∈ R2n, σ(X,FY ) = q(X ;Y ) = q(Y ;X) = σ(Y, FX) = −σ(FX, Y ).
Under this assumption of ellipticity, the numerical range of a quadratic form can
only take some very particular shapes. It is a consequence of the following result
proved by J. Sjöstrand (Lemma 3.1 in [14]),
Proposition 2.1.1. Let q : Rnx × Rnξ → C a complex-valued elliptic quadratic form.
If n ≥ 2, then there exists z ∈ C∗ such that Re(zq) is a positive definite quadratic
form. If n = 1, the same result is fulfilled if we assume besides that Σ(q) 6= C.
This proposition shows that the numerical range of an elliptic quadratic form can only
take two shapes. The first possible shape is when Σ(q) is equal to the whole complex
plane. This case can only occur in dimension n = 1. The second possible shape is
when Σ(q) is equal to a closed angular sector with a top in 0 and an opening strictly
lower than π.
Figure 2. Shape of the numerical range Σ(q) when Σ(q) 6= C.
Σ(zq)
Indeed, if Σ(q) 6= C, using that the set Σ(q) is a semi-cone
tq(x, ξ) = q(
tξ), t ∈ R+, (x, ξ) ∈ R2n,
because q is a quadratic form, we have
Σ(q) = R+z
if z is the non-zero complex number given by the proposition 2.1.1 and I is the compact
interval
I = 1 + i Im(zq)(K),
where K is the following compact subset of R2n,
(x, ξ) ∈ R2n : Re(zq)(x, ξ) = 1
The compactness of K is a direct consequence of the fact that Re(zq) is a positive
definite quadratic form.
Elliptic quadratic differential operators define some Fredholm operators (see Lemma 3.1
in [6] or Theorem 3.5 in [14]),
(2.1.5) q(x, ξ)w + z : B → L2(Rn),
where B is the Hilbert space
(2.1.6)
u ∈ L2(Rn) : xαDβxu ∈ L2(Rn) if |α+ β| ≤ 2
with the norm
‖u‖2B =
|α+β|≤2
‖xαDβxu‖2L2(Rn).
The Fredholm index of the operator q(x, ξ)w + z is independent of z and is equal to 0
if n ≥ 2. In the case where n = 1, this index can take the values −2, 0 or 2. More
precisely, this index is always equal to 0 if Σ(q) 6= C.
In the following, we will always assume that Σ(q) 6= C. Under this assumption,
J. Sjöstrand has proved in the theorem 3.5 in [14] (see also Lemma 3.2 and Theorem 3.3
in [6]) that the spectrum of an elliptic quadratic differential operator
q(x, ξ)w : B → L2(Rn),
is only composed of eigenvalues with finite multiplicity
(2.1.7) σ
q(x, ξ)w
λ∈σ(F ),
−iλ∈Σ(q)\{0}
rλ + 2kλ
(−iλ) : kλ ∈ N
where F is the Hamilton map associated to the quadratic form q and rλ is the dimen-
sion of the space of generalized eigenvectors of F in C2n belonging to the eigenvalue
λ ∈ C. Let us notice that the spectra of these operators is always included in the
numerical range of their Weyl symbols.
To end this review of preliminary properties of elliptic quadratic differential oper-
ators, let us underline that the property of normality in this class of operators can be
easily checked by computing the Poisson bracket of the real part and the imaginary
part of their symbols
(2.1.8) {Re q, Im q} =
∂Re q
∂Im q
∂Re q
∂Im q
Proposition 2.1.2. An elliptic quadratic differential operator
q(x, ξ)w : B → L2(Rn), n ∈ N∗,
is normal if and only if the quadratic form defined by the Poisson bracket of the real
part and the imaginary part of its symbol is equal to zero
(2.1.9) ∀(x, ξ) ∈ R2n, {Re q, Im q}(x, ξ) = 0.
Proof of Proposition 2.1.2. This proposition is a direct consequence of the composition
formula in Weyl calculus (see Theorem 18.5.4 in [7]), which induces that the Weyl
symbol of the commutator
[qw, (qw)∗] = [qw, qw] = −2i[(Re q)w, (Im q)w],
is equal to
−2i(Re q ♯ Im q − Im q ♯ Re q) = −2{Re q, Im q},
because Re q and Im q are some quadratic forms. The notation Re q ♯ Im q stands
for the Weyl symbol of the operator obtained by composition (Req)w(Imq)w. �
Remark. Let us notice that the symplectic invariance of the Poisson bracket (see
(21.1.4) in [7]),
(2.1.10) {(Re q) ◦ χ, (Im q) ◦ χ} = {Re q, Im q} ◦ χ,
if χ stands for a linear symplectic transformation of R2n, implies that the condition
(2.1.9) is symplectically invariant.
2.2. Statement of the main results. Let us consider an elliptic quadratic differ-
ential operator
q(x, ξ)w : B → L2(Rn).
We know from (2.1.7) that the spectrum of this operator is contained in the numerical
range of its symbol Σ(q). The following proposition gives a first localization of the
regions where the resolvent can blow up in norm and where spectral instabilities can
occur.
Proposition 2.2.1. Let q : Rn × Rn → C, n ∈ N∗, be a complex-valued elliptic
quadratic form. We have
∀z 6∈ Σ(q),
q(x, ξ)w − z
∥ ≤ 1
z,Σ(q)
where d
z,Σ(q)
stands for the distance from z to the numerical range Σ(q).
This result shows that the resolvent of an elliptic quadratic differential operator
cannot blow up in norm far from the numerical range of its symbol. We are now
going to study what kind of phenomena can occur in this particular set. There are
two cases to separate according to the property of normality or non-normality of the
operator.
2.2.1. Case of a normal operator. Let us consider a normal elliptic quadratic differ-
ential operator
q(x, ξ)w : B → L2(Rn).
Let us recall that according to the proposition 2.1.2 this property of normality is
exactly equivalent to the fact that
∀(x, ξ) ∈ R2n, {Re q, Im q}(x, ξ) = 0.
In this case, we have the classical formula (1.1.1) for its resolvent’s norm
(2.2.1) ∀z 6∈ σ
q(x, ξ)w
q(x, ξ)w − z
z, σ(q(x, ξ)w)
which induces that the ε-pseudospectrum of this operator is exactly equal to the
ε-neighbourhood of its spectrum
q(x, ξ)w
z ∈ C : d
z, σ(q(x, ξ)w)
, ε > 0.
This classical formula (2.2.1) ensures that the resolvent cannot blow up in norm far
from the spectrum and induces that the spectrum of such an operator is stable under
small perturbations.
Example 1. The operator
(2.2.2) q1(x, ξ)
w = −(1 + i)∂2x1 − ∂
+ 4(−1 + i)x1∂x1 + 2(−1 + i)x2∂x1 + 6ix2∂x2
+ 2ix1∂x2 + (6 + 5i)x
1 + (11 + i)x
2 + (10 + 4i)x1x2 − 2 + 5i,
is an example of a normal elliptic quadratic differential operator. Its spectrum is given
q1(x, ξ)
(2k1 + 1) + (2k2 + 1)
4 : (k1, k2) ∈ N2
Figure 3. Spectrum and a ε-pseudospectrum of the operator q1(x, ξ)
∗ ∗ ∗ ∗ ∗ ∗ ∗
∗ ∗ ∗ ∗ ∗ ∗
∗ ∗ ∗ ∗ ∗
∗ ∗ ∗ ∗
∗ ∗ ∗
Σ(q1)
Example 2. Let us notice that when the numerical range Σ(q) is reduced to a closed
half-line, the elliptic quadratic differential operator q(x, ξ)w is always normal since
{Re q, Im q} = |z|2{Re(z−1q), Im(z−1q)} = 0,
if z ∈ C∗ is chosen such that Im(z−1q) = 0. In fact, the operator q(x, ξ)w can in this
particular case be reduced after a conjugation by a unitary operator on L2(Rn) to the
operator
+ x2j),
where λj > 0 for all j = 1, ..., n.
Figure 4. Example of a normal elliptic quadratic differential operator.
2.2.2. Case of a non-normal operator. Let us consider a non-normal elliptic quadratic
differential operator
q(x, ξ)w : B → L2(Rn), n ∈ N∗.
We assume in the following that the numerical range Σ(q) is distinct from the whole
complex plane
(2.2.3) Σ(q) 6= C.
As mentioned in the section 2.1, this additional assumption is always fulfilled in
dimension n ≥ 2. It only excludes some very particular one-dimensional elliptic
quadratic differential operators (see the remark following the proposition 2.2.2 for
more precision about these operators).
Under this additional assumption, the numerical range Σ(q) is always a closed
angular sector with a top in 0 and a positive opening strictly lower than π.
2.2.2.a. On the pseudospectrum at the interior of the numerical range. Let us consider
the associated semiclassical elliptic quadratic differential operator
(q(x, hξ)w)0<h≤1.
We can build in every point of the interior of the numerical range Σ̊(q) some semi-
classical quasimodes.
Theorem 2.2.1. If the elliptic quadratic differential operator
q(x, ξ)w : B → L2(Rn), n ∈ N∗,
is non-normal and verifies Σ(q) 6= C then for all z ∈ Σ̊(q) and N ∈ N, there exist
h0 > 0 and a semiclassical family (uh)0<h≤h0 ∈ S(Rn) such that
‖uh‖L2(Rn) = 1 and ‖q(x, hξ)wuh − zuh‖L2(Rn) = O(hN ) when h → 0+.
This result induces the existence of semiclassical pseudospectrum of infinite index in
every point of the interior of the numerical range Σ̊(q).
According to (1.3.4), this result in the semiclassical setting induces that the resol-
vent’s norm of the quantum operator q(x, ξ)w blows up fastly along all the half-lines
belonging to the interior of the numerical range Σ̊(q),
(2.2.4) ∀z ∈ Σ̊(q), ∀N ∈ N, ∀C > 0, ∀η0 ≥ 1, ∃η ≥ η0, ‖
q(x, ξ)w − zη
)−1‖ ≥ CηN .
We deduce from (2.1.7) that as soon as an elliptic quadratic differential operator is
non-normal its resolvent blows up in norm in some regions of the resolvent set far
from its spectrum. This fact induces that the high energies of such an operator are
very unstable under small perturbations as we have already noticed on the numerical
computation performed for the rotated harmonic oscillator. It follows that in the class
of elliptic quadratic differential operators1 the property of spectral stability is exactly
equivalent to the property of normality:
σ(q(x, ξ)w) is stable under ⇔ q(x, ξ)w is a normal ⇔ {Re q, Im q} = 0.
small perturbations operator
By spectral stability, we mean here that the resolvent of these operators cannot blow
up in norm far from their spectra. Let us add that it is not very surprising to have
this property of spectral stability under the assumption of normality, but it is worth
1If we exclude the one-dimensional particular cases previously mentioned.
noticing that as soon as this property is violated, it occurs in this class of operators
some strong spectral instabilities under small perturbations for their high energies.
Examples. The two following operators
(2.2.5) q2(x, ξ)
w = −∂2x1 − 2∂
+ 4ix2∂x2 + 2x
1 + (4 + i)x
2 + 4x1x2 + 2i
(2.2.6) q3(x, ξ)
w = −(1 + i)∂2x1 − 2∂
+ 4(−1 + i)x1∂x1 + 2(1− i)x2∂x1 − 4ix1∂x2
+ (9 + 4i)x21 + (2 + i)x
2 − 4(1 + i)x1x2 − 2 + 2i,
are some examples of non-normal elliptic quadratic differential operators.
2.2.2.b. On the pseudospectrum at the boundary of the numerical range. Let us now
study what occurs on the boundary of the numerical range ∂Σ(q) for a non-normal
elliptic quadratic differential operator
q(x, ξ)w : B → L2(Rn).
Let us mention that we always assume that Σ(q) 6= C. Under these assumptions, the
boundary of the numerical range is composed of the union of the origin 0 and two
half-lines ∆1 and ∆2,
(2.2.7) ∂Σ(q) = {0} ⊔∆1 ⊔∆2,
that we can write
(2.2.8) ∆1 = z1R
+ and ∆2 = z2R
+ with z1, z2 ∈ ∂Σ(q) \ {0}.
We need to define a notion of order for the symbol q(x, ξ) on these two half-lines ∆j ,
j = 1, 2. Let us begin by recalling the classical definition of the order k(x0, ξ0) of a
symbol p(x, ξ) at a point (x0, ξ0) ∈ R2n (see section 27.2, chapter 27 in [7]). This
order k(x0, ξ0) is an element of the set N ∪ {+∞} defined by
(2.2.9) k(x0, ξ0) = sup
j ∈ Z : pI(x0, ξ0) = 0, ∀ 1 ≤ |I| ≤ j
where I = (i1, i2, ..., ik) ∈ {1, 2}k, |I| = k and pI stands for the iterated Poisson
brackets
pI = Hpi1Hpi2 ...Hpik−1 pik ,
where p1 and p2 are respectively the real and the imaginary part of the symbol p,
p = p1 + ip2. The order of a symbol q at a point z is then defined as the maximal
order of the symbol p = q − z at every point (x0, ξ0) ∈ R2n verifying
p(x0, ξ0) = q(x0, ξ0)− z = 0.
Let us underline that the symplectic invariance of the Poisson bracket (2.1.10) induces
the same property for the order of a symbol at a point.
Since here the symbol q is a quadratic form, all the iterated Poisson brackets are
also some quadratic forms. This property of degree two homogeneity of these Poisson
brackets induces that the symbol q has the same order at every point of each half-line
∆j , j = 1, 2. This allows to define the order of the symbol q on the half-line ∆j by
defining this order by this common value. Let us mention that this order can be finite
or infinite.
Examples. One can easily check that the Weyl symbol
ξ2 + eiθx2, 0 < θ < π,
of the rotated harmonic oscillator has an order equal to 2 on the both half-lines R∗+
and eiθR∗+, which composes the boundary of its numerical range. The symbol q2 of
the operator defined in (2.2.5) has an order equal to 2 on iR∗+ and to 6 on R
Σ(q2) = {z ∈ C : Re z ≥ 0, Im z ≥ 0}.
On the other hand, we can verify that the symbol q3 of the operator defined in (2.2.6)
is of infinite order on the half-line R∗+ and has an order equal to 2 on e
iπ/4R∗+,
Σ(q3) = {0} ∪ {z ∈ C∗ : 0 ≤ arg z ≤ π/4}.
In the case where the symbol is of finite order on a half-line ∆j , j = 1, 2, we have
the following result.
Theorem 2.2.2. If the Weyl symbol q(x, ξ) of a non-normal elliptic quadratic differ-
ential operator is of finite order kj on the half-line
∆j , j ∈ {1, 2}, ∆j ⊂ ∂Σ(q) \ {0},
then this order is necessary even and there is no semiclassical pseudospectrum of
index kj/(kj + 1) on ∆j for the associated semiclassical operator
∆j ⊂ C \ Λsckj/(kj+1)
q(x, hξ)w
Remark. Let us mention that we can more precisely establish that in dimension n ≥ 1,
the order kj is an even integer verifying
2 ≤ kj ≤ 4n− 2.
This result is proved in [12].
By rephrasing this result in a quantum setting, it follows from (1.3.5) and (2.1.7)
that when the symbol q of a non-normal elliptic quadratic differential operator q(x, ξ)w
is of finite order kj on a half-line
∆j , j ∈ {1, 2}, ∆j ⊂ ∂Σ(q) \ {0},
then the resolvent of this operator remains bounded in norm in a set of the following
(2.2.10)
u ∈ C : |u| ≥ C1, d(∆j , u) ≤ C2|proj∆ju|
where C1 and C2 are some positive constants.
As we will see in its proof, this absence of semiclassical pseudospectrum is linked
to some properties of subellipticity. Let us just underline for the moment that the
index kj/(kj + 1), which appears in this result is exactly equal to the loss appearing
in the subelliptic estimate hidden behind this result.
About the case of infinite order, the situation is much more complicated. Never-
theless, we can first notice in this case that we cannot expect to prove a stronger result
than an absence of semiclassical pseudospectrum of index 1. Indeed, we can easily
check on the example of the operator q3(x, ξ)
w defined in (2.2.6) that its spectrum is
given by
q3(x, ξ)
(2k1 + 1)
2 + (2k2 + 1)3
8 : (k1, k2) ∈ N2
We recall that the spectrum of this operator is only composed of eigenvalues and that
its symbol is of infinite order on R∗+. It follows from the structure of the spectrum and
(1.3.5) that if there is no semiclassical pseudospectrum of infinite index in a point of
the half-line R∗+, there is necessary no semiclassical pseudospectrum of index µ with
an index µ ≥ 1. In fact, we can prove by using a result of exponential decay in time
for the norm of contraction semigroups generated by elliptic quadratic differential
operators (see [12]) that there is never some semiclassical pseudospectrum of index 1
on all these half-lines of infinite order. Let us mention that this result of exponential
decay will not be proved here but it will be explained in the following how it induces
the absence of semiclassical pseudospectrum of index 1.
2.2.3. About the geometry of ε-pseudospectra for elliptic quadratic differential opera-
tors. Let us now explain what are the consequences of these results on the geometry
of ε-pseudospectra for elliptic quadratic differential operators. Let us begin by con-
sidering the one-dimensional case which is a bit particular. In dimension n = 1, an
elliptic quadratic differential operator can be reduced after a similitude and a conju-
gation by a unitary operator to the harmonic oscillator or to the rotated harmonic
oscillator.
Proposition 2.2.2. Let us consider q : R×R → C a complex-valued elliptic quadratic
form such that Σ(q) 6= C. For all h > 0, there exist a unitary operator (more precisely
a metaplectic operator) Uh on L
2(R), which is an automorphism of the spaces S(R)
and B, z ∈ C∗ and θ ∈ [0, π[ such that
∀h > 0, q(x, hξ)w = zUh
(hDx)
2 + eiθx2
U−1h .
Remark. In the case where Σ(q) = C, an elliptic quadratic differential operator
q(x, ξ)w can be reduced after a similitude and a conjugation by a unitary operator on
L2(Rn) to the operator defined in the Weyl quantization by the symbol
(ξ + ix)(ξ + ηx) with η ∈ C, Im η > 0,
(ξ − ix)(ξ + ηx) with η ∈ C, Im η < 0,
depending on the value of its Fredholm index, which is equal to −2 in the first case
and to 2 in the second one.
As we will see in the following, this proposition allows us to reduce the study of a
one-dimensional non-normal elliptic quadratic differential operator verifying
Σ(q) 6= C,
to the one of the rotated harmonic oscillator
Hθ = D
x + e
iθx2, 0 < θ < π.
Let us mention that the previous results (Theorem 2.2.1 and Theorem 2.2.2) were
already known in the particular case of the rotated harmonic oscillator. Indeed,
the existence of semiclassical quasimodes inducing the presence of semiclassical pseu-
dospectrum of infinite index in every point of the interior of the numerical range for
the associated semiclassical operator, is a direct consequence of a result proved by
E.B. Davies in [4] (Theorem 1). About the absence of semiclassical pseudospectrum
of index 2/3 on the boundary of the numerical range, this result has been proved for
the rotated harmonic oscillator in [10]2.
As proved in [10], this absence of semiclassical pseudospectrum allows to give a
proof of a conjecture stated by L.S. Boulton in [1]. It deals with the geometry of ε-
pseudospectra for the rotated harmonic oscillator. Let us now recall some facts about
this conjecture and some results proved by L.S. Boulton in [1].
2Let us recall that the value of the order is equal to 2 in this case.
L.S. Boulton has first proved (Theorem 3.3 in [1]) that the resolvent of the rotated
harmonic oscillator blows up in norm along all a family of curves of the following form
η 7→ bη + eiθηp,
where b and p are some positive constants verifying 1/3 < p < 3,
(2.2.11)
Hθ − (bη + eiθηp)
∥ → +∞ when η → +∞.
On the other hand, he also proved that the resolvent of this operator remains bounded
in norm on two half-stripes parallel to the half-lines R+ or e
iθR+. More precisely, he
proved that there exist some positive constants d and Md such that
(2.2.12) sup
, 0≤b≤d
Hθ − (η + ib)
∥ ≤ Md,
(2.2.13) sup
, 0≤b≤d
Hθ − eiθ(η − ib)
∥ ≤ Md.
These bounds provide some information about the shape of ε-pseudospectra of the
operator Hθ. Indeed, L.S. Boulton has proved using these results that for all suffi-
ciently small value of the positive parameter ε, the ε-pseudospectra of the rotated
harmonic oscillator is contained in the shaded set appearing on the following figure.
The eigenvalues appear on this figure marked by some ⋄.
Figure 5. A first localization of the ε-pseudospectra of the rotated
harmonic oscillator.
More precisely, L.S. Boulton proved that for all 0 < δ < 1 and m ∈ N, there exists
a positive constant ε0 such that for all 0 < ε < ε0,
(2.2.14) σε(Hθ) ⊂
{z ∈ C : |z − λn| < δ} ∪
λm+1 − δeiθ/2 + Sθ
where
λn = e
iθ/2(2n+ 1), n ∈ N
Sθ = {z ∈ C∗ : 0 ≤ arg z ≤ θ} ∪ {0}.
In fact, in view of some numerical calculations performed by E.B. Davies in [3],
L.S. Boulton has conjectured that the index p = 1/3 appearing in (2.2.11) is the
critical one in the following sense:
Let us consider 0 < p < 1/3, 0 < δ < 1 and m ∈ N. If bm,p and E are some positive
constants verifying
bm,pE + e
iθEp = λm and ∀η > E, arg zη < θ/2,
where zη = bm,pη + e
iθηp, let us set
Ωm,p =
|zη|eiα ∈ C : η ≥ E, arg zη ≤ α ≤ arg(zηeiθ)
L.S. Boulton has conjectured the following result.
Boulton’s conjecture. There exists ε0 > 0 such that for all 0 < ε < ε0,
(2.2.15) σε(Hθ) ⊂
{z ∈ C : |z − λn| < δ} ∪ Ωm,p.
The absence of semiclassical pseudospectrum of index 2/3 on the boundary of the
numerical range ∂Σ(q)\{0} for the rotated harmonic oscillator3 given by the theorem
2.2.2 shows that this index 1/3 is actually the critical one. Indeed, we can deduce
(2.2.15) from (2.2.10) (see [10] for more details) since here kj = 2, j ∈ {1, 2}. As
we will see, this theorem 2.2.2 is a consequence of a subelliptic estimate for gen-
eral semiclassical pseudodifferential operators proved by N. Dencker, J. Sjöstrand and
M. Zworski in [5] (Theorem 1.4). In the particular case of the rotated harmonic oscil-
lator, a more elementary proof of this result using only some non-trivial localization
scheme in the frequency variable is given in [10].
Let us notice that this inclusion (2.2.15) allows to give a sharp description of the ε-
pseudospectra of the rotated harmonic oscillator, which is optimal in view of (2.2.11).
Figure 6. Shape of the ε-pseudospectra of the rotated harmonic oscillator.
By coming back to the case of an arbitrary dimension n ≥ 1, let us finally underline
that using the theorem 2.2.2, we can give similar descriptions of the ε-pseudospectra
for non-normal elliptic quadratic differential operators, to the one given by L.S. Boul-
ton for the rotated harmonic oscillator, when the symbols of these operators are of
finite order on the two open half-lines, which compose the boundary of their numerical
ranges. The only difference with the particular case of the rotated harmonic oscillator
is that the critical indices, which appear in this description can be different. Indeed,
3The order of the rotated harmonic oscillator’s symbol is equal to 2 on ∂Σ(q) \ {0}.
these critical indices depend directly according to (2.2.10) on the order of the symbols
on the two half-lines composing the boundary of their numerical ranges. We refer the
reader to [10] for more details about the way of getting from (2.2.10) such descriptions
of ε-pseudospectra.
3. The proofs of the results
Before giving the proofs of the results stated in the previous section, let us begin by
recalling the symplectic invariance property of the Weyl quantization (see Theorem
18.5.9 in [7]). This symplectic invariance is actually the most important property of
the Weyl quantization.
For every affine symplectic transformation χ of R2n, there exists a unitary trans-
formation U on L2(Rn), uniquely determined apart from a constant factor of modulus
1, such that U is an automorphism of the spaces S(Rn), B and S ′(Rn), where B is
the Hilbert space defined in (2.1.6), and
(3.0.1) (a ◦ χ)(x, ξ)w = U−1a(x, ξ)wU,
for all a ∈ S ′(R2n). The operator U is a metaplectic operator associated to the affine
symplectic transformation χ.
This symplectic invariance of the Weyl quantization induces the same property for
the semiclassical pseudospectra of elliptic quadratic differential operators in the sense
that if
q : Rnx × Rnξ → C,
is a complex-valued elliptic quadratic form and χ is a linear symplectic transformation
of R2n, we have for all µ ∈ [0,∞],
(3.0.2) Λscµ
(q ◦ χ)(x, hξ)w
= Λscµ
q(x, hξ)w
To prove this fact, let us begin by noticing that for all a ∈ S ′(R2n) and h > 0, we
U−1h a(x, ξ)
wUh = a(h
−1/2x, h1/2ξ)w,
where
Uhf(x) = h
n/4f(h1/2x),
since according to the proof of Theorem 18.5.9 in [7], Uh is a metaplectic operator
associated to the linear symplectic transformation
(x, ξ) 7→ (h−1/2x, h1/2ξ).
Let us now consider the case where the symbol a is a quadratic form. The homogeneity
property of such a symbol implies that
∀h > 0, a(h−1/2x, h1/2ξ) = 1
a(x, hξ),
∀h > 0, U−1h a(x, ξ)
wUh =
a(x, hξ)w.
If q : Rnx × Rnξ → C is a complex-valued elliptic quadratic form and χ is a linear
symplectic transformation of R2n, we can notice that
(q ◦ χ)(x, hξ)w , h > 0,
is actually an elliptic quadratic differential operator since the symbol q◦χ is an elliptic
quadratic form. Let z ∈ C and U be a metaplectic operator associated to the linear
symplectic transformation χ. Using that U and Uh are some automorphisms of the
Hilbert space B and
(3.0.3) U−1h U
−1Uhq(x, hξ)
wU−1h UUh = U
−1hq(x, ξ)wUUh
= hU−1h (q ◦ χ)(x, ξ)
wUh = (q ◦ χ)(x, hξ)w ,
we obtain that
U−1h U
q(x, hξ)w − z
U−1h UUh =
(q ◦ χ)(x, hξ)w − z
Using finally that U−1h U
−1Uh is a unitary transformation of L
2(Rn), this identity
implies that
(q ◦ χ)(x, hξ)w − z
q(x, hξ)w − z
which proves (3.0.2). In the following, this property of symplectic invariance will
allow us to reduce certain symbols to some normal forms by choosing new symplectic
coordinates. We can now begin to prove the results stated in the previous section.
Let us start by the proof of the proposition 2.2.1.
Proof of Proposition 2.2.1. If the numerical range is equal to the whole complex plane,
there is nothing to prove. If Σ(q) 6= C, we have seen in the previous section that the
numerical range is necessary a closed angular sector with a top in 0 and an opening
strictly lower than π.
Let us consider z 6∈ Σ(q) and denote by z0 its orthogonal projection on the non-
empty closed convex set Σ(q). According to the shape of the numerical range, it
follows that z0 belongs to its boundary and that we can find a complex number
z1 ∈ C∗, |z1| = 1 such that
Σ(z1q) ⊂
z ∈ C : Re z ≥ 0
(3.0.4) z1z ∈
z ∈ C : Re z < 0
z,Σ(q)
= d(z1z, iR).
Using now that the operator i[Im(z1q)]
w is formally skew-selfadjoint, we obtain that
for all u ∈ S(Rn),
z1q(x, ξ)
wu− z1zu, u
L2(Rn)
= d(z1z, iR)‖u‖2L2(Rn) +
z1q(x, ξ)
L2(Rn)
.(3.0.5)
Then, since the quadratic form Re(z1q) is non-negative, we deduce from the symplectic
invariance of the Weyl quantization and the theorem 21.5.3 in [7] that there exists a
metaplectic operator U such that
z1q(x, ξ)
= U−1
+ x2j) +
j=k+1
with k, l ∈ N and λj > 0 for all j = 1, ..., k. By using that U is a unitary operator on
L2(Rn), we obtain that the quantity
z1q(x, ξ)
L2(Rn)
‖DxjUu‖2L2(Rn) + ‖xjUu‖
L2(Rn)
j=k+1
‖xjUu‖2L2(Rn),
is non-negative. Then, we can deduce from the Cauchy-Schwarz inequality, (3.0.4)
and (3.0.5) that for all u ∈ S(Rn),
z,Σ(q)
‖u‖L2(Rn) ≤ |z1| ‖q(x, ξ)wu− zu‖L2(Rn).
Finally, using the density of the Schwartz space S(Rn) in B and the fact that |z1| = 1,
we obtain that
∀z 6∈ Σ(q),
q(x, ξ)w − z
∥ ≤ 1
z,Σ(q)
since according to (2.1.7), σ
q(x, ξ)w
⊂ Σ(q). �
We now consider the one-dimensional case, which is a bit particular.
3.1. The one-dimensional case. In dimension n = 1, we can reduce the study of
complex-valued elliptic quadratic forms to exactly three normal forms after a simili-
tude and a real linear symplectic transformation.
Lemma 3.1.1. Let q : Rx × Rξ → C be a complex-valued elliptic quadratic form in
dimension 1. Then, there exists a linear symplectic transformation χ of R2 such that
the symbol q ◦ χ is equal to one of the following normal forms:
(i) α(ξ2 + eiθx2) with α ∈ C∗, 0 ≤ θ < π.
(ii) α(ξ + ix)(ξ + ηx) with α ∈ C∗, η ∈ C, Im η > 0.
(iii) α(ξ − ix)(ξ + ηx) with α ∈ C∗, η ∈ C, Im η < 0.
In the two last cases (ii) and (iii), the numerical range Σ(q) is equal to the whole
complex plane, Σ(q) = C.
Proof of Lemma 3.1.1. Let q : R2 → C be a complex-valued elliptic quadratic form.
Let us first consider the case where Σ(q) 6= C. We deduce from the proposition 2.1.1
that we can reduce our study to the case where Re q is a positive definite quadratic
form. Then, using Lemma 18.6.4 in [7], we can find a real linear symplectic transfor-
mation to reduce the quadratic form Re q to the normal form
λ(x2 + ξ2), with λ > 0.
It follows that there exist some real constants a, b and c such that
q(x, ξ) = λ
x2 + ξ2 + i(ax2 + 2bxξ + cξ2)
Then, we can choose an orthogonal matrix P ∈ O(2,R) diagonalizing the real sym-
metric matrix associated to the quadratic form ax2 + 2bxξ + cξ2,
with λ1, λ2 ∈ R. If P ∈ O(2,R) \ SO(2,R), we have
if σ0 is the matrix with determinant equal to −1,
and P̃ = Pσ0. It follows that we can always diagonalize the real symmetric matrix
associated to the quadratic form λ−1Im q by conjugating it by an element of SO(2,R).
Since the symplectic group is equal in dimension 1 to the group SL(2,R), we can after
a linear symplectic transformation of R2 reduce the quadratic form q to
x2 + ξ2 + i(γ1x
2 + γ2ξ
= α(ξ2 + reiθx2),
where γ1, γ2 ∈ R, α ∈ C∗, r > 0 and θ ∈] − π, π[. Let us notice that the elliptic-
ity of q actually implies that θ 6≡ π[2π]. Finally, using the real linear symplectic
transformation (x, ξ) 7→ (r−1/4x, r1/4ξ), we get a symbol of type (i),
αr1/2(ξ2 + eiθx2),
if 0 ≤ θ < π. If −π < θ < 0, we need to use besides the real linear symplectic
transformation (x, ξ) 7→ (ξ,−x) to obtain a symbol of type (i),
2 eiθ(ξ2 + e−iθx2).
Let us now assume that Σ(q) = C. Since the dimension is equal to 1, we can factor the
symbol q on C as a polynomial function of degree 2 in the variable ξ. Thus, according
to the dependence in the variable x of the polynomial function’s coefficients, we can
find some complex numbers λ1, λ2 and α ∈ C∗ such that
q(x, ξ) = α(ξ − λ1x)(ξ − λ2x).
The ellipticity assumption for the quadratic form q induces that
Im λj 6= 0,
if j = 1, 2. Using now the linear symplectic transformation (x, ξ) 7→ (x, ξ + Re λ1x),
we can assume that
(3.1.1) q(x, ξ) = α(ξ − irx)(ξ + bx),
with r ∈ R∗ and Im b 6= 0. Let us now check that the assumption Σ(q) = C induces
that r Im b < 0. Since
(ξ − irx)(ξ + bx) = ξ2 + (b − ir)xξ − irbx2,
the condition Σ(q) = C implies that for all (v, w) ∈ R2, there exists a solution
(x0, ξ0) ∈ R2 of the system
(3.1.2)
ξ2 +Re b xξ + r Im b x2 = v
xξ(Im b− r) − r Re b x2 = w.
Let us first notice that the second equation of (3.1.2) is fulfilled for all w ∈ R only if
Im b 6= r.
If w 6= 0, it follows from the second equation of (3.1.2) that x0 6= 0 and
(3.1.3) ξ0 =
w + r Re b x20
(Im b− r)x0
Let us consider the case where v = 0. Using (3.1.3) and the first equation of (3.1.2),
we obtain that
(w + r Re b x20)
2 +Re b (Im b− r)x20(w + r Re b x20) + r Im b (Im b− r)2x40 = 0.
We can rewrite this equation as fw(X0) = 0 if we set X0 = x
0 and
(3.1.4) fw(X) = r Im b
(Re b)2 + (Im b− r)2
X2 + w Re b (Im b+ r)X + w2.
Thus, the condition Σ(q) = C implies that there exists for all w 6= 0, a non-negative
solution X0 of the equation fw(X0) = 0. Since the quantity r Im b is assumed to be
non-zero, we first study the case where r Im b > 0. In this case, since
(3.1.5) f ′w(X) = 2r Im b
(Re b)2 + (Im b− r)2
X + w Re b (Im b+ r)
2r Im b
(Re b)2 + (Im b− r)2
because Im b 6= r, we have
(3.1.6) ∀X ∈ R+, fw(X) ≥ fw(0) = w2 > 0,
if w 6= 0 and
− w Re b (Im b+ r)
2r Im b
(Re b)2 + (Im b− r)2
) ≤ 0.
The estimate (3.1.6) shows that if r Im b > 0, the equation fw(X) = 0 has no non-
negative solution for all value of the parameter w 6= 0. This proves that the condition
Σ(q) = C induces that r Im b < 0. Using the linear symplectic transformation
(x, ξ) 7→ (|r|−1/2x, |r|1/2ξ),
we obtain the normal forms (ii) and (iii),
α|r|(ξ + ix)(ξ + ηx) with Im η > 0 and α|r|(ξ − ix)(ξ + ηx) with Im η < 0,
where η = |r|−1b. Finally, we can easily check that the numerical ranges of the normal
forms (ii) and (iii) are actually equal to the whole complex plane C. �
Let us notice that the proposition 2.2.2 and the remark following its statement are
some direct consequences of the symplectic invariance property of the Weyl quanti-
zation (see (3.0.3)) and the previous lemma. We can add that as proved after the
lemma 3.1 in [6], the Fredholm indices of the one-dimensional elliptic quadratic dif-
ferential operators with symbols of type (i), (ii) and (iii) are respectively equal to 0,
−2 and 2.
As we have mentioned in the previous section, the results of Theorem 2.2.1 and
Theorem 2.2.2 are already known in the particular case of the rotated harmonic oscil-
lator. The existence of semiclassical quasimodes inducing the presence of semiclassical
pseudospectrum of infinite index in every point of the interior of the numerical range
for the associated semiclassical operator, is a direct consequence of a result proved
by E.B. Davies in [4] (Theorem 1) and; the absence of semiclassical pseudospectrum
of index 2/3 on the boundary of the numerical range has been proved for the ro-
tated harmonic oscillator in [10]4. As we have previously mentioned (see (2.1.10) and
(3.0.2)), the property of non-normality, the order of symbols and the semiclassical
pseudospectra of elliptic quadratic differential operators are symplectically invariant.
These properties allow us to reduce by any real linear symplectic transformations the
symbols of the elliptic quadratic differential operators that we consider in our proof of
the theorem 2.2.1 and the theorem 2.2.2. By using the lemma 3.1.1, we deduce from
the results of the theorem 2.2.1 and the theorem 2.2.2 proved for the rotated harmonic
oscillator that they are therefore also fulfilled by all non-normal one-dimensional el-
liptic quadratic differential operators with a numerical range different from the whole
complex plane.
We now consider the multidimensional case. As we will see in the following, there is
a real jump of complexity between the one-dimensional case and the multidimensional
one. This jump is among other things a consequence of the complexity increase of
symplectic geometry in dimension n ≥ 2 and the larger diversity appearing in the
class of elliptic quadratic differential operators.
4Let us recall that the value of the order is equal to 2 in this case.
3.2. Case of dimension n ≥ 2. We only need to study the case of a non-normal
elliptic quadratic differential operator
(3.2.1) q(x, ξ)w : B → L2(Rn),
in dimension n ≥ 2. Let us recall that in this case, the numerical range Σ(q) is a
closed angular sector with a top in 0 and a positive opening strictly lower than π, and
that the proposition 2.1.2 gives that
(3.2.2) ∃(x0, ξ0) ∈ R2n, {Re q, Im q}(x0, ξ0) 6= 0.
Let us begin by studying what occurs at the interior of the numerical range Σ̊(q).
3.2.1. On the pseudospectrum at the interior of the numerical range. To prove the
existence of semiclassical quasimodes for the associated semiclassical operator given
by the theorem 2.2.1, we need a first purely algebraic step to characterize the points
belonging to the interior of the numerical range.
Let us consider the following decomposition of the numerical range
(3.2.3) Σ(q) = Ã ⊔ B̃,
where
(3.2.4) Ã =
z ∈ Σ(q) : ∃(x0, ξ0) ∈ R2n, z = q(x0, ξ0), {Re q, Im q}(x0, ξ0) 6= 0
(3.2.5) B̃ =
z ∈ Σ(q) : z = q(x0, ξ0) ⇒ {Re q, Im q}(x0, ξ0) = 0
The next section is devoted to give a geometrical description of these two sets. We
establish using purely algebraic arguments that
(3.2.6) Ã = Σ̊(q) and B̃ = ∂Σ(q).
This result is a consequence of the geometry induced by the quadratic setting to which
the studied symbols belong.
Let us begin by noticing that the symplectic invariance of the Poisson bracket
(2.1.10) induces the same property for the sets à and B̃. We can therefore use some
real linear symplectic transformation to reduce the symbol q. Since
{Re(zq), Im(zq)} = |z|2{Re q, Im q},
we deduce from this symplectic invariance, from the proposition 2.1.1 and the lemma
18.6.4 in [7] that after a similitude, we can reduce our study to the case where
(3.2.7) Re q(x, ξ) =
j + x
with λj > 0 for all j = 1, ..., n.
3.2.1.a. Geometrical description of the sets à and B̃. We begin by proving the fol-
lowing inclusion
(3.2.8) ∂Σ(q) ⊂ B̃.
Let us consider z ∈ ∂Σ(q) and (x0, ξ0) ∈ R2n such that z = q(x0, ξ0). This is
possible because the numerical range is a closed angular sector. If z = 0, the ellipticity
property of q implies that
(x0, ξ0) = (0, 0) and {Re q, Im q}(x0, ξ0) = 0,
because this Poisson bracket is also a quadratic form. This proves that z ∈ B̃. If
z ∈ ∂Σ(q) \ {0},
let us consider the global solution Y of the linear Cauchy problem
(3.2.9)
Y ′(t) = HRe q
Y (t)
Y (0) = (x0, ξ0),
associated to the Hamilton vector field of the symbol Re q,
HRe q =
(∂Re q
− ∂Re q
It is actually a linear Cauchy problem since Re q is a quadratic form. Setting
f(t) = Im q
Y (t)
a direct computation gives that
f ′(0) = {Re q, Im q}(x0, ξ0).
If f ′(0) 6= 0, we could find t0 6= 0 such that
|f(t0)| > |f(0)| = |Im z|.
Since Y is the flow associated to the Hamilton vector field of Re q, the quadratic form
Re q is constant under it. It follows that for all t ∈ R,
Y (t)
= Re q
Y (0)
= Re z
and provides a contradiction because, since z ∈ ∂Σ(q) \ {0}, this would imply in view
of the shape of the numerical range Σ(q) (see Figure 7) that
Y (t0)
6∈ Σ(q).
It follows that the Poisson bracket {Re q, Im q}(x0, ξ0) is necessary equal to 0 and
Figure 7.
q(Y (t
that z ∈ B̃. This ends the proof of the inclusion (3.2.8).
Let us now assume that
(3.2.10) ∂Σ(q) ⊂ B̃, ∂Σ(q) 6= B̃.
In this case, we could find
(3.2.11) z ∈ B̃ \ ∂Σ(q).
Let us first notice that z is necessary non-zero since 0 ∈ ∂Σ(q), and that Re z > 0,
since from (3.2.7),
(3.2.12) Σ(q) \ {0} ⊂ {z ∈ C∗ : Re z > 0}.
The fact that z belongs to the set B̃ implies that
(3.2.13)
Re q(x, ξ) = Re z
Im q(x, ξ) = Im z
=⇒ {Re q, Im q}(x, ξ) = 0.
We also know that there exists at least one solution to the system appearing in the
left-hand-side of (3.2.13). Since from (3.2.7), the quadratic form Re q is positive
definite, we can simultaneously reduce the quadratic forms Re q and Im q by finding
an isomorphism P of R2n such that in the new coordinates y = P−1(x, ξ),
(3.2.14) Re q(Py) =
y2j and Im q(Py) =
j with α1 ≤ ... ≤ αn.
Let us now consider the following quadratic form
(3.2.15) p(y) = {Re q, Im q}(Py).
We get from (3.2.13) and (3.2.14) that
(3.2.16)
j=1 y
j = Re z
j=1 αjy
j = Im z
=⇒ p(y) = 0.
Let us underline that the isomorphism P is not a priori a symplectic transformation
and that it does not preserve the Poisson bracket {Re q, Im q}.
We consider the two following sets
(3.2.17) E1 =
y ∈ R2n : r(y) = 0
where
(3.2.18) r(y) =
(3.2.19) E2 =
y ∈ R2n : p(y) = 0
The next lemma gives a first inclusion between these two sets E1 and E2.
Lemma 3.2.1. We have
(3.2.20) E1 ⊂ E2.
Proof of Lemma 3.2.1. Let y ∈ E1. If y = 0 then y belongs to E2 since from (3.2.15),
p is a quadratic form in the variable y. If y 6= 0, we set
y2j > 0 and ∀j = 1, ..., 2n, ỹj =
We recall from (3.2.12) that z ∈ B̃ \ ∂Σ(q) implies that Re z > 0. Then, since, on
one hand
ỹ2j = Re z,
and that, on the other hand, we have from (3.2.17) and (3.2.18) that
αj ỹ
y2j = Im z,
because y ∈ E1, we deduce from (3.2.16) and the homogeneity of degree 2 of the
quadratic form p that
p(ỹ) =
p(y) = 0.
According to (3.2.19), this proves that y ∈ E2 and ends the proof of the lemma 3.2.1.�
Then, we can notice from (3.2.14) that the boundary of the numerical range ∂Σ(q)
is given by
(3.2.21) (1 + iα1)R+ ∪ (1 + iαn)R+.
Since the numerical range Σ(q) is a closed set, the assumption
z ∈ B̃ \ ∂Σ(q) ⊂ Σ(q) \ ∂Σ(q) = Σ̊(q),
induces from (3.2.21) that
∈]α1, αn[.
This implies that the signature (r1, s1) of the quadratic form r defined in (3.2.18)
fulfills
(3.2.22) (r1, s1) ∈ N∗ × N∗ and r1 + s1 ≤ 2n.
Thus, we can assume after a new labeling that
(3.2.23) r(y) = a1y
1 + ...+ ar1y
− ar1+1y2r1+1 − ...− ar1+s1y
r1+s1
with aj > 0 for all j = 1, ..., r1+ s1. It follows from (3.2.17) and (3.2.23) that in these
new coordinates, the set E1 is the direct product of a proper cone C of R
r1+s1 and
R2n−r1−s1 ,
(3.2.24) E1 = C × R2n−r1−s1 .
Figure 8.
We are now going to prove that the two sets E1 and E2 are equal
(3.2.25) E1 = E2.
Let us reason by the absurd by assuming that it is not the case. Then, we could find
from the lemma 3.2.1,
(3.2.26) y0 ∈ E2 \ E1, y0 = (y′0, y′′0 ) with y′0 ∈ Rr1+s1 , y′′0 ∈ R2n−r1−s1 .
We deduce from (3.2.24) that y′0 6∈ C. Let us now recall an elementary geometrical
fact that we will use several times. This fact is that the intersection of a real line and
a real quadric surface is reduced to either 0, 1 or 2 points, or the line is completely
contained in the quadric surface. We first begin by proving that
(3.2.27) Rr1+s1 × {y′′ = y′′0} ⊂ E2.
Indeed, let us consider the affine subspace
F = {y ∈ R2n : y = (y′, y′′) ∈ Rr1+s1 × R2n−r1−s1 , y′′ = y′′0}.
We identify for more simplicity the space F to the space Rr1+s1 . We agree to say that
a point x′0 of R
r1+s1 belongs to the set E2 to mean that the point (x
0 ) belongs to
the set E2. With this convention, it is sufficient for proving the inclusion (3.2.27) to
consider some particular lines of Rr1+s1 , containing the point y′0 defined in (3.2.26)
and, which have an intersection with the cone C in at least two other different points
u′0 and v
0 (see Figure 9). These lines are necessary contained in the quadric surface
E2 because from the lemma 3.2.1,
E1 ⊂ E2,
and that there are at least three different points of intersection between these lines
and the quadric surface E2,
(u′0, y
0 ) ∈ C × R2n−r1−s1 = E1 ⊂ E2, (v′0, y′′0 ) ∈ C × R2n−r1−s1 = E1 ⊂ E2,
and (y′0, y
0 ) ∈ E2. Thus, we prove that the shaded disc appearing on the figure 10
is completely contained in the set E2. By using the cone structure of the set E2,
we can deduce that all the interior of the cone C (see Figure 11) is contained in E2.
Then, using again other particular intersections with some lines as on the figure 12,
we deduce from our identification of the space F to Rr1+s1 that the inclusion (3.2.27)
is fulfilled.
Figure 9.
We now prove that under these conditions, we have the identity
(3.2.28) E2 = R
Indeed, let us consider (ỹ′0, ỹ
0 ) ∈ R2n = Rr1+s1 × R2n−r1−s1 . If ỹ′0 ∈ C, then
(ỹ′0, ỹ
0 ) ∈ E2,
Figure 10.
These three points belong to E2.
The line D is contained in E2.
Figure 11.
because from (3.2.20) and (3.2.24), (ỹ′0, ỹ
0 ) ∈ E1 and E1 ⊂ E2. If, on the other hand
ỹ′0 6∈ C, we can choose a point u ∈ Rr1+s1 different from ỹ′0 such that u 6∈ C, and such
that the line containing ỹ′0 and u in R
r1+s1 , has an intersection with C in at least two
other different points v and w (see Figure 13). Thus, we can find some distinct real
numbers t1, t2 ∈ R \ {0, 1} such that
v = (1− t1)ỹ′0 + t1u ∈ C and w = (1− t2)ỹ′0 + t2u ∈ C.
Considering now the line
(1− t)(ỹ′0, ỹ′′0 ) + t(u, y′′0 ) : t ∈ R
we can notice that this real line contains at least three different points of E2:
(v, (1 − t1)ỹ′′0 + t1y′′0 ), (w, (1 − t2)ỹ′′0 + t2y′′0 ) and (u, y′′0 ).
Indeed, this is a consequence of the fact that v and w belong to C, and from (3.2.20),
(3.2.24) and (3.2.27). Thus, the line D is contained in the quadric surface E2. This
implies that (ỹ′0, ỹ
0 ) ∈ D ⊂ E2.
To sum up, we have proved that if the two sets E1 and E2 are different then the
set E2 is equal to R
2n. This fact induces in view of (3.2.19) that the quadratic form p
is identically equal to zero. By coming back to the first coordinates (x, ξ) = Py, it
Figure 12.
Figure 13.
follows from (3.2.15) that the quadratic form {Re q, Im q} is also identically equal to
zero, which contradicts (3.2.2). This proves the identity (3.2.25),
E1 = E2.
With this fact, we can resume our first reasoning by the absurd, which assume in
(3.2.11) the existence of a point z ∈ B̃ \ ∂Σ(q). Let us now consider y0 6∈ E1 = E2.
This is possible according to (3.2.2), (3.2.15) and (3.2.19). We deduce from (3.2.17)
and (3.2.19) that r(y0) and p(y0) are non-zero. By considering λ ∈ R∗ such that
p(y0) = λr(y0)
(3.2.29) r̃(y) = p(y)− λr(y),
it follows from (3.2.17), (3.2.19), (3.2.25) and (3.2.29) that
(3.2.30) E1 ⊂ {y ∈ R2n : r̃(y) = 0}.
This inclusion (3.2.30) is strict since
r̃(y0) = 0 and y0 6∈ E1.
By using now exactly the same reasoning as the one previously described to prove
(3.2.25), about the intersections of real lines and quadric surfaces, we prove that the
quadratic form r̃ is necessary identically equal to zero. Then, it follows from (3.2.29)
(3.2.31) p = λr.
By coming back to the first coordinates (x, ξ) = Py, we get using (3.2.14), (3.2.15),
(3.2.18) and (3.2.31) that for all (x, ξ) ∈ R2n,
(3.2.32) {Re q, Im q}(x, ξ) = λ
Im q(x, ξ)− Im z
Re q(x, ξ)
Let us now consider (x0, ξ0) ∈ R2n such that q(x0, ξ0) ∈ ∂Σ(q) \ {0}. This is possible
since the numerical range Σ(q) is a closed angular sector with a top in 0 and a positive
opening. We deduce from (3.2.5) and (3.2.8) that we necessarily have
{Re q, Im q}(x0, ξ0) = 0.
This induces from (3.2.32) that
(3.2.33) Im q(x0, ξ0) =
Re q(x0, ξ0),
because λ ∈ R∗. Since according to the shape of the numerical range Σ(q) and
(3.2.12),
q(x0, ξ0) ∈ ∂Σ(q) \ {0} ⊂ {z ∈ C : Re z > 0},
the identity (3.2.33) proves that the point z also belongs to the set ∂Σ(q), but it
contradicts the initial assumption
z ∈ B̃ \ ∂Σ(q).
Finally, this ends our reasoning by the absurd and proves (3.2.6).
3.2.1.b. Existence of semiclassical quasimodes at the interior of the numerical range.
To prove the existence of semiclassical quasimodes for the associated semiclassical
operator
(q(x, hξ)w)0<h≤1,
in every point of the numerical range’s interior (Theorem 2.2.1), we use an existence
result of semiclassical quasimodes for general pseudodifferential operators violating
the condition (Ψ)5. Let us mention that this result generalizes the two existence
results of semiclassical quasimodes given by E.B. Davies, in the case of Schrödinger
operators (Theorem 1 in [4]), and by M. Zworski in [17] and [18], for pseudodifferential
operators.
This existence result of semiclassical quasimodes can be stated as follows. Let us
consider a semiclassical symbol P (x, ξ;h) in S(〈(x, ξ)〉m, dx2 + dξ2) with m ∈ R+,
〈(x, ξ)〉2 = 1 + x2 + ξ2,
5The definition of the condition (Ψ) is recalled below.
where S(〈(x, ξ)〉m, dx2 + dξ2) stands for the following symbol class
S(〈(x, ξ)〉m, dx2 + dξ2) =
a(x, ξ;h) ∈ C∞(Rnx × Rnξ ,C) :
∀α ∈ N2n, sup
0<h≤1
‖〈(x, ξ)〉−m∂αx,ξa(x, ξ;h)‖L∞(R2n) < +∞
with a semiclassical expansion
(3.2.34) P (x, ξ;h) ∼
hjpj(x, ξ),
where for all j ∈ N, pj is a symbol of the class S(〈(x, ξ)〉m, dx2 + dξ2) independent
from the semiclassical parameter h.
Let z ∈ C, we assume that there exists a function q0 ∈ C∞b (R2n,C), where
C∞b (R
2n,C) stands for the set of bounded complex-valued functions on R2n with
all derivatives bounded, and a bicharacteristic curve, t ∈ [a, b] 7→ γ(t), of the real part
Re(q0(p0 − z)) of the symbol q0(p0 − z), with a < b, such that
(3.2.35) ∀t ∈ [a, b], q0
6= 0 and
q0(γ(a))
p0(γ(a))− z
> 0 > Im
q0(γ(b))
p0(γ(b))− z
Theorem 3.2.1. Under these assumptions (3.2.34) and (3.2.35), for all open neigh-
bourhood V of the compact set γ([a, b]) in R2n and for all N ∈ N, there exist h0 > 0
and (uh)0<h≤h0 a semiclassical family in S(Rn) such that
‖uh‖L2(Rn) = 1, FS
(uh)0<h≤h0
⊂ V and ‖P (x, hξ;h)wuh − zuh‖L2(Rn) = O(hN ),
when h → 0+.
The notation FS
(uh)0<h≤h0
stands for the frequency set of the semiclassical fam-
ily (uh)0<h≤h0 defined as the complement in R
2n of the set composed by the points
(x0, ξ0) ∈ R2n, for which there exists a symbol χ0(x, ξ;h) ∈ S(1, dx2 + dξ2) such that
χ0(x0, ξ0;h) = 1 and ‖χ0(x, hξ;h)wuh‖L2(Rn) = O(h∞),
when h → 0+.
This existence result of semiclassical quasimodes is an adaptation in a semiclassical
setting of the proof given by L. Hörmander in [7] for proving that the condition (Ψ) is a
necessary condition for the solvability of a pseudodifferential operator (Theorem 26.4.7
in [7]). The existence of this result has been first mentioned in [5]. A complete proof of
this adaptation in a semiclassical setting is given in [11]. This result shows that when
the principal symbol p0−z of the symbol P−z violates the condition (Ψ), there exists
in this point z some semiclassical quasimodes inducing the presence of semiclassical
pseudospectrum of infinite index for the semiclassical operator P (x, hξ;h)w.
Condition (Ψ). A complex-valued function p ∈ C∞(R2n,C) fulfills the condition (Ψ)
if there is no complex-valued function q ∈ C∞(R2n,C) such that the imaginary part
Im(qp) of the function qp changes sign from positive values to negative ones along
an oriented bicharacteristic of the symbol Re(qp) on which the function q does not
vanish.
By using the characterization given in the previous section for the interior of the
numerical range Σ̊(q) (see (3.2.4) and (3.2.6)), we are now going to prove that the
principal symbol q(x, ξ) − z of the semiclassical operator
q(x, hξ)w − z,
violates the condition (Ψ) for all z in Σ̊(q). This violation of the condition (Ψ) will
induce in view of the theorem 3.2.1 that for all z ∈ Σ̊(q) and N ∈ N, we can find a
semiclassical quasimode (uh)0<h≤h0 ∈ S(Rn), with h0 > 0, verifying
‖uh‖L2(Rn) = 1 and ‖q(x, hξ)wuh − zuh‖L2(Rn) = O(hN ) when h → 0+,
which will end the proof of Theorem 2.2.1.
Let us consider z ∈ Σ̊(q). We are now going to prove that there is actually a
violation of the condition (Ψ) for the symbol q − z. According to (3.2.4) and (3.2.6),
there are two cases to separate.
Case 1. Let us assume that there exists (x0, ξ0) ∈ R2n such that
(3.2.36) z = q(x0, ξ0), {Re(q − z), Im(q − z)}(x0, ξ0) = {Re q, Im q}(x0, ξ0) < 0.
By considering the solution of the following Cauchy problem
(3.2.37)
Y ′(t) = HRe q
Y (t)
Y (0) = (x0, ξ0),
we define the following function
(3.2.38) f(t) = Im q
Y (t)
− Im q(x0, ξ0).
As mentioned before, (3.2.37) is a linear Cauchy problem. It follows that its solution
Y is global and that the function f is well-defined on R. A direct computation using
(3.2.37) and (3.2.38) gives that for all t ∈ R,
(3.2.39) f ′(t) = {Re q, Im q}
Y (t)
Since from (3.2.36), (3.2.37), (3.2.38) and (3.2.39),
f(0) = 0, f ′(0) = {Re q, Im q}(x0, ξ0) < 0
and HRe q−Re z = HRe q, we deduce in this first case that the imaginary part of the
function q − z changes sign, at the first order, from positive values to negative ones
along the oriented bicharacteristic Y of the symbol Re q−Re z. This proves that the
symbol q − z actually violates the condition (Ψ).
Case 2. Let us now assume that there exists (x0, ξ0) ∈ R2n such that
(3.2.40) z = q(x0, ξ0), {Re(q − z), Im(q − z)}(x0, ξ0) = {Re q, Im q}(x0, ξ0) > 0.
We consider as in the previous case, the global solution Y of the Cauchy problem
(3.2.37) and the function f defined in (3.2.38). Since from (3.2.37), (3.2.38), (3.2.39)
and (3.2.40),
(3.2.41) f(0) = 0, f ′(0) = {Re q, Im q}(x0, ξ0) > 0,
we deduce this time that the imaginary part of the function q − z also changes sign,
at the first order, along the oriented bicharacteristic Y of the symbol Re q − Re z.
Nevertheless, this change of sign is done in the “wrong” way. It is a change of sign
from negative values to positive ones, which does not induce directly a violation of
the condition (Ψ). To check that there is actually a violation of the condition (Ψ)
in this second case, we need to study more precisely the behaviour of the function
Im q − Im z along this bicharacteristic Y .
We deduce from (3.2.41) that there exists ε > 0 such that
∀t ∈ [−ε, ε], f ′(t) > 0,
which induces that
(3.2.42) f(ε) > 0 and f(−ε) < 0,
since from (3.2.41), f(0) = 0. By using the following lemma, we obtain that for all
δ > 0, there exists a time t0(δ) > ε such that
(3.2.43) |Y
t0(δ)
− Y (−ε)| < δ.
Figure 14.
q(Y (�"))
z = q(Y (0))
q(Y ("))
Lemma 3.2.2. If Y (t) = (x(t), ξ(t)) is the C∞(R,R2n) function solving the linear
system of ordinary differential equations
Y ′(t) = HRe q
Y (t)
where Re q is the symbol defined in (3.2.7), then we have
∀t0 ∈ R, ∀ε > 0, ∀M > 0, ∃T1 > M, ∃T2 > M,
|Y (t0)− Y (t0 + T1)| < ε and |Y (t0)− Y (t0 − T2)| < ε.
Proof of Lemma 3.2.2. If Y (t0) = (a1, ..., an, b1, ..., bn) ∈ R2n, we deduce from (3.2.7)
that the function Y (t) = (x(t), ξ(t)) solves the following Cauchy problem
∀j = 1, ..., n,
x′j(t) = 2λjξj(t)
ξ′j(t) = −2λjxj(t)
xj(t0) = aj
ξj(t0) = bj.
It follows that for all j = 1, ..., n and t ∈ R,
(3.2.44)
xj(t) = bj sin
2(t− t0)λj
+ aj cos
2(t− t0)λj
ξj(t) = bj cos
2(t− t0)λj
− aj sin
2(t− t0)λj
Setting βj = λj/π for all j = 1, ..., n, we need to study two different cases.
Case 1: ∀j ∈ {1, ..., n}, βj ∈ Q. In this case, the function Y is periodic and the
result of Lemma 3.2.2 is obvious.
Case 2: (β1, ..., βn) 6∈ Qn. In this second case, we use the following classical result of
rational approximation: ∀ε > 0, ∀(θ1, ..., θn) ∈ Rn \Qn, ∃p1, ..., pn ∈ Z, ∃q ∈ N∗ such
0 < sup
j=1,...,n
If 0 < ε1 < 1/2, we can therefore find some integers p1,1, ..., p1,n ∈ Z and qε1 ∈ N∗
such that
0 < sup
j=1,...,n
|qε1βj − p1,j | < ε1.
j=1,...,n
|qε1βj − p1,j | > 0,
using again this result of rational approximation, we can find some other integers
p2,1, ..., p2,n ∈ Z and qε2 ∈ N∗ such that
0 < sup
j=1,...,n
|qε2βj − p2,j | < ε2.
By using this process, we build some sequences (pm,j)m∈N∗ of Z for j = 1, ..., n,
(εm)m∈N∗ of R
+ and (qεm)m∈N∗ of N
∗ such that for all m ≥ 2,
(3.2.45) 0 < sup
j=1,...,n
|qεmβj − pm,j | < εm =
j=1,...,n
∣qεm−1βj − pm−1,j
(3.2.46) 0 < εm <
The elements of the sequence (qεm)m∈N∗ are necessary two by two different. Indeed,
if qεk = qεl for k < l, this would imply according to (3.2.45) and (3.2.46) that
∀j = 1, ..., n, |pk,j − pl,j| ≤ |qεkβj − pk,j |+ |qεlβj − pl,j | < εk + εl < 1,
because 0 < ε1 < 1/2, which would induce that ∀j = 1, ..., n, pk,j = pl,j because pk,j
and pl,j are some integers; and would contradict (3.2.45) because
0 < sup
j=1,...,n
|qεlβj − pl,j | < εl ≤
j=1,...,n
|qεkβj − pk,j |.
Since the sequence (qεm)m∈N∗ is composed of integers two by two different, we can
assume after a possible extraction that qεm → +∞ when m → +∞. We deduce from
(3.2.44), (3.2.45) and (3.2.46) that
Y (t0 + qεm) → Y (t0) when m → +∞.
Then, considering (β̃1, ..., β̃n) = (−β1, ...,−βn), we obtain by using the same method
a sequence (q̃εm)m∈N∗ of integers such that q̃εm → +∞ and
Y (t0 − q̃εm) → Y (t0) when m → +∞.
This ends the proof of Lemma 3.2.2. �
Since from (3.2.42), f(−ε) < 0, we deduce from (3.2.38) and (3.2.43) that there
exists t0 > ε such that f(t0) is arbitrarily close to f(−ε). It follows in particular that
we can find t0 > ε such that f(t0) < 0. Since from (3.2.42), f(ε) > 0 and f(t0) < 0,
we deduce from (3.2.38) and (3.2.40) that the function
t 7→ Im q
Y (t)
− Im z,
changes sign from positive values to negative ones on the interval [ε, t0]. This proves
that the imaginary part of the function q−z actually changes sign from positive values
to negative ones along the oriented bicharacteristic Y of the symbol Re q−Re z; and
that the symbol q − z also violates in this second case the condition (Ψ). This ends
the proof of Theorem 2.2.1.
3.2.1.c. Another proof for the existence of semiclassical quasimodes. In the following
lines, we give another proof for the existence of semiclassical quasimodes in some
points of the numerical range’s interior. The result proved in this section is weaker
than the one given by the theorem 2.2.1, since we prove the existence of semiclassical
quasimodes in every point of the numerical range’s interior without a finite number
of particular half-lines.
Let us consider a non-normal elliptic quadratic differential operator
(3.2.47) q(x, ξ)w : B → L2(Rn),
in dimension n ≥ 2. We assume, as before, that (3.2.7) is fulfilled. Using that
the quadratic form Re q is positive definite, we can simultaneously reduce the two
quadratic forms Re q and Im q by choosing an isomorphism P of R2n such that in
the new coordinates y = P−1(x, ξ),
(3.2.48) r1(y) = Re q(Py) =
y2j , r2(y) = Im q(Py) =
with α1 ≤ ... ≤ αn. Let us study when the differential forms dr1(y) and dr2(y) are
linearly dependent on R i.e. when there exist (λ, µ) ∈ R2 \ {(0, 0)} such that
(3.2.49) λdr1(y) + µdr2(y) = 0.
It follows from (3.2.48) and (3.2.49) that for all j = 1, ..., 2n,
(3.2.50) (λ+ µαj)yj = 0.
If y 6= 0, then there exists j0 ∈ {1, ..., 2n} such that yj0 6= 0. This implies that
(3.2.51) λ+ µαj0 = 0.
We deduce from (3.2.50) and (3.2.51) that yj = 0 if αj 6= αj0 . Thus, we obtain that if
z ∈ Σ̊(q) \
(1 + iα1)R
+ ∪ ... ∪ (1 + iαn)R∗+
then the differential forms dRe q and dImq are linearly independent on R in every
point of the set q−1(z).
Figure 15.
(1 + i�
(1 + i�
(1 + i�
(1 + i�
(1 + i�
Let us consider such a point
z ∈ Σ̊(q) \
(1 + iα1)R
+ ∪ ... ∪ (1 + iαn)R∗+
Since the dimension n ≥ 2, we can apply the lemma 3.1 in [5] (see also the lemma 8.1
in [9]). It follows that for any compact, connected component Γ of q−1(z), we have
(3.2.52)
{Re q, Im q}(ρ)λq,z(dρ) = 0,
where λq,z stands for the Liouville measure on q
−1(z),
λq,z ∧ dRe q ∧ dIm q =
The set q−1(z) is a non-empty submanifold of codimension 2 in R2n. We deduce from
(3.2.4) and (3.2.6) that there exist (x0, ξ0) ∈ q−1(z) such that
(3.2.53) {Re q, Im q}(x0, ξ0) 6= 0.
Then, it follows from (3.2.52) and (3.2.53) that there necessary exists (x̃0, ξ̃0) ∈ q−1(z)
such that
(3.2.54) {Re q, Im q}(x̃0, ξ̃0) < 0.
Under this condition (3.2.54), we can use the reasoning given in the first studied case
(see (3.2.36)) to prove that the imaginary part of the function q−z changes sign, at the
first order, from positive values to negative ones along an oriented bicharacteristic of
the symbol Re q−Re z. This induces that the symbol q−z violates the condition (Ψ);
and we can conclude by using the theorem 3.2.1. Let us mention that we can also
directly use the existence result of semiclassical quasimodes given by M. Zworski in
[17] and [18]. This second proof gives the existence of semiclassical quasimodes in
every point belonging to the set
Σ̊(q) \
(1 + iα1)R
+ ∪ ... ∪ (1 + iαn)R∗+
3.2.2. On the pseudospectrum at the boundary of the numerical range. In this section,
we give a proof of the theorem 2.2.2. Let us consider a non-normal elliptic quadratic
differential operator
q(x, ξ)w : B → L2(Rn),
in dimension n ≥ 1. We assume that Σ(q) 6= C, and that its Weyl symbol q(x, ξ) is of
finite order kj on a half-line ∆j , j ∈ {1, 2} (See the definition given in (2.2.9)), which
composes the boundary of its numerical range
(3.2.55) ∂Σ(q) = {0} ⊔∆1 ⊔∆2.
As we have already done several times, we can reduce our study to case where (3.2.7)
is fulfilled.
Proof of Theorem 2.2.2. Let us consider the following symbol belonging to the
C∞b (R
2n,C) space, composed of bounded complex-valued functions on R2n with all
derivatives bounded
(3.2.56) r(x, ξ) =
q(x, ξ) − z
1 + x2 + ξ2
with z ∈ ∆j . Setting Σ̃(r) = r(R2n), we can first notice that
z ∈ ∂Σ(q) \ {0} ⇒ 0 ∈ ∂Σ̃(r).
Let us also notice that the symbol r fulfills the principal-type condition in 0. Indeed,
if (x0, ξ0) ∈ R2n was such that r(x0, ξ0) = 0 and dr(x0, ξ0) = 0, we would get from
(3.2.56) that
(3.2.57) dq(x0, ξ0) = 0.
Since from (3.2.7) and (3.2.57), we have
dRe q(x0, ξ0) = 2
(x0)jdxj + (ξ0)jdξj
this would imply that
(x0, ξ0) = (0, 0), q(x0, ξ0) = 0,
because q is a quadratic form and that λj > 0 for all j = 1, ..., n. On the other hand,
since r(x0, ξ0) = 0, we get from (3.2.56) that q(x0, ξ0) = z 6= 0 because
z ∈ ∆j ⊂ ∂Σ(q) \ {0},
which induces a contradiction. It follows that the symbol r actually fulfills the
principal-type condition in 0. Let us notice that, since symbol q is of finite order kj
in z, this induces in view of (3.2.56) that the symbol r is also of finite order kj in 0.
On the other hand, we deduce from (3.2.7) and (3.2.56) that the set
{(x, ξ) ∈ R2n : r(x, ξ) = 0} = {(x, ξ) ∈ R2n : q(x, ξ) = z},
is compact. Under these conditions, we can apply the theorem 1.4 in [5], which proves
that the integer kj is even and gives the existence of positive constants h0 and C1
such that
(3.2.58) ∀ 0 < h < h0, ∀u ∈ S(Rn), ‖r(x, hξ)wu‖L2(Rn) ≥ C1h
kj+1 ‖u‖L2(Rn).
Remark. We did not check the dynamical condition (1.7) in [5], because this assump-
tion is not necessary for the proof of Theorem 1.4. Indeed, this proof only use a part
of the proof of lemma 4.1 in [5] (a part of the second paragraph), where this condition
(1.7) is not needed.
By using some results of symbolic calculus given by Theorem 18.5.4 in [7] and (3.2.56),
we can write
(3.2.59) r(x, hξ)w(1 + x2 + h2ξ2)w = q(x, hξ)w − z + hr1(x, hξ)w + h2r2(x, hξ)w ,
(3.2.60) r1(x, ξ) = −ix
(x, ξ) + iξ
(x, ξ)
(3.2.61) r2(x, ξ) = −
(x, ξ) − 1
(x, ξ).
We can easily check from (3.2.56) that these functions r1 and r2 belong to the space
C∞b (R
2n,C), and we deduce from the Calderón-Vaillancourt theorem that there exists
a positive constant C2 such that for all u ∈ S(Rn) and 0 < h ≤ 1,
(3.2.62) ‖r1(x, hξ)wu‖L2 ≤ C2‖u‖L2 and ‖r2(x, hξ)wu‖L2 ≤ C2‖u‖L2.
It follows from (3.2.58), (3.2.59), (3.2.62) and the triangular inequality that for all
u ∈ S(Rn) and 0 < h < h0,
kj+1 ‖(1 + x2 + h2ξ2)wu‖L2(Rn)
≤ ‖r(x, hξ)w(1 + x2 + h2ξ2)wu‖L2(Rn)
≤ ‖q(x, hξ)wu− zu‖L2(Rn) + C2h(1 + h)‖u‖L2(Rn).
Since from the Cauchy-Schwarz inequality, we have for all u ∈ S(Rn) and 0 < h ≤ 1,
‖u‖2L2(Rn) ≤ ‖u‖2L2(Rn) + ‖xu‖2L2(Rn) + ‖hDxu‖2L2(Rn)
(1 + x2 + h2ξ2)wu, u
L2(Rn)
≤ ‖(1 + x2 + h2ξ2)wu‖L2(Rn)‖u‖L2(Rn),
we obtain that for all u ∈ S(Rn) and 0 < h < h0,
(3.2.63) C1h
kj+1 ‖u‖L2(Rn) ≤ ‖q(x, hξ)wu− zu‖L2(Rn) + C2h(1 + h)‖u‖L2(Rn).
Since kj ≥ 1, we deduce from (3.2.63) that there exist some positive constants h′0 and
C3 such that for all 0 < h < h
0 and u ∈ S(Rn),
‖q(x, hξ)wu− zu‖L2(Rn) ≥ C3h
kj+1 ‖u‖L2(Rn).
Using that the Schwartz space S(Rn) is dense in B and that the operator
q(x, hξ)w + z,
is a Fredholm operator of index 0, we obtain that for all 0 < h < h′0,
q(x, hξ)w − z
∥ ≤ C−13 h
kj+1 ,
which ends the proof of Theorem 2.2.2. �
About the case of infinite order, the situation is much more complicated. As
mentioned before, we cannot expect to prove a stronger result than an absence of
semiclassical pseudospectrum of index 1, but we can actually prove that there is never
some semiclassical pseudospectrum of index 1 on every half-line of infinite order, by
using a result of exponential decay in time for the norm of contraction semigroups
generated by elliptic quadratic differential operators proved in [12].
The result proved in [12] shows that the norm of a contraction semigroup
‖etq(x,ξ)
‖L(L2), t ≥ 0,
generated by an elliptic quadratic differential operator q(x, ξ)w with a Weyl symbol
verifying
Re q ≤ 0, ∃(x0, ξ0) ∈ R2n, Re q(x0, ξ0) 6= 0,
decreases exponentially in time
(3.2.64) ∃M,a > 0, ∀t ≥ 0, ‖etq(x,ξ)
‖L(L2) ≤ Me−at.
Let us consider a non-normal elliptic quadratic differential operator
q(x, ξ)w : B → L2(Rn),
in dimension n ≥ 1 such that Σ(q) 6= C. We explain in the following lines how (3.2.64)
allows to prove that there is never some semiclassical pseudospectrum of index 1 on
any open half-lines composing the boundary of the numerical range ∂Σ(q) \ {0}.
Let z ∈ ∂Σ(q)\{0}. Since the numerical range Σ(q) is a closed angular sector with
a top in 0 and a positive opening strictly lower than π, we can find ε ∈ {±1} such
(3.2.65) Re(εiz−1q) ≤ 0, ∃(x0, ξ0) ∈ R2n, Re(εiz−1q)(x0, ξ0) 6= 0.
Using the theorem 2.8 in [2], we obtain that for all η ∈ R,
q(x, ξ)w − ηz
= − iz−1ε
εiη − εiz−1q(x, ξ)w
= − iz−1ε
e−iεηsesεiz
−1q(x,ξ)wds.(3.2.66)
It follows from (3.2.64) and (3.2.65) that for all η ∈ R,
q(x, ξ)w − ηz
∥ ≤ |z|−1
‖esεiz
−1q(x,ξ)w‖L(L2)ds
≤ |z|−1
Me−asds = |z|−1M
< +∞,
which proves the absence of semiclassical pseudospectrum of index 1 on the half-line
zR∗+. We can actually use the theorem 2.8 in [2] because
iR ⊂ C \ σ
εiz−1q(x, ξ)w
Indeed, if it was not the case, we would deduce from (2.1.7) that there exists u0 ∈
B \ {0} and λ0 ∈ R such that
εiz−1q(x, ξ)wu0 = iλ0u0.
Since from (3.2.65), the quadratic form −Re(εiz−1q) is non-negative, we deduce from
the symplectic invariance of the Weyl quantization and the theorem 21.5.3 in [7] that
there exists a metaplectic operator U such that
(3.2.67) −
εiz−1q(x, ξ)
= U−1
+ x2j ) +
j=k+1
with k, l ∈ N and λj > 0 for all j = 1, ..., k. By using that U is a unitary operator on
L2(Rn), we obtain that
0 = − Re(iλ0u0, u0)L2
= − Re
εiz−1q(x, ξ)wu0, u0
εiz−1q(x, ξ)
u0, u0
‖DxjUu0‖2L2 + ‖xjUu0‖2L2
j=k+1
‖xjUu0‖2L2 ,
which induces that u0 = 0, because from (3.2.65) and (3.2.67), k + l ≥ 1. It follows
from (2.1.7) that there exists ε0 > 0 such that
εiz−1q(x, ξ)w
⊂ {z ∈ C : Re z ≤ −ε0}.
References
[1] L.S.Boulton, Non-self-adjoint harmonic oscillator semigroups and pseudospectra, J. Operator
Theory, 47, 413-429 (2002).
[2] E.B.Davies, One-Parameter Semigroups, Academic Press, London (1980).
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[4] E.B.Davies, Semi-classical states for non-self-adjoint Schrödinger operators, Comm. Math.
Phys., 200, 35-41 (1999).
[5] N.Dencker, J.Sjöstrand, M.Zworski, Pseudospectra of Semiclassical (Pseudo-)Differential Op-
erators, Comm. Pure Appl. Math., 57, 384-415 (2004).
[6] L.Hörmander, A Class of Hypoelliptic Pseudodifferential Operators with Double Characteristics,
Math. Ann., 217, 165-188 (1975).
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Verlag (1985).
[8] T.Kato, Perturbation Theory for Linear Operators, Springer-Verlag, Berlin (1980).
[9] A.Melin, J.Sjöstrand, Determinants of pseudodifferential operators and complex deformations
of phase space, Methods Appl. Anal., 9, no.2, 177-237 (2002).
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J. London Math. Soc. (2) 73, 745-761 (2006).
[11] K.Pravda-Starov, Etude du pseudo-spectre d’opérateurs non auto-adjoints, PhD Thesis of the
University of Rennes 1, France (2006).
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(2007).
[13] S.Roch, B.Silbermann, C∗-algebra techniques in numerical analysis, J. Oper. Theory 35, 241-
280 (1996).
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för Mat., 12, 85-130 (1974).
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Math. Phys., 229, 293-307 (2002).
Department of Mathematics, University of California, Evans Hall, Berke-
ley, CA 94720, USA
E-mail address: [email protected]
1. Introduction
1.1. Miscellaneous facts about pseudospectrum
1.2. Elliptic quadratic differential operators
1.3. Semiclassical pseudospectrum
2. Statement of the results
2.1. Some notations and some preliminary facts about elliptic quadratic differential operators
2.2. Statement of the main results
3. The proofs of the results
3.1. The one-dimensional case
3.2. Case of dimension n 2
References
|
0704.0325 | Fluctuation-dissipation relation on a Melde string in a turbulent flow,
considerations on a "dynamical temperature" | 8 Fluctuation-dissipation relation on a Melde string in
a turbulent flow, considerations on a “dynamical
temperature”.
V Grenard, N B Garnier and A Naert.
Université de Lyon, Laboratoire de Physique, École Normale Supérieure de Lyon,
46 Allée d’Italie, 69364 Lyon Cedex 07, France.
E-mail: [email protected]
PACS numbers: 05.70.Ln
PACS numbers: 05.40.-a
PACS numbers: 05.20.Jj
Abstract. We report on measurements of the transverse fluctuations of a string
in a turbulent air jet flow. Harmonic modes are excited by the fluctuating
drag force, at different wave-numbers. This simple mechanical probe makes
it possible to measure excitations of the flow at specific scales, averaged over
space and time: it is a scale-resolved, global measurement. We also measure
the dissipation associated to the string motion, and we consider the ratio of the
fluctuations over dissipation (FDR). In an exploratory approach, we investigate
the concept of effective temperature defined through the FDR. We compare our
observations with other definitions of temperature in turbulence. From the theory
of Kolmogorov (1941), we derive the exponent −11/3 expected for the spectrum
of the fluctuations. This simple model and our experimental results are in good
agreement, over the range of wave-numbers, and Reynolds number accessible
(74000 ≤ Re ≤ 170000).
1. Introduction
Turbulent flows exhibit a notoriously complex and unpredictable dynamics: they
present a huge number of degrees of freedom, and their dynamics are both far from
equilibrium and dissipative [1, 2, 3]. The kinetic energy injected at large scale by shear
instability mecanisms is dissipated into heat by the molecular viscosity at small scales.
That is, dissipation and injection scales are distinct. Therefore, a transport process
through scales is necessary for a flow to be stationary. It is suspected that instability
mechanisms associated with non-linearities generate harmonics, therefore transfering
energy to smaller scales almost without dissipation. An equivalent picture would
consist in vortices stretching each other in such a way that a non-zero energy transfer
occurs toward smaller scales. This picture of cascade process was first proposed by
Richardson [4]. The cascade stops approximately in the range of scales where the
viscosity becomes efficient to damp velocity gradients. In the late thirties, Kolmogorov
derived from this idea a phenomenological theory accounting for the fluctuations of
various observables in fully developed turbulence [5]. In the present work, we are
http://arxiv.org/abs/0704.0325v3
Measurements of a dynamical temperature in turbulence. 2
neither concerned by the large (energy injection) scales, nor by the small (dissipation)
scales, but by the intermediate range. In this intermediate inertial range, we study
the transport process through scales, expected to be universal. Instead of scale l, one
often refers to the wave-number k = 2π/l.
The control parameter of the flow is the Reynolds number: Re = V L
, where L is the
macroscopic scale of the flow (integral scale, or correlation length), V is a characteristic
shear velocity at large scale, and ν is the kinematic viscosity of the fluid. It is also
the mean ratio of the inertial by the dissipative contribution of the forcing over a fluid
particle. Interesting predictions were derived by Kolmogorov (1941), that we use in the
following. Especially, the range of scales over which fluctuations occur scales as Re3/4.
The prediction for the exponent of the power spectral density as 〈|ṽ|2〉 ∝ k−5/3 is
among the most famous successes of this theory [1, 2, 3].
Our experimental system is discribed in detail in the next section. It is a thin string
held by its ends at constant tension across a turbulent flow. To formalize briefly, it
is an oscillator with multiple resonances, coupled to a particular ’thermostat’: the
turbulent flow. This string is used to probe the inertial range of a flow of high enough
Reynolds numbers. The device is ’calibrated’ by measuring the average (complex)
response to an external perturbation, and then used to measure the free fluctuations
caused by turbulence alone. Measurement of the displacement r(t) caused by the
turbulent forcing f(t) is performed with small piezoelectric transducers. We measure
the average response, i.e. the displacement on one end caused by a known broad band
forcing on the other end. Then, measurements of the displacement on one end alone
give information on the forcing fluctuations. Our study goes a step forward, in an
exploratory way. Knowing the average response function of the string and measuring
r(t), we invoque a version of the Fluctuation-Dissipation Theorem extended out of
equilibrium, to define an effective temperature of the turbulent flow. This effective
temperature happends to be scale-dependant.
In this work, fully developped turbulence is addressed from the point of view of
statistical mechanics. We first recall one important break-through: the statement of
the Fluctuation-Dissipation Theorem (FDT). Consider a pair of conjugate variables
(displacement r and force f) of a small system in thermal contact with a large
heat reservoir. In the present case the small system is the string, coupled to the
turbulent flow which is the reservoir. Displacement r and force f are conjugate in
the sense that their product is the work exerted by the flow on the string. The
theorem originates from the idea that spontaneous fluctuations r(t) should have the
same statistical properties as the relaxation of r(t) after the removal of an external
forcing perturbation. The main hypothesis needed to derive this theorem are: –
linear response between f and r, – thermal equilibrium between the system under
consideration and the thermostat, – thermal equilibrium of the thermostat itself. The
response function Hr,f is such that: r(t) =
Hx,f (t − t
′)f(t′)dt′. Equivalently it
can be written in the Fourier space as: r̃(ω) = H̃r,f f̃(ω). Under some hypothesis, the
fluctuations of r (its 2-times correlation function) are linked by a very simple relation
with the dissipative response of the system to a perturbation of the conjugate variable
f (imaginary part of the average response function). It is simply proportional, and
the coefficient is nothing but the temperature multiplied by the Boltzman constant:
kBT [6]. The validity of the hypothesis has to be discussed in each case. If they are
satisfied, the correlation function of the spontaneous fluctuations is proportional to
the response function, i.e. the factor is unique and constant. Moreover, this factor
Measurements of a dynamical temperature in turbulence. 3
is the same for all couples of conjugate variables, and this factor is kBT , where T is
the temperature of the system. The Boltzman constant kB ≃ 1.38 10
−23JK−1 is an
universal constant. This relation can be expressed in spectral variables:
〈|r̃(ω)|2〉 =
2 kBT
Im[H̃r,f (ω)]. (1)
In this expression of the FDT, 〈|r̃(ω)|2〉 is the power spectral density of the fluctuations
of the displacement r, as H̃r,f(ω) is the response function on r to the conjugate
variable f . Because the string is very thin, the drag is purely viscous. It is therefore
proportional to the velocity, which is in quadrature with the displacement. The
dissipation is therefore proportional to the imaginary part of the average response
function: Im[H̃].
In the perspective of constructing a non-equilibrium thermodynamics, the FDT has
been reconsidered by L. Cugliandolo and J. Kurchan, while investigating amorphous
materials relaxing after a thermal quench through the glass transition [7, 8].
We present in the following an exploratory approach of the question of turbulent
fluctuations using their extended formalism. The Fluctuation-Dissipation Ratio
(FDR) can be rewritten:
ω 〈r̃(ω)2〉
Im[H̃r,f(ω)]
= 2 kBTeff.(ω), (2)
where the temperature is replaced by an ’effective’ temperature Teff., function of
frequency ω. The frequency dependence of Teff. expresses the fact that different degrees
of freedom are not at equilibrium with each other, resulting in internal energy fluxes.
In other words, in our system, each (independent) mode of the string couples to
(non-independent) scale of the flow. As the flow is stationary, we average our
measurements on time, and finally obtain the frequency dependance of Teff. as defined
by equation 2. Measurements of the fluctuations of the string give Fourier components
of the excitation of the flow. We measure independently the fluctuations, and the
complex average response function to a specified excitation, in a way discussed below.
We propose to analyse these measurements with the criteria discussed above.
The paper is organised as follows. The next section describes the experimental setup,
turbulent flow properties, and the setting of the string. General properties of a
vibrating Melde string are also discussed. The measurements are shown in section 3:
response, fluctuations, and the Fluctuation Dissipation Ratio of this system. In section
4, we derive from Kolmogorov’s theory a simple scaling model for the fluctuations of
the drag, and therefore the FDR, which accounts for the exponent observed in the
whole range of accessible Re. The section 5 is devoted to a discussion of our results,
especially in comparison to several definitions of temperature in turbulence proposed
in the literature.
2. The Melde string and the experimental setup
The experimental setup is sketched in Fig. 1. A turbulent air jet originates from
a nozzle of diameter 5 cm. The flow facility we used is thoroughly described in
[9]. A thin stainless steel string of length 60 cm is located 2 m downstream the
nozzle, perpendicular to the axis of the flow. At this distance, the length of the
string is about the diameter of the turbulent jet. The displacement of the string is
measured using piezoelectric multi-layer ceramics at each end of the string. A piezo
Measurements of a dynamical temperature in turbulence. 4
is deformed by a voltage. Reciprocally, if the ceramic in compressed, a voltage is
generated. The relation between voltage and deformation is linear, and the frequency
response is almost flat in the frequency range we consider here. It can be used as
actuator or sensor. We have two piezos, one on each end of the string. The two
different measurements we perform are the following. 1) complex response function:
one (input) piezo is feeded with a white noise voltage through a power amplifier. The
source is that of a HP3562A signal analyser. Standing transverse waves appear in
the string, weakly perturbed by the turbulent fluctuations. Mecanical displacement
on the other end is transformed into a voltage by the other (output) piezo. It must
be amplified, and both input and output voltages are recorded synchronously with a
24 bits A/D converter. The acquisition frequency is 50 kHz. We call response the
time averaged ratio of the voltage amplitudes on input and output piezos, recorded
simultaneously. Voltages in and out are proportional respectively to the displacement
and the constraint (on the piezos). The dimension of the actual response is the inverse
of a stiffness, as what we measure is the ratio of voltages. Dimentional prefactors are
omited for simplicity, as they are constant for the same setup (string and transducers).
The diameter of the string is 100 µm, less than the viscous scale of the flow which
is about η ≃ 170 µm at the largest Re accessible. The equation of motion of the
PIEZOS
STAND
Figure 1. Eperimental setup: the thin steel wire is pulled across a turbulent air
jet by a 4 Kg weight on a rigid stand. Piezoelectric transducers are in mecanical
contact with the wire at each end.
undamped and unforced string is a linear wave equation. Its solutions with fixed ends
are standing waves r(x, t) = A cos(ωn t − knx), where A is the amplitude, t is time
and x is position along the wire. The discrete wave numbers are kn = n
, where L
Measurements of a dynamical temperature in turbulence. 5
is the length of the string and n is a positive integer. In a first approximation, the
waves are not dispersive: ωn = c kn, where c is the phase velocity. T is the tension of
the string and µ its mass per unit length, c =
T/µ ≃ 300 m/s. With a 4 kg weight
on one end, the string’s fundamental frequency is f0 = 344 Hz.
Dissipation is mainly due to friction on air, and causes little dispersion. More precise
treatment would require terms of dissipation in the wire itself and in the piezoelectric
transducers that fix the ends. We neglect this, as the amplitude remains small (a
few tens of micrometers) if compared to the length of the ceramic pile (3mm), or
even the wire diameter (100µm). The possible coupling with compression wave is not
relevant, as the range of frequency is distinct. (Compression wave speed in steel is a
few thousands of m/s, larger than what we consider here: c ≃ 300 m/s.) When this
wire is immersed into the turbulent flow, the resonant modes are excited by the drag
forcing. The quantities measured are averaged along the wire. They are therefore
global in space but local in scale, or more precisely in Fourier-space. The vortices
at scale l are expected to excite modes of wave-number k = 2π/l. In that sense,
the string is acting like a mechanical spectrometer, almost exactly like a Fabry-Perot
interferometer.
3. Measurements
Modulus of the response function is plotted in Fig. 2. It shows that the resonance
peaks are indeed very narrow, ensuring a very precise selection of wave-numbers:
the quality factor is approximately Q ≃ 4000. The imaginary part of the response
function is giving the dissipation. The width of the peaks in the modulus is also
Figure 2. Modulus of the response function versus the harmonic number, at
Re = 154000. The abscissa is given in non-dimensional coordinates, normalised
by the fundamental frequency.
linked to the dissipation, as well as the damping time after a perturbation. We used
in the following the measurement of the imaginary part of the response, but checked
that these different methods coincide. Only the resonant frequencies are considered
in this study, as they are much more sensitive to the velocity fluctuations. This is
Measurements of a dynamical temperature in turbulence. 6
especially important at large k, as the kinetic energy of the flow is small. Spectrum
of the fluctuation excited by the turbulent drag is shown in Fig. 3. Fluctuations
resonance peaks are clearly identified. Spurious vibrations are visible, mainly caused
by the vibrations of the stand. Because the peaks are very thin, long acquisitions
are necessary, as well as large windows for the FFT calculations (150000 points), in
order to achieve a sufficient resolution (0.33Hz). The protocol we used to find the
resonance frequencies, the value of the amplitude of fluctuations, and imaginary part
of the response, is the following. Resonance frequency is obtained by spline smoothing
each peak around the maximum amplitude of the response. Then, imaginary part is
measured after being also smoothed. The amplitude of the fluctuations peaks are
collected on the spectrum, after local smoothing around the maxima. One can see the
Figure 3. Spectrum of the resonance modes of the string excited by turbulent
drag fluctuations, at Re = 154000.
FDR in Fig. 4, called kBTeff., for several values of Re. Uncertainties on this ratio have
multiple origins. Errors indicated by the size of the symbols are those coming from
the determination of the resonance frequencies. Spurious vibrations of the stand are
difficult to handle: we perform measurements of response and fluctuations in the same
conditions, to reduce its influence on the ratio. We believe the scattering of the points
in Fig. 4 comes mainly from the weakening of signal/noise ratio for large frequencies,
simply because there is less energy in the flow at large k, especially at small Re.
The only possible escape on this point is to improve the coupling between the string
and the sensors. The wave-number has been rescaled with the internal viscous scale
η ∝ Re−3/4. The ordinates have been rescaled by an estimated number of degrees
of freedom: (L/η)3 ∝ Re9/4. These Re scalings are both usual consequences from
Kolmogorov’s theory. In other words, the “thermal energy” kBTeff. that the FDR is
representing in the framework of Cugliandolo et al ’s theory, is given per degree of
freedom. Assuming the number of degrees of freedom is the total number of particles
of size η in the total volume is usual, but crude. A more realistic description should
involve correlations between them, reducing this number. However, all the curves
collapse to a single power-law with this scaling. The exponent is discussed in the
Measurements of a dynamical temperature in turbulence. 7
Figure 4. Spectrum of the FDR, labelled as thermal agitation per degree of
freedom. Axis are rescaled with proper Reynolds number dependence, between
74000 and 170000. The size of the symbols represents the uncertainty in the
determination of the maxima of the peaks. The solid line is a k−11/3 power-law
given as an eye guide.
following section.
Please note that the equipartition of energy at equilibrium would require this spectrum
to be constant. There is no equilibrium between the Fourier modes, because of
the energy flux through scales. Moreover, they are not independent, and probably
not Gaussian. There is no reason to expect equipartition. Considering a kinematik
temperature as poportional to the kinetic energy, like in the kinetic theory of gases, it
would be: T ∝ 〈ṽ2〉. And, because of Kolmogorov’s theory it would scale as k−5/3.
The dependance we observe with our definition is much steeper.
4. Scaling law
Because the susceptibility of the string is very high at resonance, the half-wave-length
modes nλ/2 match with velocity structures of scale l (n is an integer). Therefore, the
wave number of the standing wave in the string k = n 2π/λ is the same as k = 2π/l.
The necessary condition for this matching is resonance. It also ensures that velocities
of the string and fluid equalise, which is crucial for the following argument.
Displacement is proportional to the drag forcing, itself proportional to velocity, as
drag is viscous: the string diameter-based Reynolds number is small (about 10).
The Melde string is not dispersive: ω = 2πf = ck, c being the wave velocity.
Therefore, the displacement is r = v/ω = v/(ck), and its power spectrum is:
〈r̃(ω)2〉 = 〈ṽ(ω)2〉(ck)−2 ∝ k−11/3. Because the viscous dissipation at each
resonance is proportional to frequency, the FDR of Eq. 2 is simply proportional to
c k 〈r̃(ω)2〉 ∝ k−11/3. Following Eq. 2, an effective “thermal agitation” defined by the
FDR would be: kBTeff. ∝ k
−11/3, in the inertial range of fully developed turbulence.
This exponent is compatible with the spectrum we measured, as can be seen in Fig.
Measurements of a dynamical temperature in turbulence. 8
5. Discussion
Theoretical characterisation of turbulence in terms of temperature were proposed in
the past by several authors. The temperatures as defined by T. M. Brown [10] and
B. Castaing [11] do not depend on k throughout the inertial range. The qualitative
idea is that the cascade transport process is efficient enough to equalise a quantity
they call temperature. In another model invoking an extremum principle, B. Castaing
proposed a definition of temperature, which might depend on scale [12]. In any case,
none of these theories invoke the FDR. On different basis, R. Robert and J. Sommeria
proposed a definition of temperature [13], only valid for 2D turbulence. It is not
expected to apply in a 3D flow.
Now, let’s consider our experimental results from the perspective of the three points of
reflexion we proposed in the first section, in relation with the FDT. 1- Linear response:
as we mentioned, the coupling between the string and the flow is purely viscous.
Therefore, drag force is proportional to velocity: f(t) = γ v(t), γ being a friction
coefficient. It is also the time-derivative of the position f(t) = γ ω r(t). Response is
linear in r, but the coefficient depends on frequency. 2- Are fluctuations and dissipation
proportional ? As we have seen, the measurements of the FDR are consistent with a
k−11/3 scaling, it is definitely not constant with respect to k. As our system is out
of equilibrium but stationary, there is no time evolution like the relaxation of glasses.
3- Setting a string in a turbulent flow allows to perform measurements on a couple
of conjugate force-displacement variables. We have no other set of observables to
compare with, for now.
We may ask whether what we measure is actually a temperature, in a dynamical
sense. If one assumes that each mode of the string is a harmonic oscillator, and that
a harmonic oscillator at equilibrium with a bath gives the temperature of this bath
through the FDR, then equilibrium between modes of the string and modes of the flow
means the temperature is equal: measurements give the temperature of the flow at this
corresponding scale. Such interpretation still rely on the assumption that FDR on the
oscillator gives the temperature of the oscilaror: this is our working hypothesis. By
equilibrium between modes of the string and the flow, we mean a ’no-flux’ condition
on energy. This is ensured by the high susceptibility of the string at resonance. In
other words, the probe and the reservoir are in equilibrium with each other for each
k, but equilibium is obviously not expected between one scale and another.
We have performed measurements on a turbulent flow, coupling to it a set of harmonic
oscillators: a Melde string. At equilibrium with the flow, in the sense that each mode
of the string couples with the fluid at scale l = πc/ω. It gives informations much
like a spectrometer, even though the flow itself is strongly out of equilibrium. This
is true, of course, as long as the response of the string is fast enough compared to
the frequencies of the velocity fluctuations. The displacement spectra are recorded at
different values of Re, as well as the complex response of the string over an excitation
(contributions of all the standing waves).
The matching of the string’s modes and hydrodynamic structures, what we call
equilibrium between the string and the flow, is still a questionable working hypothesis.
However, drawing inspiration from Cugliandolo et al ’s theory of non-equilibrium
temperature based on the FDR, we measured the Fluctuation over Dissipation Ratio
of our string in a turbulent flow, for different values of Re. The FDR, multiplied
by an appropriate power of the Reynolds number exhibits a unique power law, when
Reynolds number is between 74000 and 170000. The exponent is consistent with a
Measurements of a dynamical temperature in turbulence. 9
value −11/3 given by a very simple model derived from Kolmogorov 1941 theory.
Acknowledgments
We acknowledge B. Castaing, E. Leveque, P. Borgnat, F. Delduc, S. Ciliberto,
E. Bertin, and K. Gawedzki for many discussions. We also thank V. Bergeron,
T. Divoux, and V. Vidal for corrections on the manuscript and for many discussions.
Thanks to F. Dumas for his help in the construction of positioning devices. As this
system became a teaching experiment, several students contributed to this study as
part of their graduate lab-course. They are gratefully acknowledged: A. Louvet,
G. Bordes, I. Dossmann, J. Perret, C. Cohen, and M. Mathieu. We also thank the
guitar maker D. Teyssot, from Lyon, who gently gave us his thinnest E strings.
[1] L.D. Landau and E.M. Lifshitz. Course of Theoretical Physics: Fluid mechanics. Mir, 1971.
[2] A.S. Monin and A.M. Yaglom. Statistical fluid mechanics. MIT Press, Cambridge, 1975.
[3] U. Frisch. Turbulence: the legacy of A.N. Kolmogorov. Cambridge Univ. Press., 1995.
[4] L.F. Richardson. Weather prediction by numerical process. Cambridge Univ. Press, 1922.
[5] A.N. Kolmogorov. C. R. Acad. Sci. U.S.S.R., 30, 1941.
[6] M. Toda R. Kubo and N. Hashitsume. Statistical Physics II: Nonequilibrium Statistical
Mechanics, volume II. Springer, 1985.
[7] L. Cugliandolo and J. Kurchan. Phys. Rev. Lett., 71, 1993.
[8] J. Kurchan L. Cugliandolo and L. Peliti. Phys. Rev. E, 55, 1997.
[9] P. Marcq and A. Naert. Phys. of Fluids, 13, 2001.
[10] T.M. Brown. J. Phys. I, 15, 1982.
[11] B. Castaing. J. Phys. II, 6, 1996.
[12] B. Castaing. J. Phys. II, 50, 1989.
[13] J. Sommeria and R. Robert. J. Fluid Mech., 229, 1991.
Introduction
The Melde string and the experimental setup
Measurements
Scaling law
Discussion
|
0704.0326 | On generalized entropy measures and pathways | ON GENERALIZED ENTROPY MEASURES AND PATHWAYS
A.M. MATHAI
Department of Mathematics and Statistics, McGill University, Montreal,
Canada H3A 2K6, and
Centre for Mathematical Sciences, Pala Campus, Arunapuram P.O., Pala-686
574, Kerala, India
H.J. HAUBOLD
Office for Outer Space Affairs, United Nations, Vienna International Centre,
P.O. Box 500, A-1400 Vienna, Austria and
Centre for Mathematical Sciences, Pala Campus, Arunapuram P.O., Pala-686
574, Kerala, India
Abstract. Product probability property, known in the literature as statistical
independence, is examined first. Then generalized entropies are introduced, all
of which give generalizations to Shannon entropy. It is shown that the nature
of the recursivity postulate automatically determines the logarithmic functional
form for Shannon entropy. Due to the logarithmic nature, Shannon entropy
naturally gives rise to additivity, when applied to situations having product
probability property. It is argued that the natural process is non-additivity,
important, for example, in statistical mechanics (Tsallis 2004, Cohen 2005),
even in product probability property situations and additivity can hold due to
the involvement of a recursivity postulate leading to a logarithmic function.
Generalized entropies are introduced and some of their properties are exam-
ined. Situations are examined where a generalized entropy of order α leads to
pathway models, exponential and power law behavior and related differential
equations. Connection of this entropy to Kerridge’s measure of “inaccuracy” is
also explored.
1. Introduction
Mathai and Rathie (1975) consider various generalizations of Shannon en-
tropy (Shannon, 1948), called entropies of order α, and give various properties,
including additivity property, and characterization theorems. Recently, Mathai
and Haubold (2006, 2006a) explored a generalized entropy of order α, which
is connected to a measure of uncertainty in a probability scheme, Kerridge’s
(Kerridge, 1961) concept of inaccuracy in a scheme, and pathway models that
are considered in this paper.
As defined in Mathai and Haubold (2006, 2006a) the entropy Mk,α(P ) is a
non-additive entropy and his measure M∗k,α(P ) is an additive entropy. It is also
shown that maximization of the continuous analogue of Mk,α(P ), denoted by
Mα(f), gives rise to various functional forms for f , depending upon the types
of constraints on f .
http://arxiv.org/abs/0704.0326v2
Occasionally, emphasis is placed on the fact that Shannon entropy satisfies
the additivity property, leading to extensivity. It will be shown that when
the product probability property (PPP) holds then a logarithmic function can
give a sum and a logarithmic function enters into Shannon entropy due to the
assumption introduced through a certain type of recursivity postulate. The
concept of statistical independence will be examined in Section 1 to illustrate
that simply because of PPP one need not expect additivity to hold or that
one should not expect this PPP should lead to extensivity. The types of non-
extensivity, associated with a number of generalized entropies, are pointed out
even when PPP holds. The nature of non-extensivity that can be expected
from a multivariate distribution, when PPP holds or when there is statistical
independence of the random variables, is illustrated by taking a trivariate case.
Maximum entropy principle is examined in Section 2. It is shown that
optimization of measures of entropies, in the continuous populations, under
selected constraints, leads to various types of models. It is shown that the
generalized entropy of order α is a convenient one to obtain various probability
models.
Section 3 examines the types of differential equations satisfied by the various
special cases of the pathway model.
1.1. Product probability property (PPP) or statistical independence
of events
Let P (A) denote the probability of the event A. If the definition P (A∩B) =
P (A)P (B) is taken as the definition of independence of the events A and B then
any event A ∈ S, and S the sure event are independent. But A is contained in S
and then the definition of independence becomes inconsistent with the common
man’s vision of independence. Even if the trivial cases of the sure event S and
the impossible event φ are deleted, still this definition becomes a resultant of
some properties of positive numbers. Consider a sample space of n distinct
elementary events. If symmetry in the outcomes is assumed then we will assign
equal probabilities 1
each to the elementary events. Let C = A ∩B. If A and
B are independent then P (C) = P (A)P (B). Let
P (A) =
, P (B) =
, P (C) =
⇒ nz = xy, x, y, z = 1, 2, ..., n− 1, z < x, y (1)
deleting S and φ. There is no solution for x, y, z for a large number of n, for
example, n = 3, 5, 7. This means that there are no independent events in such
cases and it sounds strange from a common man’s point of view.
The term “independence” of events is a misnomer. This property should
have been called product probability property or PPP of events. There is no
reason to expect the information or entropy in a joint distribution to be the sum
of the information contents of the marginal distributions when the PPP holds
for the distributions, that is when the joint density or probability function is
a product of the marginal densities or probability functions. We may expect a
term due to the product probability to enter into the expression for the entropy
in the joint distribution in such cases. But if the information or entropy is
defined in terms of a logarithm, then naturally, logarithm of a product being
the sum of logarithms, we can expect a sum coming in such situations. This is
not due to independence or due to the PPP of the densities but due to the fact
that a functional involving logarithm is taken thereby a product has become
a sum. Hence not too much importance should be put on whether or not the
entropy on the joint distribution becomes sum of the entropies on marginal
distributions or additivity property when PPP holds.
1.2. How is logarithm coming in Shannon’s entropy?
Several characterization theorems for Shannon entropy and its various gen-
eralizations are given in Mathai and Rathie (1975. Modified and refined versions
of Shannon’s own postulates are given as postulates for the first theorem charac-
terizing Shannon entropy in Mathai and Rathie (1975). Apart from continuity,
symmetry, zero-indifference and normalization postulates the main postulate
in the theorem is a recursivity postulate, which in essence says that when the
PPP holds then the entropy will be a weighted sum of the entropies, thus in
effect, assuming a logarithmic functional form. The crucial postulate is stated
here. Consider a multinomial population P = (p1, ..., pm), pi > 0, i = 1, ...,m,
p1 + ... + pm = 1, that is, pi = P (Ai), i = 1, ...,m, A1 ∪ ... ∪ Am = S,
Ai ∩ Aj = φ, i 6= j. If any pi can take a zero value also then zero-indifferent
postulate, namely that the entropy remains the same when an impossible event
is incorporated into the scheme, is to be added. Let Hn(p1, ..., pn) denote the
entropy to be defined. Then the crucial recursivity postulate says that
Hn(p1, ..., pm−1, pmq1, .., pmqn−m+1)
= Hm(p1, ..., pm) + pmHn−m+1(q1, ..., qn−m+1) (2)
i=1 pi = 1,
∑n−m+1
i=1 qi = 1. This says that if the m-th event Am is par-
titioned into independent events P (Am ∩ Bj) = P (Am)P (Bj) = pmqj , j =
1, ..., n − m + 1 so that pm = pmq1 + ... + pmqn−m+1 then the entropy Hn(·)
becomes a weighted sum. Naturally, the result will be a logarithmic function
for the measure of entropy.
There are several modifications to this crucial recursivity postulate. One
suggested by Tverberg is that n−m+ 1 = 2 and q1 = q, q2 = 1− q, 0 < q < 1
and H2(q, 1 − q) is assumed to be Lebesgue integrable in 0 ≤ q ≤ 1. Again
a characterization of Shannon entropy is obtained. In all the characterization
theorems for Shannon entropy this recursivity property enters in one form or the
other as a postulate, which in effect implies a logarithmic form for the entropy
measure. Shannon entropy Sk has the following form:
Sk = −A
pi ln pi, pi > 0, i = 1, ..., k, p1 + ...+ pk = 1, (3)
where A is a constant. If any pi is assumed to be zero then 0 ln 0 is to be
interpreted as zero. Since the constant A is present, logarithm can be taken to
any base. Usually the logarithm is taken to the base 2 for ready application to
binary systems. We will take logarithm to the base e.
1.3. Generalization of Shannon entropy
Consider again a multinomial population P = (p1, ..., pk), pi > 0, i =
1, ..., k, p1 + ... + pk = 1. The following are some of the generalizations of
Shannon entropy Sk.
Rk,α(P ) =
i=1 p
, α 6= 1, α > 0, (4)
(Rényi entropy of order α of 1961)
Hk,α(P ) =
i=1 p
i − 1
21−α − 1
, α 6= 1, α > 0 (5)
(Havrda-Charvát entropy of order α of 1967)
Tk,α(P ) =
i=1 p
i − 1
, α 6= 1, α > 0 (6)
(Tsallis entropy of 1988)
Mk,α(P ) =
i=1 p
i − 1
, α 6= 1, −∞ < α < 2 (7)
(entropic form of order α)
M∗k,α(P ) =
i=1 p
, α 6= 1, −∞ < α < 2, (8)
(additive entropic form of order α).
When α → 1 all the entropies of order α described above in (4) to (7) go to
Shannon entropy Sk.
Rk,α(P ) = lim
Hk,α(P ) = lim
Tk,α(P ) = lim
Mk,α(P ) = lim
M∗k,α(P ) = Sk.
Hence all the above measures are called generalized entropies of order α.
Let us examine to see what happens to the above entropies in the case of a
joint distribution. Let pij > 0, i = 1, ...,m, j = 1, ..., n such that
j=1 pij =
1. This is a bivariate situation of a discrete distribution. Then the entropy in
the joint distribution, for example,
Mm,n,α(P,Q) =
j=1 p
ij − 1
. (10)
If the PPP holds and if pij = piqj , p1 + ... + pm = 1, q1 + ... + qn = 1,
pi > 0, i = 1, ...,m, qj > 0, j = 1, ..., n and if P = (p1, ..., pm), Q = (q1, ..., qn)
(α− 1)Mm,α (P ) Mn,α(Q) =
i − 1
j − 1
j + 1
= Mm,n,α(P,Q) −Mm,α(P )−Mn,α(Q).
Therefore
Mm,n,α(P,Q) = Mm,α(P ) +Mn,α(Q) + (α− 1)Mm,α(P )Mn,α(Q). (11)
If any one of the above mentioned generalized entropies in (4) to (8) is written
as Fm,n,α(P,Q) then we have the relation
Fm,n,α(P,Q) = Fm,α(P ) + Fn,α(Q) + a(α)Fm,α(P )Fn,α(Q). (12)
where
a(α) = 0 (Rényi entropy Rk,α(P ))
= 21−α − 1 (Havrda-Charvát entropy Hk,α(P ))
= 1− α (Tsallis entropy Tk,α(P ))
= α− 1 (entropic form of order α, i.e., Mk,α(P ))
= 0 (additive entropic form of order α, i.e., M∗k,α(P )). (13)
When a(α) = 0 the entropy is called additive and when a(α) 6= 0 the entropy
is called non-additive. As can be expected, when a logarithmic function is
involved, as in the cases of Sk(P ), Rk,α(P ),M
k,α(P ), the entropy is additive
and a(α) = 0.
1.4. Extensions to higher dimensional joint distributions
Consider a trivariate population or a trivariate discrete distribution pijk >
0, i = 1, ...,m, j = 1, ..., n, k = 1, ..., r such that
k=1 pijk = 1. If
the PPP holds mutually, that is, pair-wise as well as jointly, which then will
imply that
pijk = piqjsk,
pi = 1,
qj = 1,
sk = 1,
P = (p1, ..., pm), Q = (q1, ..., qn), S = (s1, ..., sr).
Then proceeding as before, we have for any of the measures described above in
(4) to (8), calling it F (·),
Fm,n,r,α(P,Q, S) = Fm,α(P ) + Fn,α(Q) + Fr,α(S) + a(α)[Fm,α(P )Fn,α(Q)
+Fm,α(P )Fr,α(S) + Fn,α(Q)Fr,α(S)]
+[a(α)]2Fm,α(P )Fn,α(Q)Fr,α(S) (14)
where a(α) is the same as in (13). The same procedure can be extended to any
multivariable situation. If a(α) = 0 we may call the entropy additive and if
a(α) 6= 0 then the entropy is non-additive.
1.5. Crucial recursivity postulate
Consider the multinomial population P = (p1, ..., pk), pi > 0, i = 1, ..., k, p1+
... + pk = 1. Let the entropy measure to be determined through appropriate
postulates be denoted by Hk(P ) = Hk(p1, ..., pk). For k = 2 let
f(x) = H2(x, 1− x), 0 ≤ x ≤ 1 or x ∈ [0, 1]. (15)
If another parameter α is to be involved in H2(x, 1−x) then we will denote f(x)
by fα(x). From (5) to (7) it can be seen that the generalized entropies of order
α of Havrda-Charvát (1967), Tsallis (1988, 2004) and Shannon (1948) entropy
satisfy the functional equation
fα(x) + bα(x)fα
= fα(y) + bα(x)f
for x, y ∈ [0, ) with x+ y ∈ [0, 1], with the boundary condition
fα(0) = fα(1) (17)
where
bα(x) = 1− x (Shannon entropy Sk(P ))
= (1− x)α (Harvda-Charvát entropy Hk,α(P ))
= (1− x)α (Tsallis entropy Tk,α(P ))
= (1− x)2−α (entropic form of order α, i.e., Mk,α(P )). (18)
Observe that the normalizing constant at x = 1
is equal to 1 for Hk,α(P ) and it
is different for other entropies. Thus equations (6),(7),(8), with the appropriate
normalizing constants fα(
), can give characterization theorems for the various
entropy measures. The form of bα(x) is coming from the crucial recursivity
postulate, assumed as a desirable property for the measures.
1.6. Continuous analogues
In the continuous case let f(x) be the density function of a real random
variable x. Then the various entropy measures, corresponding to the ones in (4)
to (8) are the following:
Rα(f) =
[f(x)]αdx
, α 6= 1, α > 0 (19)
(Rényi entropy of order α)
Hα(f) =
21−α − 1
[f(x)]αdx− 1
, α 6= 1, α > 0 (20)
(Havrda-Charvát entropy of order α)
Tα(f) =
[f(x)]αdx− 1
, α 6= 1, α > 0, (21)
(Tsallis entropy of order α)
Mα(f) =
[f(x)]2−αdx− 1
, α 6= 1, α < 2 (22)
(entropic form of order α)
M∗α(f) =
[f(x)]2−αdx
, α 6= 1, α < 2 (23)
(additive entropic form of order α).
As expected, Shannon entropy in this case is given by
S(f) = −A
f(x) ln f(x)dx (24)
where A is a constant.
Note that when PPP (product probability property) or statistical indepen-
dence holds then in the continuous case also we have the property in (12) and
(14) and then non-additivity holds for the measures analogous to the ones in
(3),(5),(6),(7) with a(α) remaining the same. Since the steps are parallel a
separate derivation is not given here.
2. Maximum Entropy Principle
If we have a multinomial population P = (p1, ..., pk), pi > 0, i = 1, ..., k, p1+
...+ pk = 1 or the scheme P (Ai) = pi, A1 ∪ ... ∪ Ak = S, P (S) = 1, Ai ∩ Aj =
φ, i 6= j then we know that the maximum uncertainty in the scheme or the
minimum information from the scheme is obtained when we cannot give any
preference to the occurrence of any particular event or when the events are
equally likely or when p1 = p2 = ... = pk =
. In this case, Shannon entropy
becomes,
Sk(P ) = Sk(
, ...,
) = −A
= A ln k (25)
and this is the maximum uncertainty or maximum Shannon entropy in this
scheme. If the arbitrary functional f is to be fixed by maximizing the entropy
then in (19) to (21) we have to optimize
[f(x)]αdx for fixed α, over all
functional f , subject to the condition
f(x)dx = 1 and f(x) ≥ 0 for all x.
For applying calculus of variation procedure we consider the functional
U = [f(x)]α − λ[f(x)]
where λ is a Lagrangian multiplier. Then the Euler equation is the following:
= 0 ⇒ αfα−1 − λ = 0 ⇒ f =
= constant. (26)
Hence f is the uniform density in this case, analogous to the equally likely
situation in the multinomial case. If the first moment E(x) =
xf(x)dx
is assumed to be a given quantity for all functional f then U will become the
following for (19) to (21).
U = [f(x)]α − λ1[f(x)]− λ2xf(x)
and the Euler equation leads to the power law. That is,
= 0 ⇒ αfα−1 − λ1 − λ2x = 0 ⇒ f = c1
. (27)
By selecting c1, λ1, λ2 appropriately we can create a density out of (27). For
α > 1 and λ2
> 0 the right side in (27) increases exponentially. If α = q > 1 and
= q − 1 then we have Tsallis’ q-exponential function from the right side of
(27). If α > 1 and λ2
= −(α−1) then (27) can produce a density in the category
of a type-1 beta. From (27) it is seen that the form of the entropies of Havrda-
CharvátHk,α(P ) and Tsallis Tk,α(P ) need special attention to produce densities
(Ferri et al. 2005). However, Tsallis has considered a different constraint on
E(x). If the density f(x) is replaced by its escort density, namely, µ[f(x)]α
where µ−1 =
[f(x)]αdx and if the expected value of x in this escort density
is assumed to be fixed for all functional f then the U of (26) becomes
U = fα − λ1f + µλ2xf
= 0 ⇒ αfα−1[1 + µλ2x] = λ1 ⇒ f =
(1+λ3x)
f = λ∗1[1 + λ3x]
where λ3 is a constant and λ
1 is the normalizing constant. If λ3 is taken as
λ3 = α− 1 then
f = λ∗1[1 + (α− 1)x]
α−1 . (28)
Then (28) for α > 1 is Tsallis statistics (Tsallis 2004, Cohen 2005). Then for
α < 1 also by writing α − 1 = −(1 − α) one gets the case of Tsallis statistics
for α < 1 (Ferri et al. 2005). These modifications and the consideration of
escort distribution are not necessary if we take the generalized entropy of order
α. Thus if we consider Mα(f) and if we assume that the first moment in f(x)
itself is fixed for all functional f then the Euler equation gives
(2− α)f1−α − λ1 + λ2x = 0 ⇒ f = λ̄
and for λ2
= 1− α we have Tsallis statistics (Tsallis 2004, Cohen 2005)
f = λ̄[1− (1− α)x]
1−α (29)
coming directly, where λ̄ is the normalizing constant.
Let us start with Mα(f) of (20) under the assumptions that f(x) ≥ 0 for all
f(x)dx = 1,
xδf(x)dx is fixed for all functional f and for a specified
δ > 0, f(a) is the same for all functional f , f(b) is the same for all functional
f , for some limits a and b, then the Euler equation becomes
(2 − α)f1−α − λ1 − λ2x
δ = 0 ⇒ f = c1[1 + c
1−α . (30)
If c∗1 is written as −s(1− α), s > 0 then we have, writing f1 for f ,
f1 = c1[1− s(1 − α)x
1−α , δ > 0, α < 1, 0 ≤ x ≤
[s(1− α)]
where 1 − s(1 − α)xδ > 0. For α < 1 or −∞ < α < 1 the right side of (31)
remains as a generalized type-1 beta model with the corresponding normalizing
constant c1. For α > 1, writing 1 − α = −(α − 1) the model in (31) goes to a
generalized type-2 beta form, namely,
f2 = c2[1 + s(α− 1)x
α−1 . (32)
When α → 1 in (31) or in (32) we have an extended or stretched exponential
form,
f3 = c3e
. (33)
If c∗1 in (30) is taken as positive then (30) for α < 1, α > 1, α → 1 will be
increasing exponentially. Hence all possible forms are available from (30). The
model in (31) is a special case of the distributional pathway model and for a
discussion of the matrix-variate pathway model see Mathai (2005). Special cases
of (31) and (32) for δ = 1 are Tsallis statistics (Gell-Mann and Tsallis, 2004;
Ferri et al. 2005).
Instead of optimizing Mα(f) of (22) under the conditions that f(x) ≥ 0
for all x,
f(x)dx = 1 and
xδf(x)dx is fixed, let us optimize under the
following conditions: f(x) ≥ 0 for all x,
f(x)dx < ∞ and the following two
moment-like expressions are fixed quantities for all functional f ,
x(γ−1)(1−α)f(x)dx = fixed ,
x(γ−1)(1−α)+δf(x)dx = fixed.
Then the Euler equation becomes
(2− α)f1−α −λ1x
(γ−1)(1−α)
− λ2x
(γ−1)(1−α)+δ = 0 ⇒
f = c xγ−1[1 + c∗xδ]
and for c∗ = −s(1 − α), s > 0, we have the distributional pathway model for
the real scalar case, namely
f(x) = c xγ−1[1− s(1− α)xδ ]
1−α , δ > 0, s > 0 (34)
where c is the normalizing constant. For α < 1, (34) gives a generalized type-1
beta form, for α > 1 it gives a generalized type-2 beta form and for α → 1
we have a generalized gamma form. For α > 1, (34) gives the superstatistics
of Beck (2006) and Beck and Cohen (2003). For γ = 1, δ = 1, (34) gives
Tsallis statistics (Tsallis 2004, Cohen 2005). Densities appearing in a number
of physical problems are seen to be special cases of (34), a discussion of which
may be seen from Mathai and Haubold (2006a). For example, (34) for δ =
2, γ = 3, α → 1, x > 0 is the Maxwell-Boltzmann density; for δ = 2, γ = 1, α →
1,−∞ < x < ∞ is the Gaussian density; for γ = δ, α → 1 is the Weibull density.
For γ = 1, δ = 2, 1 < q < 3 we have the Wigner function W (p) giving the atomic
moment distribution in the framework of Fokker-Planck equation, see Douglas,
Bergamini, and Renzoni (2006) where
W (p) = z−1q [1− β(1 − q)p
1−q , 1 < q < 3. (35)
Before closing this section we may observe one more property for Mα(f). As
an expected value
Mα(f) =
E[f(x)]1−α − 1
. (36)
But Kerridge’s (Kerridge, 1961) measure of “inaccuracy” in assigning q(x) for
the true density f(x), in the generalized form is
Hα(f : q) =
(21−α − 1)
E[q(x)]α−1 − 1
, (37)
which is also connected to the measure of directed divergence between q(x) and
f(x). In (37) the normalizing constant is 21−α−1, the same factor appearing in
Havrda-Charvt́ entropy. With different normalizing constants, as seen before,
(36) and (37) have the same forms as an expected value with q(x) replaced
by f(x) in (36). Hence Mα(f) can also be looked upon as a type of directed
divergence or “inaccuracy” measure.
3. Differential Equations
The functional part in (34), for a more general exponent, namely
g(x) =
= xγ−1[1− s(1 − α)xδ]
1−α , α 6= 1, δ > 0, β > 0, s > 0 (38)
is seen to satisfy the following differential equation for γ 6= 1 which defines the
differential pathway.
g(x) = (γ − 1)xγ−1[1− s(1− α)xδ]
−sβδxδ+γ−1[1− s(1− α)xδ]
(1−α)
. (39)
Then for δ =
(γ−1)(α−1)
, γ 6= 1, α > 1 we have
g(x) = (γ − 1)g(x)− sβδ[g(x)]1−
(1−α)
β (40)
= (γ − 1)g(x)− sδ[g(x)]α (41)
for β = 1, γ 6= 1, δ = (γ − 1)(α− 1), α > 1.
For γ = 1, δ = 1 in (38) we have
g(x) = −s[g(x)]η, η = 1−
(1 − α)
= −s[g(x)]α for β = 1. (43)
Here (43) is the power law coming from Tsallis statistics (Gell-Mann and Tsallis,
2004).
Acknowledgement The authors would like to thank the Department of Science
and Technology, Government of India, New Delhi, for the financial assistance for
this work under project No. SR/S4/MS:287/05 which enabled this collaboration
possible.
4. References
Beck, C. (2006). Stretched exponentials from superstatistics. Physica A, 365,
96-101.
Beck, C. and Cohen, E.G.D. (2003). Superstatistics. Physica A, 322, 267-275.
Cohen, E.G.D. (2005). Boltzmann and Einstein: Statistics and dynamics - An
unsolved problem. Pramana, 64, 635-643.
Douglas, P., Bergamini, S., and Renzoni, F. (2006). Tunable Tsallis distribution
in dissipative optical lattices. Physical Review Letters, 96, 110601-1-4.
Ferri, G.L., Martinez, S., and Plastino, A. (2005). Equivalence of the four
versions of Tsallis’s statistics. Journal of Statistical Mechanics: Theory and
Experiment, PO4009.
Gell-Mann, M. and Tsallis, C. (Eds.) (2004). Nonextensive Statistical Mechan-
ics: Interdisciplinary Applications. Oxford University Press, Oxford.
Havrda, J. and Charvát, F. (1967). Quantification method of classification pro-
cedures: Concept of structural α-entropy. Kybernetika, 3, 30-35.
Kerridge, D.F. (1961). Inaccuracy and inference. Journal of the Royal Statisti-
cal Society Series B, 23, 184-194.
Mathai, A.M. (2005). A pathway to matrix-variate gamma and normal densi-
ties. Linear Algebra and Its Applications, 396, 317-328.
Mathai, A.M. and Haubold, H.J. (2006). Pathway model, Tsallis statistics, su-
perstatistics and a generalized measure of entropy. Physica A , 375), 110-122.
Mathai,A.M. and Haubold, H.J. (2006a). On generalized distributions and path-
ways. arXiv:cond-mat/0609526v2.
Mathai, A.M. and Rathie, P.N. (1975). Basic Concepts in Information Theory
and Statistics: Axiomatic Foundations and Applications, Wiley Halstead, New
York and Wiley Eastern, New Delhi.
Rényi, A. (1961). On measure of entropy and information. Proceedings of the
Fourth Berkeley Symposium on Mathematical Statistics and Probability, 1960,
University of California Press, 1961, Vol. 1, 547-561.
Shannon, C.E. (1948). A mathematical theory of communication. Bell System
Technical Journal, 27, 379-423, 547-561.
Tsallis, C. (1988). Possible generalization of Boltzmann-Gibbs statistics. Jour-
nal of Statistical Physics, 52, 479-487.
Tsallis, C. (2004). What should a statistical mechanics satisfy to reflect nature?,
Physica D, 193, 3-34.
http://arxiv.org/abs/cond-mat/0609526
|
0704.0327 | Evolution of a band insulating phase from a correlated metallic phase | Evolution of a band insulating phase from a correlated metallic phase
Kalobaran Maiti,∗ Ravi Shankar Singh, and V.R.R. Medicherla
Department of Condensed Matter Physics and Materials’ Science,
Tata Institute of Fundamental Research, Homi Bhabha Road, Colaba, Mumbai - 400 005, INDIA
(Dated: October 30, 2018)
We investigate the evolution of the electronic structure in SrRu1−xTixO3 as a function of x using
high resolution photoemission spectroscopy, where SrRuO3 is a weakly correlated metal and SrTiO3
is a band insulator. The surface spectra exhibit a metal-insulator transition at x = 0.5 by opening
up a soft gap. A hard gap appears at higher x values consistent with the transport properties. In
contrast, the bulk spectra reveal a pseudogap at the Fermi level, and unusual evolution exhibiting
an apparent broadening of the coherent feature and subsequent decrease in intensity of the lower
Hubbard band with the increase in x. Interestingly, the first principle approaches are found to be
sufficient to capture anomalous evolutions at high energy scale. Analysis of the spectral lineshape
indicates strong interplay between disorder and electron correlation in the electronic properties of
this system.
PACS numbers: 71.10.Hf, 71.20.-b, 71.30.+h
The investigation of the role of electron correlation in
various electronic properties is a paradigmatic problem
in solid state physics. Numerous experimental and the-
oretical studies are being performed on correlated elec-
tron systems revealing exotic phenomena such as high
temperature superconductivity, giant magnetoresistance
etc. Electron correlation essentially localizes the valence
electrons leading the system towards insulating phase.
Correlation induced insulators, known as Mott insulators
are characterized by a gapped electronic excitations in a
system where effective single particle approaches provide
a metallic ground state. The band insulators represent
insulating phase described within the single particle ap-
proaches. Strikingly, some recent theoretical studies re-
veal a correlation induced metallic ground state in a band
insulator using ionic Hubbard model [1, 2, 3, 4]. Such un-
usual transition has been observed in two dimensions by
tuning effective electron correlation strength, U/W (U
= electron-electron Coulomb repulsion strength, W =
bandwidth) and the local potential, ∆.
In order to realize such effect experimentally, we in-
vestigate the evolution of the electronic structure in
SrRu1−xTixO3 as a function of x, where the end mem-
bers, SrRuO3 and SrTiO3 are correlated ferromagnetic
metal and band insulator, respectively. Ti remains in
tetravalent state in the whole composition range having
no electron in the 3d band[5, 6]. Thus, in addition to
the introduction of disorder in the Ru-O sublattice, Ti-
substitution at the Ru-sites dilutes Ru-O-Ru connectiv-
ity leading to a reduction in Ru 4d bandwidth, W and
hence, U/W will increase. Transport measurements[7]
exhibit plethora of novel phases such as correlated metal
(x ∼ 0.0 ), disordered metal (x ∼ 0.3), Anderson insu-
lator (x ∼ 0.5), soft Coulomb gap insulator (x ∼ 0.6),
disordered correlated insulator (x ∼ 0.8), and band insu-
lator (x = 1.0).
In this study, we have used high resolution photoemis-
sion spectroscopy to probe the density function in the
vicinity of the Fermi level, ǫF and at higher energy scale
as well. Considering the fact that escape depth of the
photoelectrons is small, we have extracted the surface
and bulk spectra in every case by varying the surface
sensitivity of the technique. The surface spectra exhibit
signature of disorder at lower x values in SrRu1−xTixO3,
a metal-insulator transition exhibiting a soft gap at ǫF
for x = 0.5 and a hard gap for higher x. The bulk spec-
tra, on the other hand, reveal an unusual spectral weight
transfer and signature of a pseudogap at ǫF at higher x.
Photoemission measurements were performed using
Gammadata Scienta analyzer, SES2002 and monochro-
matized photon sources. The energy resolution for x-
ray photoemission (XP) and He II photoemission mea-
surements were set at 300 meV and 4 meV, respectively.
High quality samples of SrRu1−xTixO3 with large grain
size were prepared following solid state reaction route
using high purity ingredients[8] followed by a long sin-
tering (for about 72 hours) at the final preparation tem-
perature. Sharp x-ray diffraction patterns reveal single
phase in each composition with no signature of impurity
feature. Magnetic measurements using a high sensitiv-
ity vibrating sample magnetometer exhibit distinct fer-
romagnetic transition at each x up to x = 0.6 studied,
as also evidenced by the Curie-Weiss fits in the param-
agnetic region. The fits provide an estimation of effec-
tive magnetic moment (µ = 2.8 µB , 2.54 µB, 2.45 µB,
2.18 µB, 2.19 µB, 1.95 µB and 1.93 µB) and Curie tem-
perature (θP = 164 K, 156.6 K, 150.6 K, 145.3 K, 139 K,
138.6 K and 100 K) for x = 0.0, 0.15, 0.2, 0.3, 0.4, 0.5
and 0.6, respectively. The values of µ and θP for SrRuO3
are observed to be the largest among those available in
the literature and corresponds to well characterized sin-
gle crystalline materials[9].
In Fig. 1(a), we show the XP valence band spectra
exhibiting 4 distinct features marked by A, B, C and
D. The features C and D appear beyond 2.5 eV and
have large O 2p character as confirmed experimentally
http://arxiv.org/abs/0704.0327v1
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FIG. 1: (color online) (a) XP valence band spectra of
SrRu1−xTixO3 for various values of x. Solid line represents
the O 2p part for x = 0.6. (b) Ru 4d spectra after the sub-
traction of the O 2p contributions as shown in (a). (c) Ru 4d
band obtained from He II spectra.
by changing photoemission cross-sections [10] and theo-
retically by band structure calculations [11]. The peaks
A and B appear primarily due to the photoemission from
electronic states having Ru 4d character. The O 2p part
remains almost the same in the whole composition range
as expected. While Ru 4d intensity gradually diminishes
with the decrease in Ru-concentrations, the lineshape of
Ru 4d band exhibits significant redistribution in inten-
sity. In order to bring out the clarity, we delineate the
Ru 4d band by subtracting O 2p contributions. The sub-
tracted spectra, normalized by integrated intensity un-
der the curve, exhibit two distinct features as evident in
Fig. 1(b). The feature A corresponds to the delocalized
electronic density of states (DOS) observed in ab initio
results and is termed as coherent feature. The feature B,
absent in the ab initio results[11], is often attributed to
the signature of correlation induced localized electronic
states forming the lower Hubbard band and is known
as incoherent feature. The increase in x leads to a de-
crease in intensity of A and subsequently, the intensity
of B grows gradually. Since the bulk sensitivity of valence
electrons at 1486.6 eV photon energy is high (∼ 60%),
the spectral evolution in Fig. 1(b) manifests primarily
the changes in the bulk electronic structure.
In order to discuss the effect due to the surface elec-
tronic structure, we show the Ru 4d contributions ex-
tracted from the He II spectra in Fig. 1(c), where the
surface sensitivity is about 80%. Interestingly, all the
spectra are dominated by the peak at higher binding en-
ergies (> 1 eV) corresponding to the surface electronic
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FIG. 2: (color online) S(ǫ) obtained from (a) He II and (b)
XP spectra of SrRu1−xTixO3. (b) S(ǫ) in (a) are plotted as
a function of (c) |ǫ− ǫF |
0.5 (d) and |ǫ− ǫF |
1.25. S(ǫ) obtained
from (e) XP and (f) He II spectra of Ca1−xSrxRuO3.
structure as reported in the case of SrRuO3 and the co-
herent feature intensity corresponds essentially to the
bulk electronic structure[10, 12]. The coherent feature
intensity reduces drastically with the increase in x and
becomes almost negligible at x = 0.6. This can be vi-
sualized clearly in the spectral density of states (SDOS)
obtained by symmetrizing (S(ǫ) = I(ǫ) + I(−ǫ); I(ǫ) =
photoemission spectra, ǫ = binding energy) the He II and
XP spectra. The SDOS corresponding to He II spectrum
of SrRuO3 shown in Fig. 2(a) exhibits a sharp dip at ǫF ,
which increases gradually with the increase in x. The
SDOS corresponding to XP spectra in Fig. 2(b), how-
ever, exhibits a peak in SrRuO3 presumably due to large
resolution broadening and intense coherent feature. This
peak loses its intensity and becomes almost flat for x =
0.15 and 0.2. Further increase in x leads to a pseudogap
at ǫF , which gradually increases with the increase in x.
Both these results clearly indicate gradual depletion of
SDOS at ǫF with the increase in Ti-substitution.
The effect of resolution broadening of 4 meV in the He
II spectra is not significant in the energy scale shown in
the figure. The electron and hole lifetime broadening is
also negligible in the vicinity of ǫF . Thus, S(ǫ) in Fig.
2(a) provide a good testing ground to investigate evolu-
tion of the spectral lineshape at ǫF . The lineshape of
S(ǫ) in Fig. 2(a) exhibits significant modification with
the increase in x. We, thus, replot S(ǫ) as a function
of |ǫ − ǫF |
α for various values of α. Two extremal cases
representing α = 0.5 and 1.25 are shown in Fig. 2(c) and
2(d), respectively. It is evident that S(ǫ) of SrRuO3 ex-
hibit a straight line behavior in Fig. 2(c) suggesting sig-
nificant role of disorder in the electronic structure. The
influence of disorder can also be verified by substitutions
at the A-sites in the ABO3 structure. This has been ver-
ified by plotting SDOS obtained from the XP and He II
spectra of Ca1−xSrxRuO3 in Fig. 2(e) and 2(f), respec-
tively. Here, the electronic properties of the end mem-
bers, SrRuO3 and CaRuO3 are known to be strongly in-
fluenced by the disorder[13]. Substitution of Sr at the
Ca-sites is expected to enhance the disorder effect. The
lineshape of S(ǫ) in both Fig. 2(e) and 2(f) remains al-
most the same across the whole composition range. Such
disorder induced spectral dependence is consistent with
the observations in other systems[14, 15] as well.
Interestingly, the lineshape modifies significantly with
the increase in x and becomes 1.25 in the 60% Ti substi-
tuted sample. Ti substitution introduces defects in the
Ru-O network, where Ti4+ having no d-electron, does
not contribute in the valence band. Thus, in addition to
the disorder effects, the reduced degree of Ru-O-Ru con-
nectivity leads to a decrease in bandwidth, W , which in
turn enhances U/W . In systems consisting of localized
electronic states in the vicinity of ǫF , a soft Coulomb gap
opens up due to electron-electron Coulomb repulsion; in
such a situation, the ground state is stable with respect to
single-particle excitations, when SDOS is characterized
by (ǫ − ǫF )
2-dependence [16, 17]. Here, gradual increase
in α with the increase in x in the intermediate compo-
sitions is curious and indicates strong interplay between
correlation effect and disorder in this system.
The extraction of surface and bulk spectra requires
both the XP and He II spectra collected at significantly
different surface sensitivities. Thus, we broaden the He II
spectra upto 300 meV and extract the surface and bulk
spectra analytically using the same parameters as used
before for CaSrRuO3 system[10]. The surface spectra
shown in Fig. 3(b) exhibit a gradual decrease in coherent
feature intensity with the increase in x and subsequently,
the feature around 1.5 eV becomes intense, narrower and
slightly shifted towards higher binding energies. The de-
crease in intensity at ǫF is clearly visible in the sym-
metrized spectra, S(ǫ) shown in Fig. 3(d). Interestingly,
S(ǫ) of x = 0.5 sample exhibits a soft gap at ǫF and
a hard gap appears in S(ǫ) corresponding to higher x.
This spectral evolution is remarkably consistent with the
transport properties[7]. These results corresponding to 2-
dimensional surface states presumably have strong impli-
cation in realizing recent theoretical predictions[1, 2, 3, 4]
and the bulk properties of this system.
The picture is strikingly different in the bulk spectra
where the electronic structure is 3-dimensional. The bulk
spectrum of SrRuO3 exhibits an intense and sharp coher-
ent feature in the vicinity of ǫF and the incoherent fea-
ture appears around 2 eV. The enhancement of U/W due
to Ti substitution is expected to increase the incoherent
feature intensity. In sharp contrast, the intensity of the
+ , - . / 0 1 2
345 678 9:; <=>?
OPQ RST UVW XYZ[
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úûüýþÿ� ������ ����
FIG. 3: (color online) Extracted (a) bulk and (b) surface
spectra of SrRu1−xTixO3 for various values of x. The SDOS
obtained from bulk and surface spectra are shown in (c) and
(d), respectively.
2 eV feature reduces significantly and the coherent fea-
ture becomes broad. In addition, the bulk spectra of all
the intermediate compositions appear very similar. The
symmetrized bulk spectra shown in Fig. 3(c) exhibit a
small lowering of intensity at ǫF with the increase in x.
Since, U is weak in these highly extended 4d
systems[10, 12], a perturbative approach may be use-
ful to understand the role of electron correlation in the
spectral lineshape. We have calculated the bare density
of states (DOS) for SrRuO3 and SrRu0.5Ti0.5O3 using
state-of-the-art full potential linearized augmented plane
wave method[11, 18]. The self energy and spectral func-
tions were calculated using this t2g partial DOS as done
before[19]. The real and imaginary parts of the self en-
ergy are shown in Fig. 4(a) and 4(b), and the spectral
functions for different U values are shown in Fig. 4(c) and
4(d) for SrRuO3 and SrRu0.5Ti0.5O3, respectively. The
increase in U leads to a spectral weight transfer outside
the LDA DOS width creating the lower and upper Hub-
bard bands. Subsequently, the total width of the LDA
DOS diminishes gradually. While these results exhibit
similar scenario as that observed in the most sophisti-
cated calculations using dynamical mean field theory, the
separation between the lower and upper Hubbard bands
is significantly larger than the corresponding values of U .
It is important to note here that the band structure cal-
culations include the electron-electron interaction term
within the local density approximations. The perturba-
tion calculations in the present case essentially provide
an estimation of the correction in U already included in
the effective single particle Hamiltonian.
In order to compare with the experimental spectra, the
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FIG. 4: (color online) Real and imaginary parts of the self
energy of (a) SrRuO3 and (b) SrRu0.5Ti0.5O3 obtained by
second order perturbation method following the method of
Treglia et al.[19]. Spectral functions for various values of U of
(c) SrRuO3 and (d) SrRu0.5Ti0.5O3. Calculated experimen-
tal spectra for different values of U of (e) SrRuO3 and (f)
SrRu0.5Ti0.5O3.
calculated spectral functions are convoluted by Fermi-
Dirac distribution function and the gaussian representing
the resolution broadening of 300 meV. The comparison is
shown in Figs. 4(e) and 4(f). Interestingly, the spectral
shape corresponding to U = 0.6 ± 0.1 exhibits remark-
able representation of the experimental bulk spectra in
both the cases. These results clearly establish that per-
turbative approaches and local description of the corre-
lation effects are sufficient to capture electronic structure
of these weakly correlated systems. The overall narrow-
ing of the valence band observed in the substituted com-
pounds are essentially a single particle effect and can be
attributed to the reduced degree of Ru-O-Ru connectiv-
ity in these systems. While the high energy scale features
are reproduced remarkably well within this picture, the
occurrence of a pseudogap at ǫF with increasing x (not
visible in Fig.4 due to large energy scale) suggests in-
creasing role of disorder.
In summary, the high resolution spectra of SrRuO3 ex-
hibit signature of disorder in the vicinity of the Fermi
level. Introduction of the Ti4+ sublattice within the
Ru4+ sublattice provides a paradigmatic example, where
the charge density near Ti4+ sites is close to zero and each
Ru4+ site contributes 4 electrons in the valence band.
Such large charge fluctuation leads to a significant change
in spectral lineshape and a dip appears at ǫF (pseudo-
gap). Interestingly, the effects are much stronger in the
two dimensional (surface) electronic structure leading to
a soft gap at 50% substitution and eventually a hard gap
appears. Bulk electronic structure (3-dimensional), how-
ever, remains less influenced. A theoretical understand-
ing of these effects needs consideration of strong disorder
in addition to the electron correlation effects.
∗ Corresponding author: [email protected]
[1] A. Fuhrmann, D. Heilmann, and H. Monien, Phys. Rev.
B 73, 245118 (2006).
[2] S.S. Kancharla and E. Dagotto, Phys. Rev. Lett. 98,
016402 (2007).
[3] Arti Garg, H.R. Krishnamurthy, and Mohit Randeria,
Phys. Rev. Lett. 97, 046403 (2006).
[4] N. Paris, K. Bouadim, F. Hebert, G.G. Batrouni, and
R.T. Scalettar, Phys. Rev. Lett. 98, 046403 (2007).
[5] J. Kim, J.-Y. Kim, B.-G. Park, and S.-J. Oh, Phys. Rev.
B 73, 235109 (2006), M. Abbate, J.A. Guevara, S.L.
Cuffini, Y.P. Mascarenhas, and E. Morikawa, Eur. Phys.
J. B 25, 203 (2002).
[6] S. Ray, D.D. Sarma, and R. Vijayaraghavan, Phys. Rev.
B 73, 165105 (2006).
[7] K.W. Kim, J.S. Lee, T.W. Noh, S.R. Lee, and K. Char,
Phys. Rev. B 71, 125104 (2005).
[8] R.S. Singh and K. Maiti, Solid State Commun, 140, 188
(2006).
[9] G. Cao, S. McCall, M. Shepard, J.E. Crow, and R.P.
Guertin, Phys. Rev. B 56, 321 (1997).
[10] K. Maiti and R.S. Singh, Phys. Rev. B 71, 161102(R)
(2005).
[11] K. Maiti, Phys. Rev. B 73, 235110 (2006).
[12] M. Takizawa, D. Toyota, H. Wadati, A. Chikamatsu, H.
Kumigashira, A. Fujimori, M. Oshima, Z. Fang, M. Lipp-
maa, M. Kawasaki, and H. Koinuma, Phys. Rev. B 72,
060404(R) (2005).
[13] K. Maiti, R.S. Singh, and V.R.R. Medicherla, Europhys.
Lett. (in print); Condmat/0604648.
[14] B.L. Altshuler and A.G. Aronov, Solid State Commun.
30, 115 (1979).
[15] D.D. Sarma et al., Phys. Rev. Lett. 80, 4004 (1998).
[16] A.L. Efros and B.I. Shklovskii, J. Phys. C: Solid State
Phys. 8, L49 (1975).
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(1995).
[18] P. Blaha, K. Schwarz, G.K.H. Madsen, D. Kvasnicka, and
J. Luitz, WIEN2k, An Augmented Plane Wave + Lo-
cal Orbitals Program for Calculating Crystal Properties
(Karlheinz Schwarz, Techn. Universität Wien, Austria),
2001. ISBN 3-9501031-1-2.
[19] G. Treglia et. al., J. Physique 41, 281 (1980); ibid, Phys.
Rev. B 21, 3729 (1980); D.D. Sarma et al., Phys. Rev.
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|
0704.0328 | Electroweak phase transitions in the MSSM with an extra $U(1)'$ | Electroweak phase transitions in the MSSM with an
extra U (1)′
S.W. Ham(1), E.J. Yoo(2), and S.K. Oh(1,2)
(1) Center for High Energy Physics, Kyungpook National University,
Daegu 702-701, Korea
(2) Department of Physics, Konkuk University, Seoul 143-701, Korea
Abstract
We investigate the possibility of electroweak phase transition in the minimal
supersymmetric standard model (MSSM) with an extra U(1)′. This model has two
Higgs doublets and a singlet, in addition to a singlet exotic quark superfield. We
find that at the one-loop level this model may accommodate the electroweak phase
transitions that are strongly first-order in a reasonably large region of the parameter
space. In the parameter region where the phase transitions take place, we observe
that the lightest scalar Higgs boson has a smaller mass when the strength of the
phase transition becomes weaker. Also, the other three heavier neutral Higgs bosons
get more large masses when the strength of the phase transition becomes weaker.
http://arxiv.org/abs/0704.0328v1
I. INTRODUCTION
The baryon asymmetry of the universe can be dynamically generated during the evolution
of the universe, if the mechanism of baryogenesis satisfies the three Sakharov conditions
[1]. The three Sakharov conditions are: the presence of baryon number violation, the
violation of both C and CP, and a deviation from thermal equilibrium. It is known that
the universe can escape out of the thermal equilibrium by means of electroweak phase
transition, which should be strongly first-order in order to ensure sufficient deviation from
thermal equilibrium to generate the baryon asymmetry that is observed today. However,
it has been already recognized that the Standard Model (SM) has some difficulty to realize
the desired electroweak phase transition. The present experimental lower bound on the
mass of the SM Higgs boson does not allow the electroweak phase transition to be strongly
first-order [2, 3]. The electroweak phase transition is weakly first-order or higher order in
the SM. Thus, the SM is inadequate to generate sufficient baryon asymmetry. Moreover,
the amount CP violation in the Cabibbo-Kobayashi-Maskawa (CKM) matrix is too small
to account for the baryon asymmetry of the observed universe [4].
Consequently, new physical models beyond the SM have extensively been studied for
the possibility of reasonable explanation of the baryon asymmetry of the universe. Espe-
cially, the low energy supersymmetric models have been studied widely within the context
of electroweak baryogenesis [5-7]. The simplest supersymmetric model that includes the
SM is the minimal supersymmetric standard model (MSSM), which possesses in its su-
perpotential the µ term that accounts for the mixing between two Higgs doublets. The
µ parameter, which has the mass dimension, causes some problem with respect to its
energy scale [8]. Several possibilities have been investigated in the literature to solve the
so-called µ problem [9-12]. Introducing an additional U(1)′ to the MSSM is one of the
plausible explanations for the µ problem of the MSSM.
The MSSM with an extra U(1)′ can not only solve the µ problem but we will show
that it can also overcome the difficulties that the SM encounters when the SM tries to
satisfy the Sakharov conditions. This model can accommodate sufficient CP violation,
because it possesses other sources of CP violation besides the CKM matrix. It is possible
to realize the explicit CP violation in this model by means of complex CP phases arising
from the soft SUSY breaking terms [12].
Then, it is the purpose of this paper to show that this model indeed allows the strongly
first-order electroweak phase transitions such that it can successfully explain the baryo-
genesis. The characteristics of the electroweak phase transitions are determined essen-
tially by the temperature-dependent part of the Higgs potential. We construct the full
temperature-dependent Higgs potential at the one-loop level, and examine if the elec-
troweak phase transition may be strongly first-order. Two methods are employed for the
construction of the temperature-dependent Higgs potential. One method assumes that
the critical temperature at which the electroweak phase transition occurs is relatively
high, thus the temperature-dependent effective potential is approximated by retaining
only terms proportional to T 2, whereas the other method carries out numerically exact
integrations of the temperature-dependent effective potential. The thermal effects of par-
ticles whose masses are comparatively smaller than the critical temperature are included
at the one-loop level in the former method, whereas the particle content is different in the
latter method.
Either way, we obtain almost the same physical results. Unlike the MSSM, this model
allows a strongly first-order electroweak phase transition in a wide region of the parame-
ter space, and the first-order electroweak phase transition can be strong enough without
requiring a light stop quark. An interesting behavior of this model with respect to the
strongly first-order electroweak phase transition is that the mass of the lightest neutral
Higgs boson becomes larger when the phase transition gets stronger. On the other hand,
the masses of the other three neutral Higgs bosons become smaller when the phase tran-
sition gets stronger.
II. ZERO TEMPERATURE
The MSSM with an extra U(1)′ accommodates in its Higgs sector two Higgs doublets
H1 = (H
1 , H
1 ), H2 = (H
2 , H
2 ), and one Higgs singlet, S. In terms of these Higgs fields,
the relevant part of the superpotential of this model may be written as
W ≈ htQH2t
R + hbQH1b
R + hkSDLD̄R − λSH
ǫH2 , (1)
where we take into account only the third generation: tcR and b
R are, respectively, the
right-handed singlet top and bottom quark superfields, DR is the right-handed singlet
exotic quark (a vector-like down quark) superfield, Q is the left-handed SU(2) doublet
quark superfield of the third generation, and DL is the left-handed singlet exotic quark
superfield. Further, ht, hb and hk are, respectively, the dimensionless Yukawa coupling
coefficients of top, bottom, and exotic quark superfields, and ǫ is an antisymmetric 2× 2
matrix with ǫ12 = 1.
From the superpotential, at zero temperature, we can construct the Higgs potential
at the tree level, which may be read as
V0 = VF + VD + VS , (2)
where
VF = |λ|
2[(|H1|
2 + |H2|
2)|S|2 + |HT
1~σH1 +H
2~σH2)
(|H1|
2 − |H2|
(Q̃1|H1|
2 + Q̃2|H2|
2 + Q̃3|S|
2)2 ,
VS = m
2 +m2
2 +m2
|S|2 − [λAλ(H
ǫH2)S +H.c.] , (3)
where ~σ denotes the three Pauli matrices, g1, g2, and g
are the U(1), SU(2), and U(1)′
gauge coupling constants, respectively, Q̃1, Q̃2, and Q̃3 are the U(1)
′ hypercharges of H1,
H2, and S, respectively, and m
i (i = 1, 2, 3) are the soft SUSY breaking masses. In the
Higgs potential, λ and Aλ may in general be complex numbers. However, they will be
assumed to be real in the subsequent discussions, as we do not consider CP violation in
the Higgs sector. The soft masses are also assumed to be real, without loss of generality,
and they are eventually eliminated by imposing minimum conditions with respect to the
neutral Higgs fields, The gauge invariance of the superpotential under of U(1)′ requires
that the three U(1)′ hypercharges should satisfy Q̃1 + Q̃2 + Q̃3 = 0.
The above Higgs potential at the tree level would allow the three neutral Higgs fields
, and S to develop the vacuum expectation values (VEVs) v1(0), v2(0), and s(0),
respectively. Remark that these VEVs are obtained at zero temperature. However, for
simplicity, we omit the temperature dependence of these VEVs until next section where
we take into account the finite temperature effect.
The tree-level Higgs potential should now be corrected by the radiative one-loop effects.
In SUSY models, the radiative corrections due to the top and stop quarks contribute
most dominantly to the tree-level Higgs sector. Besides, if tanβ = v2/v1 is very large,
the radiative corrections due to the bottom and sbottom quarks should also be included
since they become no longer negligible. Furthermore, the radiative corrections due to the
exotic quark and squark may be important since the Yukawa coupling of the exotic quark
to the singlet field S can be large at the electroweak scale [11]. Therefore, we take into
account all the contributions from the top, bottom, exotic quark sector to the tree-level
Higgs potential.
The one-loop radiative corrections are evaluated by the effective potential method [13].
We assume that the squark masses are degenerate. Ignoring the mixings in the masses of
the squarks [14], the one-loop effective potential is given by
l=t,b,k
+ log
m̃2 +M2l
, (4)
where t, b, and k, respectively are top, bottom, and exotic quark fields including the
corresponding squark fields, Mt = ht|H2|, Mb = hb|H1|, Mk = hk|S| are the field-
dependent quark masses, and m̃ is the soft SUSY breaking mass, which is assumed that
m̃ = 1000 GeV ≫ mq (q= t, b, or k).
The Higgs sector of the present model consists of six physical Higgs bosons: a pair
of charged Higgs boson, one neutral pseudoscalar Higgs boson, and three neutral scalar
Higgs bosons. The tree-level mass of the charged Higgs boson is given by
m2C± = m
W − λ
2v2 +
2λAλs
sin 2β
, (5)
where v =
v21 + v
2 = 175 GeV and m
W = g
2/2 is the squared mass of the W boson.
At the tree level, the mass of the charged Higgs boson might be either smaller or larger
than the W boson mass.
The tree-level mass of the neutral pseudoscalar Higgs boson is given by
m2A =
2λAλv
sin 2α
, (6)
where tanα = (v/2s) sin 2β implies the splitting between the electroweak symmetry break-
ing scale and the extra U(1)′ symmetry breaking scale. Note that these tree-level masses
of both the neutral pseudoscalar and the charged Higgs bosons do not receive any radiative
corrections, because the squark masses are degenerate.
The tree-level squared masses of the three neutral scalar Higgs bosons are considerably
affected by the radiative corrections. Their squared masses at the one-loop level are given
as the eigenvalues of the 3×3 one-loop level mass matrix, whose elements may be written
M11 = m
Z cos
2 β + 2g
v2 cos2 β +m2A sin
2 β cos2 α + fa(m
M22 = m
Z sin
2 β + 2g
v2 sin2 β +m2A cos
2 β cos2 α + fa(m
t ) ,
M33 = 2g
2 +m2A sin
2 α + fa(m
M12 = g
Q̃1Q̃2v
2 sin 2β + (λ2v2 −m2Z/2) sin 2β −m
A cos β sin β cos
2 α ,
M13 = 2g
1 Q̃1Q̃3vs cos β + 2λ
2vs cosβ −m2A sin β cosα sinα ,
M23 = 2g
Q̃2Q̃3vs sin β + 2λ
2vs sinβ −m2A cos β cosα sinα , (7)
where m2Z = (g
)v2/2 is the squared mass of the Z boson, and the function fa(m
is defined as
3h2qm
m̃2 +m2q
4h2qm
m̃2 +m2q
(m̃2 +m2q)
. (8)
We assume that the masses of three scalar Higgs bosons Si are sorted such that mS1 ≤
mS2 ≤ mS3 .
III. FINITE TEMPERATURE
Now, let us study the temperature dependence of the Higgs potential in order to inves-
tigate the nature of the electroweak phase transition in the MSSM with an extra U(1)′.
We evaluate VT , the temperature-dependent part of the Higgs potential at the one-loop
level, using the effective potential method. It is given as [15]
l=B,F
dx x2 log
1± exp
x2 +m2l (φi)/T
, (9)
where B and F stand for bosons (t̃, b̃, and k̃) and fermions (t, b, and k), and nt = nb =
nk = −12 and nt̃ = nb̃ = nk̃ = 12. The negative sign is for bosons and the positive sign
is for fermions. Thus, the full Higgs potential at finite temperature at the one-loop level
is given by
V (T ) = V0 + V1 + VT (10)
For numerical analysis, we need to set the values of the relevant parameters of the
model. As in the previous section, the soft SUSY breaking mass is set as m̃ = 1000 GeV.
The quark masses are set as mt = 175 GeV, mb = 4 GeV, and mk = 400 GeV. From
these values, mq̃ =
m̃2 +m2q (q = t, b, k) yield the squark masses as mt̃ = 1015 GeV,
= 1000 GeV, and m
= 1077 GeV.
Some caution should be taken for setting the values of Q̃i (i=1, 2, 3), the U(1)
hypercharges of the Higgs doublets and the Higgs singlet. In the MSSM with an extra
U(1)′, the extra neutral gauge boson mass (mZ′) and the mixing angle (αZZ′) between the
two neutral gauge bosons (Z,Z ′) may impose strong constraints on the parameter values.
For our numerical analysis, mZ′ is estimated to be larger than 600 GeV, and αZZ′ smaller
than 2 × 10−3, for tan β = 3 and s(T = 0) = 500 GeV. Besides, as recent research has
suggested [10], we impose the constraint of Q̃1Q̃2 > 0. Further, the U(1)
′ gauge invariance
condition requires that Q̃3 = −(Q̃1 + Q̃2).
In this paper, we define new charges Qi = g
1Q̃i since Q̃i appear always together with
. Then, one may establish the allowed area in the (Q1, Q2)-plane by imposing the above
constraints. For tanβ = 3 and s(T = 0) = 500 GeV, the result is shown in Fig. 1, where
the small area near the point (Q1, Q2) = (-1, 0) and the upper right corner of Fig. 1
are the allowed areas. The hatched region is the excluded area. There are two specific
points in Fig. 1, marked by a star (∗) and a cross (+). The values of Q1 and Q2 at
the star-marked point correspond to the ν-model of E6 gauge group realizations [11]. We
would take the values of Q1 and Q2 at the cross-marked point, namely, (Q1, Q2) = (-1,
-0.1), and hence Q3 =1.1.
With these parameter values at hand, we would investigate the possibility of the
strongly first-order electroweak phase transition by using two different ways. The first
method is to retain only the dominant T 2-proportional part from the high-temperature
approximation of VT , and to take account only those particles whose masses are relatively
small [6]. The second method is to perform the integration in VT in numerically exact way,
and to consider only the contributions of top, bottom, and exotic quarks and squarks.
1. Method A
Let us start with the high temperature approximation of VT , which is expressed as [3]
VT ≈ −
i=t,b,k
T 2m2i (φi)
m4i (φi)
m2i (φi)
cFT 2
i=t̃,b̃,k̃
T 2m2i (φi)
Tm3i (φi)
m4i (φi)
m2i (φi)
cBT 2
, (11)
where log cF = 2.64 and log cB = 5.41. It is known that in the SM the high temperature
approximation is consistent with the exact integration of VT within 5 % at temperature
T for mF/T < 1.6 and mB/T < 2.2, where mF and mB are respectively the fermion mass
and the boson mass that participate in the potential.
We select those terms that are proportional to T 2 in the above expression, which
become most dominant at high temperature. Thus, we assume that the temperature at
which the electroweak phase transition takes place is sufficiently high. We also assume
that the U(1) and SU(2) gaugino masses M1 and M2 in the chargino and neutralino
sectors are very much larger than the other mass parameters. We take into account the
thermal effects due to the Higgs bosons, W , Z, and the extra U(1) gauge boson in the
boson sector, and t, b, k quarks, the lighter chargino, and the three light neutralinos in the
fermion sector, because their masses are relatively small as compared with temperature,
similarly to the analyses of previous articles [6]. Explicitly, the T 2 terms in the high
temperature approximation of VT can be expressed as
+ 4m2
+ 2m2
+ (2g2
+ 6g2
+ 6λ2)(|H1|
2 + |H2|
2) + 12λ2|S|2
+ 12g
2 + Q̃2
2 + Q̃2
|S|2) + 2g
Q̃1Q̃2(|H1|
2 + |H2|
Q̃2Q̃3(|H2|
2 + |S|2) + 2g
Q̃1Q̃3(|H1|
2 + |S|2)
1 (Q̃1 + Q̃2)(Q̃1|H1|
2 + Q̃2|H2|
2 + Q̃3|S|
+6(h2t |H2|
2 + h2b |H1|
2 + h2k|S|
. (12)
Now, the neutral scalar Higgs fields develop the temperature-dependent VEVs, v1(T ),
v2(T ), and s(T ), which we will simply denote v1, v2, and s, respectively. In terms of
these temperature-dependent VEVs, the vacuum at finite temperature is defined as the
minimum of V (T ) as
〈V (v1, v2, s, T )〉 = 〈V0〉+ 〈V1〉+ 〈VT 〉 , (13)
where
〈V0〉 = m
g21 + g
)2 + λ2(v2
s2 + v2
− 2λAλv1v2s+
(Q̃1v
+ Q̃2v
+ Q̃3s
2)2 ,
〈V1〉 = fb(m
t ) + fb(m
b) + fb(m
〈VT 〉 =
+ 4m2
+ 2m2
+ (2g2
+ 6g2
+ 6λ2)(v2
) + 12λ2s2
+ 12g
+ Q̃2
+ Q̃2
s2) + 2g
Q̃1Q̃2(v
1 Q̃2Q̃3(v
2 + s
2) + 2g
1 Q̃1Q̃3(v
1 + s
1 (Q̃1 + Q̃2)(Q̃1v
1 + Q̃2v
2 + Q̃3s
2) + 6(h2tv
2 + h
1 + k
. (14)
In the above expressions, the function fb is defined as
+ log
m̃2 +m2q
, (15)
and the soft SUSY breaking masses at the one-loop level are given as
cos 2β − λ2(s(0)2 + v(0)2 sin2 β) + λAλs(0) tanβ
Q̃1(Q̃1v(0)
2 cos2 β + Q̃2v(0)
2 sin2 β + Q̃3s(0)
2)− fc(m
b(0))
cos 2β − λ2(s(0)2 + v(0)2 cos2 β) + λAλs(0) cotβ
Q̃2(Q̃1v(0)
2 cos2 β + Q̃2v(0)
2 sin2 β + Q̃3s(0)
2)− fc(m
t (0))
= − λ2v(0)2 +
2s(0)
v(0)2Aλ sin 2β
Q̃3(Q̃1v(0)
2 cos2 β + Q̃2v(0)
2 sin2 β + Q̃3s(0)
2)− fc(m
k(0)) , (16)
where v1(0), v2(0), and s(0) are the VEVs evaluated at zero temperature in the preceding
section, tan β = v2(0)/v1(0), v(0) =
v1(0)2 + v2(0)2 = 175 GeV, and the function fc is
defined as
3h2qm
2 + 2 log
m̃2 +m2q
m̃2 +m2q
. (17)
Now, let us determine the critical temperature at which the electroweak phase tran-
sition takes place. In our analysis, the critical temperature is defined by a temperature
at which 〈V (T )〉 has two distinct minima with equal value, that is, a pair of degenerate
vacua. In order to have a pair of degenerate vacua, the potential 〈V (T )〉 should satisfy
the minimum condition of
0 = 2m2
s− 2λAλv1v2 + 2λ
1 Q̃3s(Q̃1v
1 + Q̃2v
2 + Q̃3s
2) + 2h2kmkfc(m
s[24λ2 + 24g
+ 20g
Q̃3(Q̃1 + Q̃2) + 12k
2] , (18)
which is obtained by calculating the first derivative of the full effective potential at the
finite temperature with respect to s.
For given parameter values at given temperature, one may solve the above minimum
condition to express s in terms of the other two VEVs, v1 and v2. Then, by substituting
s into 〈V (v1, v2, s, T )〉, one may obtain 〈V (v1, v2, T )〉 which depends only on v1 and v2.
By inspecting the shape of 〈V (v1, v2, T )〉 on the (v1, v2)-plane for given parameter values
at given temperature, we may determine whether it possess a pair of degenerate vacua or
In Fig. 2, the equipotential contours of 〈V (v1, v2, T )〉 are plotted on the (v1, v2)-
plane, where the parameter values are set as tanβ = 3, λ = 0.8, s(0) = 500 GeV,
mA = 1830 GeV, and the temperature is set as T = 100 GeV, which is actually the
critical temperature Tc. One can easily spot two distinct minima of 〈V (v1, v2, T )〉 on
the (v1, v2)-plane, namely, one at (0, 0) and the other at (275, 640) GeV. The phase of
the state is symmetric at the minimum point (0, 0) on the (v1, v2)-plane, whereas it is
broken at (275, 640) GeV. The electroweak phase transition may take place from (0, 0) to
(275, 640) GeV on the (v1, v2)-plane, which is evidently discontinuous and therefore it is
first-order.
The distance on the (v1, v2)-plane between the two minima of 〈V (v1, v2, T )〉, defined
as vc, determines the strength of the electroweak phase transition. The electroweak phase
transition is said to be strong if vc/Tc > 1, and weak otherwise. In Fig. 2, the distance is
calculated to be
(275− 0)2 + (640− 0)2 = 696 (GeV) . (19)
In Fig. 2, the strength of the electroweak phase transition is about vc/Tc = 6.9, which
definitely tells that the electroweak phase transition is a strong one. Therefore, the
particular parameter values set for Fig. 2 yields an electroweak phase transition which
is first-order as well as strong. Note that vc does not depend on s, that is, we need not
to know the values of s at the two minima to calculate vc. Actually, vc is the VEV at
the broken phase. The masses of the neutral scalar Higgs bosons at zero temperature
for the parameter values of Fig. 2 are obtained as mS1 = 56 GeV, mS2 = 807 GeV, and
mS3 = 1827 GeV.
We repeat the above job of analysis, varying the values of the relevant parameters. We
find that there are a large number of sets of parameter values that allow strongly first-order
electroweak phase transitions. Thus, the MSSM with an extra U(1)′ may accommodate
TABLE 1: Some sets of λ and mA that allow strongly first-order electroweak phase
transitions in the MSSM with an extra U(1)′, obtained by Method A. The values of other
parameters are fixed as tanβ = 3, s(0) = 500 GeV, m̃ = 1000 GeV, and Tc = 100
GeV. The pair of numbers in the third column are the coordinates of the broken-phase
minimum of 〈V (v1, v2, T )〉. The coordinates of its symmetric-phase minimum is (0, 0)
for all sets. The three numbers in the fourth column are the masses of S1, S2, and S3,
respectively. The number in the last column is the strength of the first-order electroweak
phase transition.
λ mA (GeV) (v1, v2) (GeV) mS1 , mS2 , mS3 (GeV) vc/Tc
0.1 478 (1750, 1650) 120, 524, 792 26
0.2 675 (1400, 1500) 118, 674, 796 23
0.3 900 (1200, 1400) 112, 786, 908 18
0.4 1109 (870, 1200) 104, 792, 1112 15
0.5 1306 (600, 1000) 93, 796, 1307 12
0.6 1486 (430, 850) 82, 800, 1485 8
0.7 1660 (340, 700) 70, 803, 1658 7
0.8 1830 (275, 640) 56, 807, 1827 6.9
the desired phase transitions for a wide region in its parameter space. Some of the results
are listed in Table 1, where tanβ = 3, s(0) = 500 GeV, and T = 100 GeV are fixed as
the values set in Fig. 2, whereas λ and mA have different values. The set of numbers in
the last row of Table 1 is the numerical result of Fig. 2.
Every set of numbers in each row of Table 1 gives 〈V (v1, v2, T )〉 a pair of degenerate
minima, the minimum of symmetric phase at (0, 0) on the (v1, v2)-plane, and the one of
broken phase at a different point on the (v1, v2)-plane as given in Table 1. The electroweak
phase transition is strongly first-order. One may easily observe in Table 1 that, as the
value of λ increases, a larger value of mA allow desired phase transitions. On the other
hand, the strength of the phase transition is reinforced if the value of λ decreases.
The masses of the neutral scalar Higgs bosons exhibit some interesting behavior. For a
larger value of mA, both S2 and S3 have also larger masses whereas S1 has a smaller mass.
The tendency is that the strength of the phase transition is reinforced if mS1 increases
and if mA, mS2 , and mS3 decrease. In the SM, the strength of the first order electroweak
phase transition decreases if its single Higgs boson mass is increased. Also, in the MSSM,
we have a weaker phase transition if the lighter one of its two scalar Higgs bosons has a
larger mass. In this regard, the tendency of our model is opposite to those of the SM or
the MSSM. One can see that this strange behavior also occurs in some parameter region
of a non-minimal SUSY model, as shown in Fig. 3 of Ref. [7].
2. Method B
The second method evaluates VT by exact integration to obtain the temperature-dependent
full potential V (T ) at one-loop level, where the thermal effects of top, bottom, and exotic
quarks and squarks are taken into account. The thermal effects of the gauge bosons can
be a help for strengthening the first-order electroweak phase transition, but we would
omit them, since the strength of the phase transition is already strong enough.
This method starts with the exact integral expression for 〈VT 〉 after replacing the
neutral Higgs fields by their VEVs as
〈VT 〉 = −
l=t,b,k
dx x2 log
1− exp
m2l (v1, v2, s)
l=t̃,b̃,k̃
dx x2 log
1 + exp
m̃2 +m2l (v1, v2, s)
,(20)
which is different from 〈VT 〉 of Method A, while 〈V0〉 and 〈V1〉 are the same as those of
Method A. From the full 〈V (T )〉 = 〈V0〉 + 〈V1〉 + 〈VT 〉, we obtain a minimum condition
for degenerate vacua as
0 = 2m2
s− 2λAλv1v2 + 2λ
)s+ 2g
Q3s(Q̃1v
+ Q̃2v
+ Q̃3s
+ 2h2kmkfc(m
dx x2
2h2ks exp(−
x2 +m2k/T
x2 +m2k/T
1 + exp(−
x2 +m2k/T
dx x2
2h2ks exp(−
x2 + (m̃2 +m2k)/T
x2 + (m̃2 +m2k)/T
1 + exp(−
x2 + (m̃2 +m2k)/T
] , (21)
where mk depends only on s and is independent from v1 and v2.
Solving the above minimum condition is harder than solving the corresponding mini-
mum condition of Method A. Nevertheless, we can solve it by using the bisection method
to express s in terms of the other parameters. Then, eliminating s from 〈V (T )〉, we can
obtain the expression for 〈V (v1, v2, T )〉 which depends only on v1 and v2. Subsequent
steps of numerical analysis are the same as the previous method.
In Fig. 3, equipotential contours of 〈V (v1, v2, T )〉 obtained by the present method is
plotted on the (v1, v2)-plane, where the parameter values are set slightly different from
the previous method: tan β = 3, λ = 0.8, s(0) = 500 GeV, mA = 1780 GeV, and T = 100
GeV. The shape of the equipotential contours of Fig. 3 is almost the same as that of
Fig. 2. One can see that there are two distinct minima in Fig. 3, just like Fig. 2:
one at (0, 0), and the other at (165, 440) GeV on the (v1, v2)-plane, indicating that the
phase transition is first order. The strength of the first-order phase transition is strong,
since vc/Tc = 4.7. The masses of the three scalar Higgs bosons are evaluated at zero
temperature as mS1 = 82 GeV, mS2 = 804 GeV, and mS3 = 1777 GeV.
TABLE 2: Some sets of λ and mA that allow strongly first-order electroweak phase
transitions in the MSSM with an extra U(1)′, obtained by Method B. Other descriptions
are the same as Table 1.
λ mA GeV (v1B, v2B) GeV mSi GeV vc/Tc
0.1 462 (1600, 1600) 121, 468, 791 22
0.2 663 (1400, 1400) 118, 662, 795 19
0.3 885 (1100, 1100) 113, 785, 894 15
0.4 1095 (800, 1200) 106, 792, 1098 14
0.5 1287 (680, 990) 97, 796, 1288 12
0.6 1457 (400, 750) 91, 799, 1456 8
0.7 1620 (300, 600) 86, 801, 1618 6
0.8 1780 (165, 440) 82, 804, 1777 4.7
Comparing Fig. 3 with Fig. 2, one may safely remark that Method A and Method
B lead qualitatively the same results. Either method, whether 〈VT 〉 is calculated by
direct integration or is simplified by high-temperature approximation, and whether the
participating particles at the one-loop level are somewhat exhaustive or selective, we
find that the MSSM with and extra U(1)′ allows strongly first-order electroweak phase
transitions for certain region in its parameter space.
We repeat the numerical analysis by varying the parameter values. and some of the
results are listed in Table 2. Like in Table 1, tan β = 3, s(0) = 500 GeV, and T = 100
GeV are fixed, whereas λ and mA are varied. The set of numbers in the last row of Table
2 is the numerical result of Fig. 3. Comparing Table 2 with Table 1, one may notice that
the numbers are slightly different from each other but the general behavior of the two
tables is exactly the same.
IV. DISCUSSIONS AND CONCLUSIONS
We investigate the MSSM with an extra U(1)′ if it could accommodate strongly first-
order electroweak phase transitions to provide sufficient baryon asymmetry, for reasonable
masses of scalar Higgs bosons. To do so, we need the temperature-dependent part of
the Higgs potential at the one-loop level. Explicitly, its expression is obtained by two
complementary methods: Method A employs high-temperature approximation and retains
only the most dominant T 2 terms, and takes into account the thermal effects at the one-
loop level of various participating particles. On the other hand, method B performs
numerical integrations, and the thermal effects of top, bottom, and exotic quarks and
squarks are accounted for.
Both methods lead us to essentially the same conclusion: the strongly first-order
electroweak phase transition is possible in the MSSM with an extra U(1)′, for a wide
region in its parameter space. The masses of the scalar Higgs bosons are obtained within
reasonably acceptable ranges. Accordingly, we may expect that the MSSM with an extra
U(1)′ can explain the baryon asymmetry of the universe.
We remark that the MSSM with an extra U(1)′ exhibits an interesting behavior with
respect to the correlation between the strength of the phase transition and the Higgs
boson masses. The MSSM with an extra U(1)′ is opposite to the SM or to the MSSM
in the sense that the mass of the lightest scalar Higgs boson increases when the strength
of the strongly first-order electroweak phase transition becomes stronger. In the SM, its
single Higgs boson has a larger mass when the strength of the first order electroweak
phase transition decreases. In the MSSM, we also have a larger mass for the lighter one
of its two scalar Higgs bosons when the phase transition becomes weaker.
ACKNOWLEDGMENTS
This research is supported by KOSEF through CHEP. The authors would like to
acknowledge the support from KISTI (Korea Institute of Science and Technology Infor-
mation) under ”The Strategic Supercomputing Support Program” with Dr. Kihyeon Cho
as the technical supporter. The use of the computing system of the Supercomputing
Center is also greatly appreciated.
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http://arxiv.org/abs/hep-ph/0703041
FIGURE CAPTION
FIG. 1. : The allowed area in the (Q1, Q2)-plane. For tanβ = 3 and s(T = 0) = 500
GeV, the small area near the point (Q1, Q2) = (-1, 0) and the upper right corner are the
allowed areas, whereas the hatched region is the excluded area. There are two specific
points, marked by a star (∗) and a cross (+). The values of Q1 and Q2 at the star-marked
point correspond to the ν-model of E6 gauge group realizations. The values of Q1 and Q2
at the cross-marked point are (Q1, Q2) = (-1, -0.1), and hence Q3 =1.1. In our discussions,
we choose this point.
FIG. 2. : The plot of the equipotential contours of 〈V (v1, v2, T )〉 on the (v1, v2)-plane,
obtained by Method A. The parameter values are set as tan β = 3, λ = 0.8, s(0) = 500
GeV, mA = 1830 GeV, and the temperature is set as T = 100 GeV, which is actually the
critical temperature Tc. Notice two distinct minima of 〈V (v1, v2, T )〉 on the (v1, v2)-plane:
(0, 0) where the phase of the state is symmetric, and (275, 640) GeV, where the phase
of the state is broken. The electroweak phase transition may take place from (0, 0) to
(275, 640) GeV on the (v1, v2)-plane, which is evidently discontinuous and therefore it is
first order. The distance between the two minima is vc = 696 GeV, indicating that the
strength of the first-order phase transition is strong (vc/Tc > 1). The masses of the three
scalar Higgs bosons are obtained as mS1 = 56 GeV, mS2 = 807 GeV, and mS3 = 1827
FIG. 3. : The plot of the equipotential contours of 〈V (v1, v2, T )〉 on (v1, v2)-plane, ob-
tained by Method B. The parameter values are set as tan β = 3, λ = 0.8, s(0) = 500
GeV, mA = 1780 GeV, and Tc = 100 GeV. The coordinates of two minima are: (0, 0)
and (165, 440) GeV. The distance between the two minima is vc = 470 GeV, thus the
electroweak phase transition between the two minima is strongly first-order. The masses
of the three scalar Higgs bosons are obtained as mS1 = 82 GeV, mS2 = 804 GeV, and
mS3 = 1777 GeV.
-1 -0.8 -0.6 -0.4 -0.2 0 0.2 0.4 0.6 0.8 1
FIG. 1: The allowed area in the (Q1, Q2)-plane. For tanβ = 3 and s(T = 0) = 500
GeV, the small area near the point (Q1, Q2) = (-1, 0) and the upper right corner are the
allowed areas, whereas the hatched region is the excluded area. There are two specific
points, marked by a star (∗) and a cross (+). The values of Q1 and Q2 at the star-marked
point correspond to the ν-model of E6 gauge group realizations. The values of Q1 and Q2
at the cross-marked point are (Q1, Q2) = (-1, -0.1), and hence Q3 =1.1. In our discussions,
we choose this point.
0 50 100 150 200 250 300 350 400
V1 (GeV)
V2 (GeV)
FIG. 2: The plot of the equipotential contours of 〈V (v1, v2, T )〉 on the (v1, v2)-plane,
obtained by Method A. The parameter values are set as tan β = 3, λ = 0.8, s(0) = 500
GeV, mA = 1830 GeV, and the temperature is set as T = 100 GeV, which is actually the
critical temperature Tc. Notice two distinct minima of 〈V (v1, v2, T )〉 on the (v1, v2)-plane:
(0, 0) where the phase of the state is symmetric, and (275, 640) GeV, where the phase
of the state is broken. The electroweak phase transition may take place from (0, 0) to
(275, 640) GeV on the (v1, v2)-plane, which is evidently discontinuous and therefore it is
first order. The distance between the two minima is vc = 696 GeV, indicating that the
strength of the first-order phase transition is strong (vc/Tc > 1). The masses of the three
scalar Higgs bosons are obtained as mS1 = 56 GeV, mS2 = 807 GeV, and mS3 = 1827
0 25 50 75 100 125 150 175 200 225 250
V1 (GeV)
V2 (GeV)
FIG. 3: The plot of the equipotential contours of 〈V (v1, v2, T )〉 on (v1, v2)-plane, obtained
by Method B. The parameter values are set as tan β = 3, λ = 0.8, s(0) = 500 GeV, mA =
1780 GeV, and Tc = 100 GeV. The coordinates of two minima are: (0, 0) and (165, 440)
GeV. The distance between the two minima is vc = 470 GeV, thus the electroweak phase
transition between the two minima is strongly first-order. The masses of the three scalar
Higgs bosons are obtained as mS1 = 82 GeV, mS2 = 804 GeV, and mS3 = 1777 GeV.
INTRODUCTION
ZERO TEMPERATURE
FINITE TEMPERATURE
Method A
Method B
DISCUSSIONS AND CONCLUSIONS
|
0704.0329 | Solutions of fractional reaction-diffusion equations in terms of the
H-function | arXiv:0704.0329v2 [math.PR] 7 Aug 2007
SOLUTIONS OF FRACTIONAL REACTION-DIFFUSION
EQUATIONS IN TERMS OF THE H-FUNCTION
H.J. HAUBOLD
Office for Outer Space Affairs, United Nations, Vienna International Centre
P.O. Box 500, A-1400, Vienna, Austria and
Centre for Mathematical Sciences, Pala Campus
Arunapuram P.O., Pala-686 574, Kerala, India
A.M .MATHAI
Department of Mathematics and Statistics, McGill University
Montreal, Canada H3A 2K6 and
Centre for Mathematical Sciences, Pala Campus
Arunapuram P.O., Pala-686 574, Kerala, India
R.K. SAXENA
Department of Mathematics and Statistics, Jai Narain Vyas University
Jodhpur-342004, India
Abstract. This paper deals with the investigation of the solution of an unified
fractional reaction-diffusion equation associated with the Caputo derivative as
the time-derivative and Riesz-Feller fractional derivative as the space-derivative.
The solution is derived by the application of the Laplace and Fourier transforms
in closed form in terms of the H-function. The results derived are of general
nature and include the results investigated earlier by many authors, notably
by Mainardi et al. (2001, 2005) for the fundamental solution of the space-time
fractional diffusion equation, and Saxena et al. (2006a, b) for fractional reaction-
diffusion equations. The advantage of using Riesz-Feller derivative lies in the
fact that the solution of the fractional reaction-diffusion equation containing this
derivative includes the fundamental solution for space-time fractional diffusion,
which itself is a generalization of neutral fractional diffusion, space-fractional
diffusion, and time-fractional diffusion. These specialized types of diffusion can
be interpreted as spatial probability density functions evolving in time and are
expressible in terms of the H-functions in compact form.
1 Introduction
The review of the theory and applications of reaction-diffusion systems is con-
tained in many books and articles. In recent work authors have demonstrated
the depth of mathematics and related physical issues of reaction-diffusion equa-
tions such as nonlinear phenomena, stationary and spatio-temporal dissipative
pattern formation, oscillations, waves etc. (Frank, 2005; Grafiychuk, Datsko,
http://arxiv.org/abs/0704.0329v2
and Meleshko, 2006, 20076). In recent time, interest in fractional reaction-
diffusion equations has increased because the equation exhibits self-organization
phenomena and introduces a new parameter, the fractional index, into the equa-
tion. Additionally, the analysis of fractional reaction-diffusion equations is of
great interest from the analytical and numerical point of view.
The objective of this paper is to derive the solution of an unified model of
reaction-diffusion system (14), associated with the Caputo derivative and the
Riesz-Feller derivative. This new model provides the extension of the models
discussed earlier by Mainardi, Luchko, and Pagnini (2001), Mainardi, Pagnini,
and Saxena (2005), and Saxena, Mathai, and Haubold (2006a). The present
study is in continuation of our earlier work, Haubold and Mathai (1995, 2000)
and Saxena, Mathai, and Haubold (2006a, 2006b).
2 Results Required in the Sequel
In view of the results
J−1/2(x) =
cosx. (1)
and (Mathai and Saxena, 1978, p. 49), the cosine transform of the H-function
is given by
tρ−1cos(kt)Hm,np,q
(ap,Ap)
(bq,Bq)
dt (2)
n+1,m
q+1,p+2
(1−bq,Bq),(
(ρ,µ),(1−ap,ap),(
, (3)
where Re[ρ + µmin1≤j≤m(
)] > 0, Re[ρ+ µmax1≤j≤n
] < 0, |argα| < 1
πΩ, Ω >
k > 0 and Ω =
j=1 Bj −
j=m+1 Bj +
j=1 Aj −
j=n+1 Aj .
The Riemann-Liouville fractional integral of order ν is defined by (Miller and
Ross, 1993, p. 45; Kilbas et al., 2006)
t N(x, t) =
(t − u)ν−1N(x, u)du, (4)
where Re(ν) > 0.
The following fractional derivative of order α > 0 is introduced by Caputo
(1969; see also Kilbas et al., 2006) in the form
t f(x, t) =
Γ(m − α)
f (m)(x, τ)dτ
(t − τ)α+1−m
, m − 1 < α ≤ m, Re(α) > 0, m ∈ N.
∂mf(x, t)
, if α = m. (5)
where ∂
f(x, t) is the mth partial derivative of f(x,t) with respect to t.
The Laplace transform of the Caputo derivative is given by Caputo (1969;
see also Kilbas et al., 2006) in the form
L {0D
t f(x, t); s} = s
αF (x, s)−
sα−r−1f (r)(x, 0+), (m− 1 < α ≤ m). (6)
Following Feller (1952, 1971), it is conventional to define the Riesz-Feller
space-fractional derivative of order α and skewness θ in terms of its Fourier
transform as
F {xD
θ f(x); k} = −Ψ
α(k)f
∗(k), (7)
where
Ψθα(k) = |k|
αexp[i(signk)
], 0 < α ≤ 2, |θ| ≤ min {α, 2 − α} . (8)
When θ = 0, then (8) reduces to
F {xD
0 f(x); k} = −|k|
α, (9)
which is the Fourier transform of the Weyl fractional operator, defined by
xf(t) =
Γ(n − µ)
f(u)du
(t − u)µ−n+1
. (10)
This shows that the Riesz-Feller operator may be regarded as a generalization
of the Weyl operator.
Further, when θ = 0, we have a symmetric operator with respect to x that
can be interpreted as
0 = −
This can be formally deduced by writing −(k)α = −(k2)α/2. For 0 < α < 2 and
|θ| ≤ min {α, 2 − α}, the Riesz-Feller derivative can be shown to possess the
following integral representation in the x domain:
θ f(x) =
Γ(1 + α)
sin[(α + θ)π/2]
f(x + ξ) − f(x)
+ sin[(α − θ)π/2]
f(x − ξ) − f(x)
. (12)
Finally, we need the following property of the H-function (Mathai and Sax-
ena, 1978)
Hm,np,q
(ap,ap)
(bq ,Bq)
Hm,np,q
(ap,Ap/δ)
(bq,Bq/δ
, (δ > 0). (13)
3 Unified Fractional Reaction-Diffusion Equa-
In this section, we will investigate the solution of the reaction-diffusion equation
(14) under the initial conditions (15). The result is given in the form of the
following
Theorem. Consider the unified fractional reaction-diffusion model
t N(x, t) = ηxD
θ N(x, t) + Φ(x, t), (14)
where η, t > 0, x ∈ r; α, θ, β are real parameters with the constraints
0 < α ≤ 2, |θ| ≤ min(α, 2 − α), 0 < β ≤ 2, and the initial conditions
N(x, 0) = f(x), Nt(x, 0) = g(x) ); for x ∈ R,
|x|→±∞ N(x, t) = 0, t > 0. (15)
Here Nt(x, 0) means the first partial derivative of N(x, t) with respect to t
evaluated at t = 0, η is a diffusion constant and Φ(x, t) is a nonlinear function
belonging to the area of reaction-diffusion. Further xD
θ is the Riesz-Feller
space-fractional derivative of order α and asymmetry θ. 0D
t is the Caputo
time-fractional derivative of order β. Then for the solution of (14), subject to
the above constraints, there holds the formula
N(x, t) =
f∗(k)Eβ,1(−ηt
βΨθα(k))exp(−ikx)dk (16)
tg∗(k)Eβ,2(−ηk
αtβΨθα(k))exp(−ikx)dk
ξβ−1dξ
Φ∗(k, t − ξ)Eβ,β(−ηk
αtβΨθα(k))exp(−ikx)dk.
In equation (16) and the following, Eα,β(z) denotes the generalized Mittag-
Leffler function (Saxena, Mathai, and Haubold, 2004; Berberan-Santos, 2005;
Chamati and Tonchev, 2006).
Proof. If we apply the Laplace transform with respect to the time variable t,
Fourier transform with respect to space variable x, and use the initial conditions
(15) and the formula (7), then the given equation transforms into the form
∼(k, s) − sβ−1f∗(k) − sβ−2g∗(k) = −ηΨθα(k)N
∼(k, s) + Φ
∼(k, s),
where according to the conventions followed , the symbol ∼ will stand for the
Laplace transform with respect to time variable t and * represents the Fourier
transform with respect to space variable x.
Solving for N
∼ , it yields
∼(k, s) =
f∗(k)sβ−1
sβ + ηΨθα(k)
g∗(k)sβ−2
sβ + ηΨθα(k)
sβ + ηΨθα(k)
. (17)
On taking the inverse Laplace transform of (17) and applying the formula
a + sα
= tα−βEα,α−β+1(−at
α), (18)
where Re(s) > 0, Re(α) > 0, Re(α − β) > −1; it is seen that
N∗(k, t) = f∗(k)Eβ,1(−ηt
βΨθα(k)) + g
∗(k)tEβ,2(−ηt
βΨθα(k))
Φ∗(k, t − ξ)ξβ−1Eβ,β(−ηΨ
α(k)ξ
β)dξ. (19)
The required solution (16) is now obtained by taking the inverse Fourier trans-
form of (19). This completes the proof of the theorem.
4 Special Cases
When g(x) = 0, then by the application of the convolution theorem of the
Fourier transform to the solution (16) of the theorem, it readily yields
Corollary 1. The solution of the fractional reaction-diffusion equation
N(x, t) − η
N(x, t) = Φ(x, t), x ∈ R, t > 0, η > 0, (20)
with initial conditions
N(x, 0) = f(x), Nt(x, 0) = 0 for x ∈ R, 1 < β ≤ 2,
x→±∞ N(x, t) = 0, (21)
where η is a diffusion constant and Φ(x, t) is a nonlinear function belonging to
the area of reaction-diffusion, is given by
N(x, t) =
G1(x − τ, t)f(τ)dτ
(t − ξ)β−1dξ
G2(x − τ, t − ξ)Φ(τ, ξ)dτ, (22)
where
α − θ
G1(x, t) =
exp(−ikx)Eβ,1(−η|t
β |Ψθα(k))dk (23)
η1/αtβ/α
(1,1/α),(β,β/α),(1,ρ)
(1,1/α),(1,1),(1,ρ)
, (α > 0)
G2(x, t) =
exp(−ikx)Eβ,β(−ηt
βΨθα(k))dk
η1/αtβ/α
(1,1/α),(β,β/α),(1,ρ)
(1,1/α),(1,1),(1,ρ)
, (α > 0). (24)
In deriving the above results, we have used the inverse Fourier transform formula
F−1[Eβ,γ(−ηt
βΨαθ (k)); x] =
3,3 [
η1αtβ/α
(1,1/α),(γ,β/α),(1,ρ)
(1,1/α),(1,1),(1,ρ)
], (25)
where Re(β) > 0, Re(γ) > 0, which can be established by following a procedure
similar to that employed by Mainardi, Luchko, and Pagnini (2001). Next , if
we set f(x) = δ(x), Φ = 0, g(x) = 0, where δ(x) is the Dirac delta-function,
then we arrive at the following interesting result given by Mainardi, Pagnini,
and Saxena (2005).
Corollary 2. Consider the following space-time fractional diffusion model
∂βN(x, t)
= η xD
θ N(x, t), η > 0, x ∈ R, 0 < β ≤ 2, (26)
with the initial conditions N(x, t = 0) = δ(x), Nt(x, 0) = 0,
x→±∞ N(x, t) = 0
where η is a diffusion constant and δ(x) is the Dirac delta-function. Then for
the fundamental solution of (26) with initial conditions, there holds the formula
N(x, t) =
3.3 [
(ηtβ)1/α
(1,1/α),(1,β/α),(1,ρ)
(1,1/α),(1,1),(1,ρ)
], (27)
where ρ = α−θ
Some interesting special cases of (26) are enumerated below.
(i) We note that for α = β, Mainardi, Pagnini, and Saxena (2005) have
shown that the corresponding solution of (26), denoted by Nθα, which we call as
the neutral fractional diffusion, can be expressed in terms of elementary function
and can be defined for x > 0 as
Neutral fractional diffusion: 0 < α = β < 2; θ ≤ min {α, 2 − α} ,
Nθα(x) =
xα−1sin[(π/2)(α − θ)]
1 + 2xαcos[(π/2)(α − θ)] + x2α
. (28)
The neutral fractional diffusion is not studied at length in the literature.
Next we derive some stable densities in terms of the H-functions as special
cases of the solution of the equation (26)
(ii) If we set β = 1, 0 < α < 2; θ ≤ min {α, 2 − α}then (26) reduces to space
fractional diffusion equation, which we denote by Lθα(x) is the fundamental
solution of the following space-time fractional diffusion model:
∂N(x, t)
= η xD
θ N(x, t), η > 0, x ∈ R, (29)
with the initial conditions N(x, t = 0) = δ(x), limx→±∞N(x, t) = 0,, where η is a
diffusion constant and δ(x) is the Dirac-delta function. Hence for the solution
of (29) there holds the formula
Lθα(x) =
α(ηt)1/α
(ηt)1/α
(1,1),(ρ,ρ)
),(ρ,ρ)
, 0 < α < 1, |θ| ≤ α, (30)
where ρ = α−θ
. The density represented by the above expression is known as
α-stable Lévy density. Another form of this density is given by
Lθα(x) =
α(ηt)1/α
(ηt)1/α
(1− 1
),(1−ρ,ρ)
(0,1),(1−ρ,ρ)
, 1 < α < 2, |θ| ≤ 2 − α,
(iii) Next, if we take α = 2, 0 < β < 2, θ = 0, then we obtain the time
fractional diffusion, which is governed by the following time fractional diffusion
model:
∂βN(x, t)
N(x, t), η > 0, x ∈ R, 0 < β ≤ 2, (32)
with the initial conditions N(x, t = 0) = δ(x), Nt(x, 0) = 0,
x→±∞ N(x, t) = 0
where η is a diffusion constant and δ(x) is the Dirac delta-function, whose
fundamental solution is given by the equation
N(x, t) =
(ηtβ)1/2
(1,β/2)
(1,1)
. (33)
(iv) Further, if we set α = 2, β = 1 and θ → 0 then for the fundamental
solution of the standard diffusion equation
N(x, t) = η
N(x, t), (34)
with initial condition
N(x, t = 0) = δ(x), limx→±∞N(x, t) = 0, (35)
there holds the formula
N(x, t) =
η1/2t1/2
(1,1/2)
(1,1)
= (4πηt)−1/2exp[−
], (36)
which is the classical Gaussian density. For further details of these special cases
based on the Green function, one can refer to the paper by Mainardi, Luchko,
and Pagnini (2001) and Mainardi, Pagnini, and Saxena (2005).
Remark. Fractional order moments and the asymptotic expansion of the solu-
tion (27) are discussed by Mainardi, Luchko, and Pagnini (2001).
Finally, for β = 1/2 in (14), we arrive at
Corollary 3. Consider the following fractional reaction-diffusion model
t N(x, t) = ηxD
θ N(x, t) + Φ(x, t), (37)
where η, t > 0, x ∈ R; α, θ are real parameters with the constraints
0 < α ≤ 2, |θ| ≤ min(α, 2 − α), and the initial conditions
N(x, 0) = f(x), for x ∈ R, limx→±∞N(x, t) = 0. (38)
Here η is a diffusion constant and Φ(x, t) is a nonlinear function belonging to
the area of reaction-diffusion. Further xD
θ is the Riesz-Feller space fractional
derivative of order α and asymmetry θ and D
t is the Caputo time-fractional
derivative of order 1/2. Then for the solution of (37), subject to the above
constraints, there holds the formula
N(x, t) =
f∗(k)E1/2,1(−ηt
βΨθα(k))exp(−ikx)dk (39)
ξ−1/2dξ
Φ∗(kct − ξ)E 1
(−ηkαt1.2Ψθα(k))exp(−ikx)dk.
If we set θ = 0 in (39), then it reduces to the result recently obtained by the
authors (2006a) for the fractional reaction-diffusion equation.
5 References
Berberan-Santos, M.N. (2005). Properties of the Mittag-Leffler relaxation func-
tion, Journal of Mathematical Chemistry, 38, 629-635.
Caputo, M. (1969). Elasticita e Dissipazione, Zanichelli, Bologna.
Chamati, H. and Tonchev, N.S. (2006). Generalized Mittag-Leffler functions
in the theory of finite-size scaling for systems with strong anisotropy and/or
long-range interaction, Journal of Physics A: Mathematical and General, 39,
469-478.
Feller, W. (1952). On a generalization of Marcel Riesz’ potentials and the
semi-groups generated by them, Meddeladen Lund Universitets Matematiska
Seminarium (Comm. Sém. Mathém. Université de Lund ), Tome suppl. dédié
a M. Riesz, Lund, 73-81.
Feller, W. (1966). An Introduction to Probability Theory and its Applications,
Vol. II, John Wiley and Sons, New York.
Frank, T.D. (2005). Nonlinear Fokker-Planck Equations: Fundamentals and
Applications, Springer, Berlin Heidelberg New York.
Grafiychuk, V., Datsko, B., and Meleshko, V. (2006). Mathematical model-
ing of pattern formation in sub- and superdiffusive reaction-diffusion systems,
arXiv:nlin.AO/06110005 v3.
Grafiychuk, V., Datsko, B., and Meleshko, V. (2007). Nonlinear oscillations and
stability domains in fractional reaction-diffusion systems, arXiv:nlin.PS/0702013
Haubold, H.J. and Mathai, A.M. (2000). The fractional kinetic equation and
thermonuclear functions, Astrophysics and Space Science, 273, 53-63.
Haubold, H.J. and Mathai, A.M. (1995). A heuristic remark on the periodic
variation in the number of solar neutrinos detected on Earth, Astrophysics and
Space Science, 228, 113-124.
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tions of Fractional Differential Equations, Elsevier, Amsterdam.
Mainardi, F., Luchko, Y., and Pagnini, G. (2001). The fundamental solution
of the space-time fractional diffusion equation, Fractional Calculus and Applied
Analysis. 4, 153-192.
Mainardi, F., Pagnini, G., and Saxena, R.K. (2005). Fox H-functions in frac-
tional diffusion, Journal of Computational and Applied Mathematics 178, 321-
Mathai, A.M. and Saxena, R.K. (1978). The H-function with Applications in
Statistics and Other Disciplines, John Wiley and Sons, New York, London, and
Sydney.
Miller, K.S. and Ross, B. (1993). An Introduction to the Fractional Calculus
and Fractional Differential Equations, John Wiley and Sons, New York.
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equations, Astrophysics and Space Science, 282, 281-287.
Saxena, R.K., Mathai, A.M., and Haubold, H.J. (2006a). Fractional reaction-
diffusion equations, Astrophysics and Space Science, 305, 289-296.
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systems and nonlinear waves, Astrophysics and Space Science, 305, 297-303.
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|
0704.0330 | Random Matrix Theory at Nonzero $\mu$ and $T$ | Random Matrix Theory at Nonzero µ and T
Kim Splittorff1,∗) and Jacobus Johannes Maria Verbaarschot1,2 ,3 ,∗∗)
1 The Niels Bohr Institute, Blegdamsvej 17, DK-2100, Copenhagen Ø, Denmark
2 Niels Bohr International Academy, Blegdamsvej 17, DK-2100, Copenhagen Ø
3 Department of Physics and Astronomy, SUNY, Stony Brook, New York 11794
We review applications of random matrix theory to QCD at nonzero temperature and
chemical potential. The chiral phase transition of QCD and QCD-like theories is discussed
in terms of eigenvalues of the Dirac operator. We show that for QCD at µ 6= 0, which
has a sign problem, the discontinuity in the chiral condensate is due to an alternative to
the Banks-Casher relation. The severity of the sign problem is analyzed in the microscopic
domain of QCD.
§1. Introduction
Starting from its introduction in nuclear physics by Wigner,1) random matrix
theories have been applied to a wide range of problems ranging from the physics of
proteins2) to quantum gravity (see3), 4) for a historical review). Three reasons for
the ubiquity of random matrix theory come to mind. First, eigenvalues of large ran-
dom matrices have universal properties determined by symmetries. Second, random
matrices are models for disorder present in many physical systems. Third, random
matrix theories have a topological expansion which is important for applications to
quantum field theory. One of the attractive features of random matrix theory is that
analytical information can be obtained for complex systems which otherwise only
can be studied experimentally or numerically.
In this review we discuss applications of random matrix theory to QCD at
nonzero temperature and chemical potential. Since the order parameter for the
chiral phase transition5), 6) and the deconfining phase transition7), 8) are determined
by the infrared behavior of the eigenvalues of the Dirac operator, these eigenvalues
are essential for the phase transitions in QCD. Remarkably, the distribution of the
smallest Dirac eigenvalues is given by universal functions9)–13) that depend only on
one or two parameters, the chiral condensate and the pion decay constant. This
offers an alternative way to measure these constants on the lattice.14)–22)
§2. Random Matrix Theory in QCD
Chiral Random Matrix Theory (chRMT) is a theory with the global symmetries
of QCD, but matrix elements of the Dirac operator replaced by random numbers9), 10)
iW † m
, P (W ) ∼ e−NTrW †W . (2.1)
∗) e-mail address: [email protected]
∗∗) e-mail address: [email protected]
http://arxiv.org/abs/0704.0330v1
2 K.Splittorff and J.J.M. Verbaarschot
This random matrix model has the global symmetries and topological properties of
QCD. It is confining in the sense that only color singlets have a nonzero expecta-
tion value. It is now well understood that fluctuations of low-lying eigenvalues of
the Dirac operator are described by chRMT (see23)–28) for lectures and reviews).
Philosphically, this is important because of the realization that chaotic motion dom-
inates the dynamics of quarks at low energy. Practically, this is important because
we can use powerful random matrix techniques to calculate physical observables.
The condition for the applicability of chRMT is that the Compton wavelength
of Goldstone bosons associated with the mass scale z of these eigenvalues is much
larger than the size of the box. With the squared mass of the associated Goldstone
boson given by 2zΣ/F 2π , this condition reads
≪ Λ2. (2.2)
The second condition is necessary to factorize the partition function into a contribu-
tion from the lightest degrees of freedom and all heavier degrees of freedom. These
two conditions determine the microscopic domain of QCD. We stress that z is a scale
in the Dirac spectrum so that, for sufficiently large volumes, we always have eigenval-
ues in the domain (2.2) where eigenvalues fluctuate according to chRMT. This can be
shown rigorously from the following two observations.30), 31) First, the infrared Dirac
spectrum follows from a (partially quenched) chiral Lagrangian determined by chiral
symmetry, and the inequality (2.2) is the condition for factorization of the partition
function into a factor containing the constant modes and another factor containing
the nonzero momentum modes. Second, the factor with the constant modes is equal
to the large N limit of chiral random matrix theory.
In32), 33) the condition (2.2) was imposed on the quark masses and was the bases
for a systematic expansion of the chiral Lagrangian known as the ǫ expansion.
One feature that underlies universal properties of eigenvalues is that they be-
have as repulsive confined charges. This follows from the joint probability distri-
bution ∼
k<l(λ
)2 exp(−N
). It can be shown that eigenvalues
correlations at the micrsocopic scale are universal.34) The reason is spontaneous
symmetry breaking and a mass gap so that they can be described in terms of a
chiral Lagrangian.
2.1. Chiral Random Matrix Theory at µ 6= 0 and T 6= 0
A nonzero temperature does not change the fluctuating behavior of the Dirac
eigenvalues provided that chiral symmetry remains broken. However, a transition to
a different universality class takes place at the critical temperature. A random matrix
model that reproduces this universal behavior of QCD is obtained by replacing the
off-diagonal elements in (2.1) by35)
iW → iW + t, iW † → iW † − t with t = diag(−πT, πT ). (2.3)
This model has been studied elaborately in the literature (see e.g.35)–40)).
A nonzero chemical potential can be introduced analogously to the quark mass.
The requirement is that the small µ behaviour of the QCD partition function should
Random Matrix Theory 3
0.0 1.0 2.0 3.0 4.0
2µ/mπ
m=0.10
m=0.05
m=0.01
Fig. 1. Lattice results for Nc = 2 (taken from
55)) and phase quenched QCD with Nc = 3 (taken
from56))
be reproduced by the random matrix partition function. This achieved by modifying
(2.1) by41)
iW → iW + µ, iW † → iW † + µ, (2.4)
resulting in a nonhermitean Dirac operator with eigenvalues scattered in the complex
plane. The prescription (2.4) is not unique. A random matrix model that has had a
strong impact on recent developments is defined by42)
iW → iW + µH, iW † → iW † + µH with H† = H, (2.5)
where H is drawn from a Gaussian ensemble of random matrices. This model is in
the same universality class as (2.4) but is technically simpler since it can be worked
out by means of the complex orthogonal polynomial method.42)–46)
There are other types of random matrix models that have been applied to QCD.
For example models with random gauge fields such as the Eguchi-Kawai model47) or
its 2-dimensional version.48) QCD in 1 dimension49), 50) is a random matrix model
as well, with universally fluctuating Dirac eigenvalues. Also models with random
Wilson loops51), 52) have attracted significant interest.
§3. Phases of QCD and RMT
QCD-like theories with charged Goldstone bosons have a critical chemical poten-
tial equal to mπ/2. The phase transition to the Bose condensed phase can therefore
be described completely in terms of a chiral Lagragian. At the mean field level,53)
the kinetic terms of this chiral Lagrangian do not contribute, so that these results
can also be obtained from chiral random matrix theory. Indeed, the static part of
the chiral Lagrangian53), 54)
F 2πµ
2Tr[U,B][U †, B]−
ΣTr(MU +MU †). (3.1)
can also be obtained from the large N limit of the models (2.4) or (2.5).
4 K.Splittorff and J.J.M. Verbaarschot
Tricritial point
0.1 0.2 0.3 0.4 0.5 0.6 0.7
Fig. 2. QCD phase diagram in the µTm-space (taken from58))
In Fig. 1 we display lattice results for QCD with Nc = 2
55) and phase quenched
QCD.56) They show an impressive agreement with the results from (3.1) given by
the solid curves in both figures.
3.1. Schematic RMT Phase Diagram
The phase transition in QCD with Nc = 3 at µc = mN/3 cannot be analyzed
by means of chiral Lagrangians. Because of the sign problem lattice studies are not
possible either. In such situation there is long tradition to analyze the same problem
in a much simpler theory in the hope of obtaining at least a qualitative understanding
of the problem. For example, one dimensional QCD,49), 50) or more recently, super
Yang-Mills theory and AdS-CFT duality,57) been explored as toy models for QCD.
We will use random matrix theory at T 6= 0 and µ 6= 0, introduced in (2.3) and
(2.4) to obtain a qualitive understanding of the QCD phase diagram. Lattice QCD
simulations show that the chiral phase transition at µ = 0 is of second order or a
steep cross-over. At T = 0 we expect a first order phase transition at µc = mN/3.
It is natural that the first order line ends in a critical end point or joins the second
order critical line at the tricritical point (see Fig. 3.1, left). This is indeed what
is observed in random matrix theory58), 59) (see Fig. 3.1, right). A similar phase
diagram has also been obtained from the NJL model.60)–62)
Another scenario that was discovered in RMT is the splitting of the first order
line into two at nonzero isospin chemical potential.63) This behavior was also found
in a NJL model64), 65) but might not be stable against flavor mixing interactions.66)
§4. Dirac Spectrum in Theories Without a Sign Problem
Since the spectrum of the Dirac operator determines the chiral condensate, phase
transitions in QCD can be understood in terms of its spectral flow. In this section we
discuss theories with a positive fermion determinant such as QCD with two colors and
phase quenched QCD, where a probabilistic interpretation of the eigenvalue density
is possible. The relation between chiral symmetry breaking and Dirac spectra is
much more complicated when the fermion determinant is complex and its discussion
will be postponed to the next section.
The spectrum of an anti-Hermitean Dirac operator is purely imaginary with an
eigenvalue density that is proportional to the volume. If chiral symmetry is broken
spontaneously, the chiral condensate becomes discontinuous across the imaginary
axis in the thermodynamic limit. Chiral symmetry is restored if such discontinuity
Random Matrix Theory 5
mm m m m m
T < Tc
µ = 0
T > Tc
µ = 0 T < Tc
µ < µc
T < Tc
µ = µc
T < Tc
µ > µc
T > Tc
µ > µc
Fig. 3. Critical behavior of the Dirac spectrum. µc = mπ/2 for T = 0 and increases with T .
is absent for example by the formation of a gap in the Dirac spectrum, see eg.71) .
For µ 6= 0, the Dirac spectrum broadens into a strip of width 4µ2F 2π/Σ.49), 67)
The chemical potential becomes critical when the quark mass hits the edge of this
strip. At this point the chiral condensate starts rotating into a pion condensate.
Chiral symmetry restoration takes place when a gap forms at zero. A schematic
picture of the critical behavior of Dirac eigenvalues is shown in Fig. 3 and the spectral
flow of the Dirac eigenvalues with respect to increasing µ and T is summarized in
Fig. 4. One conclusion from this behavior is that Tc(µ) is a concave function
of µ, and that µc(T ) is a convex function of T . The spectral flow discussed in this
section is supported by lattice simulations at T 6= 0 and µ 6= 0 (See Fig. 5)
4.1. Dirac spectrum in the µ-plane
We could equally well have diagonalized the Dirac operator in a representation
where µγ0 is proportional to the identity,
det(D +m+ µγ0) = det(γ0(D +m) + µ). (4.1)
These eigenvalues are relevant to the baryon number density. A gap in the spectrum
develops at m 6= 0 (see Fig. 6), and the chemical potential becomes critical, µ =
mπ/2 when it hits the inner edge of the domain of eigenvalues.
Increasing µ
Increasing T
Fig. 4. Spectral flow of the Dirac spectrum (left) and phase diagram (right) with respect to µ and
T in phase quenched QCD and QCD with two colors.
6 K.Splittorff and J.J.M. Verbaarschot
1 1.5 2 2.5
b=0.35
b=0.3525
b=0.355
b=0.3575
b=0.36
1.76(t-0.93)
0.0 0.1 0.2 0.3
β=5.5
β=5.66
β=5.71
β=5.75 β=5.9
Fig. 5. Temperature and chemical potential dependence of Dirac eigenvalues. From left to right
taken from.70), 72)–74)
4.2. Quenched Lattice QCD Dirac Spectra at µ 6= 0
Small Dirac eigenvalues at µ 6= 0 have been computed in quenched QCD. The
analytical formulas for the average density of the small Dirac eigenvalues are avail-
able.68), 69) They were first derived68) by exploiting the Toda lattice hierarchy in the
flavor index. Comparisons of random matrix predictions68) for the radial spectral
density and lattice QCD results75), 76) are shown in the left panel of Fig. 7. In other
cases, such as the overlap Dirac operator77) and QCD with Nc = 2,
78) a similar
degree of agreement was found. Both the spectral density and two-point correlations
can be derived from the Lagrangian (3.1), i.e. they are determined by two param-
eters, Fπ and Σ. This can be exploited to extract these low-energy constants. For
example, Fπ and Σ were determined
19), 21) (see also20)) from the correlators shown
in the two right panels of Fig. 7.
§5. Chiral Symmetry Breaking at µ 6= 0
The full QCD partition function at µ 6= 0 which is the average of
det(D +m+ µγ0) = |det(D +m+ µγ0)|eiθ, θ 6= 0, (5.1)
has properties which are drastically different from the phase quenched partition
function where the phase factor is absent. In particular, µc = mN/3 instead of mπ/2,
so that the free energy remains µ-independent until µ = mN/3. For µ < mN/3 the
Fig. 6. Eigenvalues of γ0(D + m) for a random matrix Dirac operator at m = 0 (left), m 6= 0
(middle) (both taken from79)), and lattice QCD at m 6= 0 (right, taken from49)).
Random Matrix Theory 7
—– Splittorff-Verbaarschot-2004
—– Wettig-2004
0 2 4 6 8
−0.15
−0.05
V = 8
10000 configs
µisoFπV
= 0.159
1.27 1.37 1.47 π/2 1.67 1.77 1.87
angle (θ)
lattice: 6
, µa = 0.006
fit: µFV
= 0.14
Fig. 7. The radial spectral density for (left, taken from75), 76)) and two-point correlations (middle
taken from19) and right taken from21)).
chiral condensate remains discontinuous at m = 0, whereas the chiral condensate
of the phase quenched theory approaches zero for m → 0 (see Fig. 5). The only
difference between the phase quenched partition function and the full QCD partition
function is the phase of the fermion determinant. We conclude that the phase factor
is responsible for the discontinuity of the chiral condensate. How can this happen if
for each configuration the support of the spectrum is approximately the same? This
problem known as the “Silver Blaze Problem”80) was solved in.6)
5.1. Unquenched Spectral Density
The spectral density for QCD with dynamical fermions is given by
ρNf (λ) = 〈
δ2(λ− λk)detNf (D +m+ µγ0)〉. (5.2)
Because of the phase of the fermion determinant, this density is in general complex
and can be decomposed as ρNf (λ) = ρNf=0(λ) + ρU (λ). The chiral condensate can
then be decomposed as ΣNf (m) = ΣNf=0(m) +ΣU (m), so that the discontinuity in
Σ(m) is due to ρU . Asymptotically it behaves as
ρU ∼ e
µ2F 2V e
iIm(λ)ΣV
and vanishes outside an ellips starting at Re(λ) = m (see Fig. 9).6) In the right part
of this figure we show the real part of the spectral density for QCD with one flavor
at nonzero chemical potential.
Scatter plot of Dirac eigenvalues
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quark mass m
Support of spectrum
Chiral condensate
condensate
Quenched chiral
in full QCD
µ2F 2
Σ(m) = 1
Fig. 8. Chiral condensate of quenched and full QCD.
8 K.Splittorff and J.J.M. Verbaarschot
Dirac spectrum for Full QCD.
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Oscillating Region
quark mass m
-1000100
0.001
0.002
2F 2µ2
µ2F 2
Fig. 9. Support (left) and real part (right, taken from27)) of Dirac spectral density for QCD with
Nf = 1 and µ 6= 0.
This result explains the mechanism of chiral symmetry breaking at nonzero
chemical potential. The phase of the fermion determinant rotates the pion conden-
sate back into a chiral condensate, but it does so in an unexpected way.6) The same
mechanism is at play for 1d QCD at µ 6= 0.82)
§6. Phase of the Fermion Determinant
The magnitude of the sign problem can be measured by means of the expectation
value of the phase factor of the fermion determiant which can be defined in two ways
〈e2iθ〉Nf =
det(D + µγ0 +m)
det∗(D + µγ0 +m)
detNf (D + µγ0 +m)
, 〈e2iθ〉1+1∗ =
ZNf=2
Z1+1∗
The average 〈· · · 〉 is with respect to the Yang-Mills action. The sign problem is
managable when the average phase factor remains finite in the thermodynamic limit.
In the microscopic domain it is possible to obtain exact analytical expressions for
the average phase factor by exploiting the equivalence between QCD and RMT in
this domain. For µ < mπ/2 the free energy of both QCD and phase quenched
QCD are independent of µ. This does not imply that the average phase factor is
µ-independent. The µ-dependence originates from the charged Goldstone bosons
with mass mπ ± 2µ, and for Nf flavors the mean field result83), 84) for 〈exp(2iθ)〉
reads (1 − 4µ2/m2π)Nf+1. The exact result for the average phase factor for Nf = 2
is shown in Fig. 10 (right), where lattice results85) are also shown (left). The exact
result has an essential singularity at µ = 0, but its thermodyanmic limit agrees with
the mean result.
0 0.5 1 1.5
2µ/mπ
mΣV = 4
mΣV >> 1
Fig. 10. Average phase factor. Lattice QCD results are shown left (taken from85)) and the exact
microscopic result83) is shown right.
Random Matrix Theory 9
§7. Conclusions
The equivalence of chiral random matrix theory and QCD has been exploited
succesfully to derive a host of analytical results. Among others, eigenvalue fluctua-
tions predicted by chRMT have been observed in lattice simulations, the phases of
QCD can be understood in terms of spectral flow, observables can be extracted from
the fluctuations of the smallest eigenvalues, the sign problem is not serious when the
quark mass is outside the domain of the eigenvalues, and mean field results can be
obtained from random matrix theory. Summarizing, chiral random matrix theory is
a powerful tool for analyzing the infrared domain of QCD.
Acknowledgements
The YITP is thanked for its hospitality. G. Akemann, J. Osborn and P.H.
Damgaard are acknowledged for valuable discussions. This work was supported by
US DOE Grant No. DE-FG-88ER40388 (JV), the Villum Kann Rasmussen Foun-
dation (JV), the Danish National Bank (JV) and the Carslberg Foundation (KS).
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|
0704.0331 | Symmetries by base substitutions in the genetic code predict 2' or 3'
aminoacylation of tRNAs | Microsoft Word - MS737.rtf
Manuscript submitted as a Letter to the Editor.
Title:
Symmetries by base substitutions in the genetic code predict 2’ or 3’ aminoacylation of tRNAs.
Authors: Jean-Luc Jestina, Christophe Souléb
Addresses:
aUnité de Chimie Organique, URA 2128 CNRS
Département de Biologie Structurale et Chimie, Institut Pasteur
25 rue du Dr. Roux, 75724 Paris 15, France
email: [email protected] (corresponding author)
tel +33 1 4438 9496; fax +33 1 4568 8404
bInstitut des Hautes Etudes Scientifiques, CNRS
35 route de Chartres, 91440 Bures-sur-Yvette, France
email: [email protected]
Key words :
Mutation; degeneracy; aminoacyl-tRNA synthetase; codon; symmetry breaking.
Understanding why the genetic code is the way it is, has been the subject of numerous
models and still remains largely a challenge (Freeland et al., 2000; Sella and Ardell,
2006). Associations between codons and amino acids were suggested to rely on RNA-
amino acid interactions (Raszka and Mandel, 1972; Yarus, 1998). Closely related
codons were put in correspondence with closely related amino acids within their
biosynthetic pathways (Wong, 2005). Codons have also been grouped into systems
characterized by interlocked thermodynamic cycles (Klump, 2006). Evolutionary
models that minimise the number of the most frequent mutations provide a rationale for
the fact that transitions at the third base of codons are mostly neutral mutations
(Goldberg and Wittes, 1966). Similarly, minimization of the deleterious effects of
sequence-dependent single-base deletions catalyzed by DNA polymerases provides a
rationale for the assignment of stop signals to codons (Jestin and Kempf, 1997). While
in-frame stop codons are strictly selected against, out-of-frame stop codons minimize
the costs of ribosomal slippages (Seligmann and Pollock, 2004). In this context, the
frequencies of codons were found to be highly dependent on the reading frame and
highlighted a symmetrical codon pattern (Koch and Lehmann, 1997). As the genetic
code is quasi-universal among living organisms, models do not need to be time-
dependent, even though time-dependent models have been suggested (Bahi and Michel,
2004; Rodin and Rodin, 2006; Sella and Ardell, 2006). Symmetries in the genetic code
are of special interest as they may highlight underlying organization principles of the
code. A supersymmetric model for the evolution of the genetic code was proposed:
successive breaking of these symmetries would provide an evolutive scenario for the
decomposition into sets of synonymous codons (Hornos and Hornos, 1993; Bashford et
al., 1997). When the amino acids are mapped to the vertices of a 28-gon, three two-fold
symmetries were identified for three subsets of the cognate aminoacyl-tRNA
synthetases (Yang, 2004).
This letter reports complete sets of two-fold symmetries between partitions of the
universal genetic code. By substituting bases at each position of the codons according
to a fixed rule, it happens that properties of the degeneracy pattern or of tRNA
aminoacylation specificity are exchanged.
First the set of sixty-four codons of the genetic code was partitionned in two groups of
thirty-two codons depending on whether the third base of triplets is necessary or not to
define unambiguously an amino acid or a stop signal (property 1). Rumer reported a
symmetry by base substitutions that alters property 1 (Rumer, 1966) . The substitutions
exchanging T and G as well as A and C are applied to all three codon bases and are
called Rumer’s transformation. If the third base is necessary to define an amino acid,
then the symmetrical codon by Rumer’s transformation does not require the third base
of codons to be defined so as to define unambiguously the amino acid. Conversely, if
the third base does not have to be defined so as to define unambiguously an amino acid,
then the symmetrical codon by Rumer’s transformation requires the third base to be
given so as to define unambiguously the amino acid. More recently, one of the authors
reported a symmetry that leaves unchanged property 1 (Jestin, 2006): this symmetry
consists in applying to the first base of codons the substitutions exchanging G and C as
well as T and A. For example, GCN codons coding for alanine are exchanged into CCN
codons coding for proline; for GCN and CCN codons, the third base does not have to
be defined so as to define unambiguously the amino acid.
Here we report a third symmetry that alters property 1 (Fig.1). This symmetry is
obtained by applying successively the two symmetries described above. It consists in
applying the substitution exchanging A and G as well as C and T (a transition) to the
first base in the codon, the substitution exchanging A and C as well as G and T (a
transversion) to the second base in the codon, and the substitution exchanging A and C
as well as G and T (a transversion) in the third base of the codon.
We show further that the only other symmetries exchanging both groups into each other
are obtained by combining the previous ones with a symmetry acting only on the third
base of the codons (here we do not include the substitution on the second base which
exchanges A and C when fixing G and T). This can be seen by counting the number of
occurrences of A, C, G, and T as first, second or third base in a codon of each group.
The result is given in Table 1.
These symmetries are valid for the standard genetic code and for other genetic codes
such as the vertebrate mitochondrial genetic code which has a higher degree of
symmetry of its degeneracy pattern as noted earlier (Lehmann, 2000; Jestin, 2006).
In addition to the existence of Rumer’s transformation, Shcherbak discussed the
following Rumer’s rule (Shcherbak, 1989), which can be read off Table 1: the ratio R =
C+G/T+A of the number of occurrences of C and G by the number of occurrences of T
and A in positions 1, 2 and 3 is equal to 3, 3 and 1 respectively in codons of the first
group (and hence it is 1/3, 1/3 and 1 for codons of the second group). Similarly, the
ratio P = T+C/A+G is 1, 3 and 1 in positions 1, 2 and 3 of the first group of codons.
Secondly, we considered another grouping of codons of the genetic code depending on
whether the amino acids are acylated by amino acyl-tRNA synthetases at the 2’ or at
the 3’ hydroxyl group of the tRNA’s last ribose (property 2) (Sprinzl and Cramer, 1975;
Arnez and Moras, 1994). This classification of amino acyl-tRNA synthetases is very
similar to the one based on sequence homology and on structural considerations (Eriani
et al., 1990; Cusack, 1997). Class I synthetases contain HIGH and KMSKS consensus
sequences, which are absent from class II amino acyl tRNA synthetases. At the
structural level, class I synthetases also contain a Rossman fold, a domain that binds
nucleotides, unlike class II synthetases. Class I enzymes catalyse acylation at the 2’
hydroxyl group of the tRNA while class II enzymes generally catalyse acylation at the
3’ hydroxyl group of the tRNA. PheRS as a class II enzyme that catalyses acylation at
the tRNA’s 2’ hydroxyl group is therefore an exception.
The case of cysteinyl-tRNACys synthetase (CysRS) is ambiguous and was investigated
recently. CysRS is a class I synthetase, but establishes contacts with the major groove
of the acceptor stem of the tRNACys as commonly found for class II enzymes. The
enzyme from Escherichia coli is able to catalyse the acylation reaction at both 2’ and 3’
hydroxyl groups of the tRNACys. The 2’ acylation is about one order of magnitude
faster than the 3’ acylation when catalysed by E. coli cysteinyl-tRNA synthetase in
vitro (Shitivelband and Hou, 2005).
The following classification was then used for 2’ acylated amino acids (Ile, Leu, Met,
Val, Trp, Tyr, Arg, Gln, Glu, Phe) and for 3’ acylated amino acids (His, Pro, Ser, Thr,
Asn, Asp, Lys, Ala, Gly). To the class of 2’ acylated amino acids we also added the
stop signals, a choice partially justified by the fact that two stop codons of the
mitochondrial code of vertebrates code for the 2’ acylated amino acid Arg in the
universal code. Note that if cysteine were not in the class 3’, or if a stop signal was not
in the class 2’, symmetries could not be identified. If cysteine is assigned to the class 2’
as suggested by the previous paragraph, the symmetries are broken. Loss of the
symmetries might have occurred during the evolution of aminoacyl-tRNA synthetases
and might be associated to the late appearance of this amino acid in the genetic code
(Brooks and Fresco, 2002).
When considering molecular properties such as polarity, volume and hydrophobicity,
no statistical differences were noted between class 2’ and class I on one hand, class 3’
and class II on the other hand (Table 3).
There exist two symmetries by base substitutions that exchange the class 2’ with the
class 3’ of the corresponding codon groups (cf. Fig.2). They consist in applying the
substitution exchanging A and C as well as G and T (a transversion) to the first base of
the codon, the substitution exchanging A and G as well as C and T (a transition) to the
second base of the codon, and the substitution exchanging A and C as well as G and T
or A and T as well as C and G (a transversion) to the third base of the codon. These two
symmetries differ by the substitution exchanging A and G as well as C and T in the
third position. They are not related to those depicted in Figures 4 and 5 (Yang, 2004) as
Yang’s three symmetries act only on three subsets of amino acids whereas the
symmetries described herein are valid for the whole codon table.
There are no other symmetries by base substitutions between the two classes 2’ and 3’,
as can be seen by counting the occurrences of A, C, G and T in each class and each
position (Table 2). Note also the following analog of the Rumer’s rule: both the ratio R
= C+G / T+A and the ratio Q = A+C / G+T are equal to 1, 1/3, 1 in positions 1, 2, 3
respectively in the class 2’ (and 1, 3, 1 in the class 3’).
In this letter we have described new symmetries by base substitutions in the genetic
code for partitions concerning the codon degeneracy level or the tRNA-aminoacylation
class. Several evolutionary models have been proposed concerning tRNAs and their
aminoacyl-tRNA synthetases (Martinez Gimenez and Tabares Seisdedos, 2002;
Klipcan and Safro, 2004; Chechetkin, 2006; Di Giulio, 2006). Newly introduced amino
acids may well have been selected to minimize the deleterious effects of
mistranslations, and possibly according to their molecular volumes (Torabi et al.,
2006). A unique serie of binary divisions of the codon table was recently noted: when
the same differentiation rule was applied at each division, the class I / class II pattern
arose consistently (Delarue, 2007). Aminoacyl-tRNA synthetases are likely to have
evolved by gene duplication and mutation of primordial synthetases within each class,
as evidenced by sequence homology (Woese et al., 2000). Consistently, the symmetries
highlighted in this manuscript require three base substitutions per codon, which are
unlikely to happen, thereby shedding some light on the duplication and divergence
mechanism of evolution among the two classes of aminoacyl-tRNA synthetases.
Acknowledgements :
We thank H. Epstein, E. Yeramian, D. Moras, B. Prum and J. Perona for their help.
References :
Arnez, J. G., Moras, D. 1994. Aminoacyl-tRNA synthetase tRNA recognition. Oxford,
IRL Press 61-81.
Bahi, J. M., Michel, C. J. 2004. A stochastic gene evolution model with time dependent
mutations. Bull. Math. Biol. 66, 763-778.
Bashford, J. D., Tsohantjis, I., Jarvis, P. D. 1997. Codon and nucleotide assignments in a
supersymmetric model of the genetic code. Phys. Lett. A 233, 481-488.
Brooks, D. J., Fresco, J. R. 2002. Increased frequency of cysteine, tyrosine, and
phenylalanine residues since the last universal ancestor. Mol. Cell. Proteomics 1,
125-131.
Chechetkin, V. R. 2006. Genetic code from tRNA point of view. J. Theor. Biol. 242, 922-
934.
Cusack, S. 1997. Aminoacyl-tRNA synthetases. Curr. Opin. Struct. Biol. 7, 881-889.
Delarue, M. 2007. An asymmetric underlying rule in the assignment of codons. RNA 13,
161-169.
Di Giulio, M. 2006. The non-monophyletic origin of the tRNA molecule and the origin of
genes only after the evolutionary stage of the last universal common ancestor. J.
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physicochemical distances between amino acids in the evolution of the genetic code.
J. Theor. Biol. 168, 43-51.
Eriani, G., Delarue, M., Poch, O., Gangloff, J., Moras, D. 1990. Partition of tRNA
synthetases into two classes based on mutually exclusive sets of sequence motifs.
Nature 347, 203-206.
Freeland, S. J., Knight, R. D., Landweber, L. F., Hurst, L. D. 2000. Early fixation of an
optimal genetic code. Mol. Biol. Evol. 17, 511-518.
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420-424.
Hornos, J. E. M., Hornos, Y. M. M. 1993. Algebraic model for the evolution of the
genetic code. Phys. Rev. Lett. 71, 4401-4404.
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Jestin, J. L., Kempf, A. 1997. Chain-termination codons and polymerase-induced
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Klipcan, L., Safro, M. 2004. Amino acid biogenesis, evolution of the genetic code and
aminoacyl-tRNA synthetases. J. Theor. Biol. 228, 389-396.
Klump, H. H. 2006. Exploring the energy landscape of the genetic code. Arch. Biochem.
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Martinez Gimenez, J. A., Tabares Seisdedos, R. 2002. On the dimerization of the
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Figure Legends :
Figure 1
Exchange of Group I (codons for which the third base does not have to be defined to
specify the amino acid) into Group II (codons for which the third base must be
defined to specify unambiguously the amino acid or the stop signal) by the
transformation (AG/CT for the first base, GT/AC for the second and third bases).
N=A,T,G or C; H=A,T or C; Y=T or C; R=A or G.
Figure 2
Exchange of the classes 2’ and 3’ by the transformation (AC/GT on the first base,
AG/CT on the second base, AC/GT on the third base). The special case of cysteine is
labelled by an asterisk and discussed in the text.
Table I
Number of occurences of the bases A, C, G and T at each position within the
codon in each group.
Table II
Number of occurences of the bases A, C, G and T at each position within the
codon in each class.
Table III
Statistical t-values computed from the data on hydrophobicity (Kyte and
Doolittle, 1982), molecular volume and polarity (Di Giulio et al., 1994)
comparing the class 2’ with class I, and the class 3’ with class II. These values are
below the threshold of significance given in the Student’s table.
A C G T
Base 1 Group I 4 12 12 4
Group II 12 4 4 12
____________________________
Base 2 Group I 0 16 8 8
Group II 16 0 8 8
____________________________
Base 3 Group I 8 8 8 8
Group II 8 8 8 8
Table 1
A C G T
Base 1 Class 2’ 6 10 6 10
Class 3’ 10 6 10 6
_____________________________
Base 2 Class 2’ 8 0 8 16
Class 3’ 8 16 8 0
_____________________________
Base 3 Class 2’ 10 6 10 6
Class 3’ 6 10 6 10
Table 2
Class 2’ / Class I Class 3’ / Class II
Hydrophobicity 0.07 0.11
Polarity 0.017 0.019
Volume 0.57 0.45
Table 3
|
0704.0333 | Optical properties of the Holstein-t-J model from dynamical mean-field
theory | Optical properties of theHolstein-t-Jmodel fromdynamicalmean-field theory
E. Cappelluti a,b,∗, S. Ciuchi c, S. Fratini d
aDipartimento di Fisica, Università “La Sapienza”, P.le A. Moro 2, 00185 Rome, Italy
bSMC Research Center and ISC, INFM-CNR, v. dei Taurini 19, 00185 Rome, Italy
cINFM and Dipartimento di Fisica, Università dell’Aquila, via Vetoio, I-67010 Coppito-L’Aquila, Italy
dInstitut Néel - CNRS & Université Joseph Fourier, BP 166, F-38042 Grenoble Cedex 9, France
Abstract
We employ dynamical mean-field theory to study the optical conductivity σ(ω) of one hole in the Holstein-t-J model. We provide
an exact solution for σ(ω) in the limit of infinite connectivity. We apply our analysis to Nd2−xCexCuO4. We show that our model
can explain many features of the optical conductivity in this compounds in terms of magnetic/lattice polaron formation.
Key words: magnetic/lattice polarons, spin fluctuations, optical conductivity, cuprates.
PACS: 71.10.Fd, 71.38.-k, 78.20.Bh, 75.30.Ds.
The problem of a single hole in the t-J model interact-
ing also with the lattice degrees of freedom has attracted
recently a notable interest in connection with the physical
properties of the underdoped high-T
cuprates [1,2,3,4]. An
important issue in this regime is the formation of lattice
or magnetic polarons (or both of them) and their mutual
interaction. Along this line, the one-particle properties (as
the effective mass, spectral function, etc.) have been widely
investigated with different techniques. Much less effort has
been however paid to the study of the optical properties.
On the analytical ground, the definition of the optical con-
ductivity (OC) in the single hole is a delicate matter which
needs particular care even for the pure t-J or Holstein
model [5,6]. On the other hand, numerical calculations on
clusters are limited by finite size effects [7]. As a general
rule, thus, the choice of a particular theoretical approach
depends on which property is under examination and on
its feasibility to investigate it.
In this paper we summarize the main results of our work
based on the dynamical mean-field theory (DMFT). Tech-
nical details will be presented in a forthcoming longer pub-
lication [8]. In the infinite coordination number limit z →
∞, we provide an exact solution for σ(ω) as a functional
of the local one-particle Green’s function at finite temper-
ature. It should be stressed that, due to the classical treat-
ment of the magnetic background, the DMFT solution for
∗ Corresponding author. Tel: (+39) 06-49937453 fax: (+39) 06-
49937440
Email address: [email protected] (E. Cappelluti).
0 1 2 3 4
Ref. [7]
this work
λ=1, J/t=0.4, ω
Fig. 1. Comparison between the optical conductivity σ(ω) obtained
by our DMFT solution and Lanczos diagonalization in two dimen-
sions on a finite cluster (Ref. [7]).
z → ∞ is purely local so that it cannot describe the coher-
ent propagation of holes due to the spin fluctuations, nor
the metallic Drude-like peak in σ(ω). On the other hand,
the local properties (as the average number of phonons, size
of the magnetic polaron, etc.) are well captured by this ap-
proach, [9] as well as the incoherent contributions to the
OC. We can explicitly show this feature by comparing in
Fig. 1 our DMFT results with numerical calculations using
Lanczos diagonalization for a single hole in the 2DHolstein-
t-J model on a
10 cluster [7].
The remarkably good agreement of the overall shape as-
sesses the feasibility of our approach to investigate the in-
coherent contributions to the finite frequency OC. This is-
sue is particularly important in light of the intensive de-
bate about the origin of the mid-infrared (MIR) band in the
underdoped high-T
cuprates. Different interpretations for
this feature have been discussed in the literature, involving
Preprint submitted to Elsevier 29 October 2018
http://arxiv.org/abs/0704.0333v1
charge/spin fluctuations, stripe ordering, and other mecha-
nisms. This spread of differentmechanisms reflects the pres-
ence in this doping regime of several actors, which makes
it difficult to isolate each effect from the others. A simpler
and ideal situation is the case of electron-doped cuprates, as
Nd2−xCexCuO4. In these compounds, the long-range anti-
ferromagnetic (AF) order extents up to x ≃ 0.14, so that
the low doping regime x . 0.1 we are interested in, lies well
within the AF phase. On the experimental side, in addi-
tion, a detailed and exhaustive study of the optical conduc-
tivity as a function of temperature T and of the doping x
was recently provided in Ref. [10]. In that work the authors
showed that the low doping OC spectra are characterized
at low temperature by a MIR pseudogap, with an absorp-
tion band edge which varies from EMIR ≃ 0.5− 0.6 for x =
0.05 to EMIR ≃ 0.3− 0.4 for x = 0.1, and is barely distin-
guishable for x = 0.125. Quite interestingly, increasing the
temperature leads to a filling of the pseudogap, rather than
a closing of it. Also remarkable is the temperature depen-
dence of the MIR spectral weight which does not present
any signature at the long-range Néel temperature TN but
rather a kink to a higher “pseudogap” temperature T ∗.
We show here that our approach is able to describe all
these features, and in particular the MIR band edge, in
terms of an optical gap due to the formation of a mag-
netic/lattice polaron. We define T ∗ as the temperature
where the size of the spin polaron becomes larger than
the AF correlation length, that is the maximum tempera-
ture where an injected charge actually probes the magnetic
background. In this perspective we can identify T ∗ with the
mean field Néel temperature of our model, which represents
the temperature above which the system is described by
a paramagnetic state (rather than the onset of long range
order). From Ref. [10] we get for instance T ∗ = 440 K at
x = 0.05 and T ∗ = 200 K at x = 0.125. Using the Curie-
Weiss relation T ∗
= J/4 we estimate respectively J = 152
meV (J/t = 0.126) and J = 69 meV (J/t = 0.057). Note
that such values of J do not represent the bare exchange
interaction but rather the effective spin-exchange coupling
which is reduced by hole doping. We also set ω0 = 84 meV,
consistent with the energy window of the optical phonons
in the cuprates. The electron-phonon (el-ph) coupling con-
stant is fixed to λ = 0.75 in order to reproduce the exper-
imental MIR band edge ≈ 0.5− 0.6 eV in the optical con-
ductivity at x = 0.05, and we assume λ to be independent
of the doping x. Note that with these choices no more free
adjustable parameters remain.
In Fig. 2 we show the temperature evolution of the MIR
optical conductivity for the representative cases x = 0.05
and x = 0.125 (note that in order to compare with the ex-
perimental data of Ref. [10] the tail of a Drude-peak should
be superimposed). Most remarkable is the behavior of σ(ω)
at low temperature, which shows a well defined gap for
x = 0.05 while no gap is found for x = 0.125. This fea-
ture reflects the formation of the lattice polaron and its in-
terplay with the spin degrees of freedom. While the el-ph
coupling λ = 0.75 alone is not strong enough at x = 0.125
0 0.5 1
ω [eV]
0.5 1 1.5
ω [eV]
0 200 400
T [K]
x=0.05
x=0.125
T=50K
T=440K
T=540K
T=540K
T=340K
T=50K
T=190K
Fig. 2. Temperature dependence of the optical conductivity σ(ω) for
x = 0.05 and x = 0.125. Solid lines are used for T ≤ T ∗, dashed
lines for T > T ∗. Inset: loss of the MIR spectral weight ∆Neff , as
defined in Ref. [10], as function of T for x = 0.05 (filled circles) and
x = 0.125 (empty squares). Arrows mark the corresponding T ∗.
(J/t = 0.057) to establish a spin/lattice polaron, the lo-
calization effects induced by the larger exchange coupling
J/t = 0.126 at x = 0.05 favor the lattice polaron forma-
tion. This leads thus to the opening of an optical gap in
σ(ω) (this key point will be extensively discussed in a forth-
coming publication[8]). Increasing T reduces the localiza-
tion effects induced by the magnetic ordering. This makes
the positive interplay with the el-ph coupling less effective,
leading to a progressive filling of the pseudogap. Note that
this effect disappears in the disordered magnetic case for
T > T ∗, and further increasing of T leads to a reduction of
the MIR optical conductivity which is spread on a larger
energy window. This is reflected in the characteristic tem-
perature behavior of the MIR spectral weight ∆Neff , as de-
fined in Ref. [10], which presents a kink at T ∗ (inset of Fig.
2)[11].
References
[1] A.S. Mishchenko and and N. Nagaosa, Phys. Rev. Lett. 93
(2004) 0236402; Phys. Rev. B 73 (2006) 092502.
[2] O. Rösch and O. Gunnarsson, Phys. Rev. Lett. 92 (2004) 146403;
Eur. Phys. J. B 43 (2005) 11.
[3] O. Gunnarsson and O. Rösch, Phys. Rev. B 73 (2006) 174521.
[4] P. Prelovšek, R. Zeyher, and P. Horsch, Phys. Rev. Lett. 96
(2006) 086402.
[5] M.P.H. Stumpf and D.E. Logan, Eur.Phys.J.B, 8 (1999) 377.
[6] S. Fratini and S. Ciuchi, Phys. Rev. B 74 (2006) 075101.
[7] B. Bäuml et al., Phys. Rev. B 58 (1998) 3663.
[8] E. Cappelluti, S. Ciuchi and S. Fratini, in preparation (2007).
[9] E. Cappelluti and S. Ciuchi, Phys. Rev. B 66 (2002) 165102.
[10] Y. Onose et., Phys. Rev. B 69 (2004) 024504.
[11] Since we do not find any isosbestic point in our calculations,
we use the experimental energy windows of Ref. [10] to define
∆Neff , namely ωmin = 0.12 eV, ωmax = 0.42 eV for x = 0.05
and ωmax = 0.21 eV for x = 0.125.
References
|
0704.0334 | A Multiphilic Descriptor for Chemical Reactivity and Selectivity | Microsoft Word - LA_Multiphilic_3-4-7.doc
A Multiphilic Descriptor for Chemical Reactivity
and Selectivity
J. Padmanabhan1,2, R. Parthasarathi2, M. Elango2, V. Subramanian2,*, B. S.
Krishnamoorthy1,3, S. Gutierrez-Oliva4, A. Toro-Labbé4,*, D. R. Roy1 and P. K.
Chattaraj1,*
1Department of Chemistry, Indian Institute of Technology, Kharagpur 721302, India.
2Chemical Laboratory, Central Leather Research Institute, Adyar, Chennai 600 020, India.
3School of Chemistry, Bharathidasan University, Tiruchirappalli-620 024, India.
4Laboratorio de Química Teórica Computacional (QTC), Facultad de Química,
Pontificia Universidad Católica de Chile, Casilla 306, Correo 22, Santiago, Chile.
Abstract
In line with the local philicity concept proposed by Chattaraj et al. (Chattaraj, P. K.;
Maiti, B.; Sarkar, U. J. Phys. Chem. A. 2003, 107, 4973) and a dual descriptor derived by
Toro-Labbé and coworkers (Morell, C.; Grand, A.; Toro-Labbé, A. J. Phys. Chem. A.
2005, 109, 205), we propose a multiphilic descriptor. It is defined as the difference
between nucleophilic (ωk+) and electrophilic (ωk-) condensed philicity functions. This
descriptor is capable of simultaneously explaining the nucleophilicity and electrophilicity
of the given atomic sites in the molecule. Variation of these quantities along the path of a
soft reaction is also analyzed. Predictive ability of this descriptor has been successfully
tested on the selected systems and reactions. Corresponding force profiles are also
analyzed in some representative cases. Also, to study the intra- and intermolecular
reactivities another related descriptor namely, the nucleophilicity excess ( ∓
gωΔ ) for a
nucleophile, over the electrophilicity in it has been defined and tested on all-metal aromatic
compounds.
*Authors for correspondence:
E-mail: [email protected], [email protected], [email protected],
1. Introduction
The understanding of chemical reactivity and site selectivity of the molecular
systems has been effectively handled by the conceptual density functional theory (DFT).1
Chemical potential, global hardness, global softness, electronegativity and electrophilicity
are global reactivity descriptors, highly successful in predicting global chemical reactivity
trends. Fukui function (FF) and local softness are extensively applied to probe the local
reactivity and site selectivity. The formal definitions of all these descriptors and working
equations for their computation have been described. 1-4 Various applications of both
global and local reactivity descriptors in the context of chemical reactivity and site
selectivity have been reviewed in detail.3
Parr et al. introduced the concept of Electrophilicity (ω) as a global reactivity index
similar to the chemical hardness and chemical potential. 5 This new reactivity index
measures the stabilization in energy when the system acquires an additional electronic
charge ΔN from the environment. The electrophilicity is defined as
ημω 2/2= (1)
In Eq. (1), μ ≈ -(I+A)/2 and η ≈ (I-A)/2 are the electronic chemical potential and the
chemical hardness of the ground state of atoms and molecules, respectively, approximated
in terms of the vertical ionization potential (I) and electron affinity (A). The
electrophilicity is a descriptor of reactivity that allows a quantitative classification of the
global electrophilic nature of a molecule within a relative scale. 5
Fukui Function (FF) 6 is one of the widely used local density functional descriptors
to model chemical reactivity and site selectivity and is defined as the derivative of the
electron density ρ ( r ) with respect to the total number of electrons N in the system, at
constant external potential ν ( r ) acting on an electron due to all the nuclei in the system
[ ] [ ] )()()()( rvN Nrrvrf ∂∂== ρδδμ . (2)
The condensed FF are calculated using the procedure proposed by Yang and
Mortier,7 based on a finite difference method
)()1( NqNqf kkk −+=
+ for nucleophilic attack (3a)
)1()( −−=− NqNqf kkk for electrophilic attack (3b)
[ ] 2)1()1( −−+= NqNqf kkok for radical attack (3c)
where kq is the electronic population of atom k in a molecule.
Chattaraj et al.8 have introduced the concept of generalized philicity. It contains
almost all information about hitherto known different global and local reactivity and
selectivity descriptors, in addition to the information regarding electrophilic/nucleophilic
power of a given atomic site in a molecule. It is possible to define a local quantity called
philicity associated with a site k in a molecule with the help of the corresponding
condensed- to- atom variants of FF, αkf as
αα ωω kk f= (4)
where (α= +, - and 0) represents local philic quantities describing nucleophilic,
electrophilic and radical attacks respectively. Eq. (4) predicts that the most electrophilic
site in a molecule is the one providing the maximum value of ωk+. When two molecules
react, which one will act as an electrophile (nucleophile) will depend on, which has a
higher (lower) electrophilicity index. This global trend originates from the local behavior
of the molecules or precisely at the atomic site(s) that is(are) prone to electrophilic
(nucleophilic) attack. Recently the usefulness of electrophilicity index in elucidating the
toxicity of polychlorinated biphenyls, benzidine and chlorophenol has been assessed in
detail. 9-11
In addition to the knowledge of global softness (S), which is the inverse of
hardness, 12 different local softnesses 13 used to describe the reactivity of atoms in
molecule, can be defined as
k ks Sf
α α= (5)
where (α= +, - and 0) represents local softness quantities describing nucleophilic,
electrophilic and radical attacks respectively. Based on local softness, relative
nucleophilicity (sk- /sk+) and relative electrophilicity (sk+ /sk-) indices have also been defined
and their usefulness to predict reactive sites also been addressed to.14 It has been
established that the quantum chemical model selected to derive wave function; population
scheme used to obtain the partial charges and basis set employed in the molecular orbital
calculations are important parameters, which significantly influence the FF values. 15-18
The condensed philicity summed over a group of relevant atoms is defined as the
“group philicity”. It can be expressed as19
αα ωω
where n is the number of atoms coordinated to the reactive atom, αωk is the local
electrophilicity of the atom k, and ωgα is the group philicity obtained by adding the local
philicity of the nearby bonded atoms. In this study19 the group nucleophilicity index (ωg+)
of the selected systems is used to compare the chemical reactivity trends.
Toro-Labbé et al20 have recently proposed a dual descriptor (Δf ( r )), which is
defined as the difference between the nucleophilic and electrophilic Fukui functions and is
given by,
Δf(r) = [ (f +(r) - (f - (r) ] (7)
If Δf(r) > 0, then the site is favored for a nucleophilic attack, whereas if Δf (r) < 0, then the
site may be favored for an electrophilic attack. The associated dual local softness have also
been defined as,19
Δsk = S (fk+ - fk-) = (sk+ - sk-) (8)
It is defined as the condensed version of Δf (r) multiplied by the molecular softness S.
2. Multiphilic Descriptor
In the light of the local philicity concept proposed by Chattaraj et al.8 and the dual
descriptor derived by Toro-Labbé and coworkers,20 we propose a multiphilic descriptor
using the unified philicity concept, which can concurrently characterize both nucleophilic
and electrophilic nature of a chemical species. It is defined as the difference between the
nucleophilic and electrophilic condensed philicity functions. It is an index of selectivity
towards nucleophilic attack, which can as well characterize an electrophilic attack and is
given by,21
Δωk = [ωk+ - ωk- ] = ω [Δƒk] (9)
where Δƒk is the condensed-to-atom variant-k of Δƒ(r) (eq 7). If Δωk > 0, then the
site k is favored for a nucleophilic attack, whereas if Δωk < 0, then the site k may be favored
for an electrophilic attack. Because FFs are positive (0 < ƒk < 1), -1 < Δƒk < 1, and the
normalization condition for Δωk is
0=Δ=Δ ∑∑
k fωω (10)
Although Δωk and Δfk will contain the same intramolecular reactivity information
the former is expected to be a better intermolecular descriptor because of its global
information content.
We may analyze the nature of ( )rωΔ in terms of that22 of ( )f rΔ as follows:
[ ]( )( )
ωω ⎛ ⎞∂∂⎛ ⎞
= ⎜ ⎟⎜ ⎟∂ ∂⎝ ⎠ ⎝ ⎠
∂ ∂⎛ ⎞ ⎛ ⎞
= +⎜ ⎟ ⎜ ⎟∂ ∂⎝ ⎠ ⎝ ⎠
( ) ( )
f r f r
= + Δ⎜ ⎟∂⎝ ⎠
( ) ( )
f r r
= + Δ⎜ ⎟∂⎝ ⎠
( ) ( )
r f r
⎡ ⎤∂ ∂⎛ ⎞ ⎛ ⎞
Δ = −⎢ ⎥⎜ ⎟ ⎜ ⎟∂ ∂⎝ ⎠ ⎝ ⎠⎣ ⎦
The multiphilicity descriptor, ( )rωΔ is a measure of the difference between local
and global (modulated by ( )f r ) reactivity variations associated with the electron
acceptance/ removal. Incidentally, the variation of
ω∂⎛ ⎞
⎜ ⎟∂⎝ ⎠
across the periodic table is
similar to that of μ.23
2v vN N
⎡ ⎤∂ ∂⎛ ⎞
=⎜ ⎟ ⎢ ⎥∂ ∂⎝ ⎠ ⎣ ⎦
24 vv N
μ μ μ η
η η η
⎛ ⎞∂ ∂⎛ ⎞
= −⎜ ⎟ ⎜ ⎟∂ ∂⎝ ⎠⎝ ⎠
= −
μ γ μ
= − = −
Since γ is generally very small,24
ω∂⎛ ⎞
⎜ ⎟∂⎝ ⎠
is expected to follow the μ trend.
Problems associated with the definition of η and the discontinuity25 in E as a
function of N will be present in the ( )f rΔ definition and the discontinuity in ( )rρ .
Similar type of differentiation has also been attempted by other research workers.26
Also, to study the intra- and intermolecular reactivities another related descriptor
namely, nucleophilicity excess ( ∓
gωΔ ) for a nucleophile, over the electrophilicity (net
nucleophilicity) in it is defined as
( )+−+− −=−=Δ ggggg ffωωωω ∓ (11)
where )(
ωω and )(
ωω are the group philicities of the
nucleophile in the molecule due to electrophilic and nucleophilic attacks respectively. It is
expected that the nucleophilicity excess ( ∓
gωΔ ) for a nucleophile should always be
positive whereas it will provide a negative value for an electrophile in a molecule.
In the present study, we use both the multiphilicity descriptor and nucleophilicity
excess to probe the nature of attack/reactivity at a particular site in the selected systems.
3. Computational Details
The geometries of HCHO, CH3CHO, CH3COCH3, C2H5COC2H5, CH2=CHCHO
CH3CH=CHCHO, NH2OH, CH3ONH2, CH3NHOH, OHCH2CH2NH2, CH3SNH2,
CH3NHSH, SHCH2CH2NH2 and all-metal aromatic molecules, viz., MAl4– (M=Li, Na, K
and Cu) are optimized by B3LYP/6-311+G** as available in the GAUSSIAN 98 package.27
Various reactivity and selectivity descriptors such as chemical hardness, chemical
potential, softness, electrophilicity and the appropriate local quantities employing natural
population analysis (NPA)28, 29 scheme are calculated. HPA scheme (Stockholder
Partitioning Scheme) 30 as implemented in the DMOL3 package 31 has also been used to
calculate the local quantities employing BLYP/DND method. For all-metal aromatic
molecules, ∆SCF method has been utilized to compute the ionization potential (IP) and
electron affinity (EA) according to the equations (I=EN-1 - EN, A=EN - EN+1, where I and A
are obtained from total electronic energy calculations on the N-1, N, N+1-electron systems
at the neutral molecule geometry).
4. Results and Discussion
A series of carbonyl compounds is selected in the present study to probe the
usefulness of the multiphilicity descriptor (Figure 1). A comparison with various other
descriptors and the recently derived dual descriptor is also probed. Due to bipolar nature of
C=O bond, both nucleophilic and electrophilic attacks are possible at C and O sites. It is
noted that the rate of nucleophilic addition on the carbonyl compound be reduced by
electron donating alkyl groups and enhanced by electron withdrawing ones. 32 Recently,
we have studied a set of these carbonyl compounds in the light of philicity and group
philicity.19 The global molecular properties of the selected series of carbonyl compounds
are presented in Table 1. Various local quantities for particular sites of the selected systems
are listed in Table 2 and Table 3. Selected compounds are grouped into two sets namely,
nonconjugated and α, β-conjugated carbonyl compounds.
For the nonconjugated carbonyl compounds, the carbon atom (C1) bearing the
carbonyl group is expected to be the most reactive site towards a nucleophilic attack. Table
2 lists the values of local reactivity descriptors using B3LYP/6-311+G** method for NPA
derived charges of the selected molecules. NPA derived local quantities predict the
expected maximum value for carbonyl carbon (C1) of all the selected molecules for fk+, sk+
and ωk+. But sk+/sk- is unable to provide the maximum value for C1 atom due to negative FF
values. One important point to note is that among the descriptors fk+, sk+, ωk+ and sk+/sk-,
+ value is capable of providing a clear distinction between carbonyl carbon (C1) and the
oxygen site for nucleophilic attack.
Since, HPA derived charges generally provide non-negative FF values, we also
made use of it for local reactivity analysis on carbonyl compounds. HPA derived local
reactivity descriptors also predict the expected maximum value for C1 atom in the case of
HCHO and CH3CHO but fails to predict for CH3COCH3 and C2H5COC2H5, where oxygen
atom is shown to be prone towards nucleophilic attack. Nevertheless, the fk+ value of
oxygen is almost same as that of carbonyl carbon (C1), thus making it difficult to make a
clear decision on the electrophilic behavior of these atoms. Under these situation, dual
descriptors Δf (r), Δs k and multiphilic descriptor Δω (r), give a helping hand. All these
quantities provide a clear difference between nucleophilic and electrophilic attacks at a
particular site with their sign. That is, they provide positive value for site prone for
nucleophilic attack and a negative value at the site prone for electrophilic attack. The
advantage of multiphilic descriptor Δω (r) is that they provide higher value in terms of
magnitude compared to other dual descriptors. For instance, values of Δf(r), Δsk and Δω(r)
for nucleophilic (electrophilic) attack at carbonyl carbon (oxygen) site of CH3CHO are
1.06 (-0.93), 0.17 (-0.15), 3.03 (-2.65) respectively for NPA derived charges. Almost the
same trend is followed in the case of HPA derived charges.
The second group of compounds namely, α, β-conjugated carbonyl is elaborately
studied in the recent past because of the presence of two reactive centers.33 The first
reactive site is the carbon (C1) of the carbonyl, and the second is the carbon in the β
position (C6). In such a case, the β carbon is activated because of the withdrawing
mesomeric effect of the adjacent carbonyl group. As seen from Table 2 and Table 3, NPA
derived charges give a maximum value for fk+ to carbonyl carbon whereas HPA derived
charges provide maximum fk+ value to the β carbon atom (C6) in the case of CH2=CHCHO
molecule. For CH3CH=CHCHO, NPA (HPA) provide maximum fk+ value of 0.44 (0.17) to
carbonyl carbon (C1) compared to the β carbon site of 0.34 (0.16). This ambiguous
behavior may be due to the dependence of local reactivity descriptors on the selection of
basis set and population schemes. Further oxygen site shows high value for fk+ and other
local descriptors, making it difficult to predict the proper electrophilic site. Even now Δω
(r) exhibits high positive value on both carbons that are supposed to be electrophilic and a
high negative value on the oxygen site disclosing clearly its nucleophilic character
compared to other dual descriptors. Also it can be noted from Tables 2 and 3 that, even for
molecules with more than one reactive sites, Δω (r) is capable of making a clear distinction
among them in terms of their magnitude. That is, for molecules 6 and 7 having two
reactive sites as carbon (C1) of the carbonyl and the carbon in the β position (C6), our
descriptors are capable of distinctly identifying the stronger site
(electrophilic/nucleophilic).
Optimized structures along with atom numbering for the selected set of amines are
presented in Figure 2. Global and local reactivity properties of the selected set of amines
calculated using B3LYP/6-311+g** and BLYP/DND methods are presented in Tables 4 to
6. Global reactivity trend based on ω, is given by
B3LYP/6-311+g** method (Table 4)
(i) CH3ONH2 > OHCH2CH2NH2 > CH3NHOH > NH2OH
(ii) CH3NHSH > SHCH2CH2NH2 > CH3SNH2
BLYP/DND method (Table 4)
(i) CH3ONH2 > OHCH2CH2NH2 > NH2OH > CH3NHOH
(ii) CH3NHSH > SHCH2CH2NH2 > CH3SNH2
Though both the methods show variation in reactivity trend for oxygen containing
systems, trends related to sulfur containing systems are same.
Based on NPA and HPA charge derived multiphilic descriptor at nitrogen site
(∆ωN), following reactivity trend has been obtained,
NPA (Table 5)
(1) OHCH2CH2NH2 > CH3NHOH > NH2OH > CH3ONH2
(2) CH3NHSH > SHCH2CH2NH2 > CH3SNH2
HPA (Table 6)
(1) OHCH2CH2NH2 > CH3ONH2 > NH2OH > CH3NHOH
(2) CH3NHSH > SHCH2CH2NH2 > CH3SNH2
It may be noted that trends are same as ω for sulfur containing systems, but shows
variations with respect to oxygen containing systems for both NPA and HPA charge
derived ∆ωN.
So for as the intramolecular reactivity trends are concerned, site with maximum
negative value of ∆ωk is the most preferred site for electrophilic attack. Chemical intuition
suggests that N site is more prone towards electrophilic attack. Table 7 lists the site with
maximum negative value for ∆ωk for the selected set of amines. It is seen that with a few
exception, N site is predicted as the most preferred site for electrophilic attack.
Further in order to test ∆ωk along intrinsic reaction coordinate (IRC), we consider a
cope rearrangement of hexa-1,5-diene. This is an example of [3,3] sigmatropic reaction.
Figure 3 provides the optimized geometrical structures with atom numbering for the
reactant, transition state and product calculated using B3LYP/6-31G* level of theory.
Table 8 gives the global reactivity parameters of the reactant, transition state and product.
As expected, hardness is minimum (2.48 eV) and the corresponding electrophilicity index
is maximum (1.57 eV) at the transition state. Variation of global reactivity parameter along
the IRC path is presented in Table 9 and Figure 4 (a-b). Variation of energy (E) and ω
along IRC path is given in Figure 5a. It is seen that both E and ω are maximum around the
transition state indicating it as the most unstable structure along the IRC path. Figure 5 b
provides the variation of hardness (η) and polarizability (α) along the IRC path. An inverse
relationship exists between them. That is, η reaches a minimum whereas α becomes
maximum at the transition state as expected.
Variation of multiphilic descriptor (∆ωk) along IRC for the important atomic sites
(C1 and C3/ C6 and C11) is presented in Figure 5. In going from reactant to product, C1 and
C3 (C6 and C11) sites change their nature and become more prone towards electrophilic
attack (nucleophilic attack) at the product side. This change in the nature of attack takes
place around the transition state.
In studying the importance of nucleophilicity excess ( ∓
gωΔ ) descriptor, a careful
analysis on the electronic structure, property and reactivity of all-metal aromatic
compounds, viz., MAl4– (M=Li, Na, K and Cu) is performed. The four membered
aluminum unit Al4 present in all the molecules may be considered as a single unit. This
unit can easily take part in charge transfer process with the M (≡Li, Na, K, Cu) atom in
those complexes.
Figure 6 shows the various stable isomers of MAl4–. The C4v isomer of the MAl4– is
reported as energetically most stable, least polarizable and hardest.34, 35 Table 10 presents
the group philicity (ωg+, ωg–) values of the Al42– nucleophile and M+ (M=Li, Na, K, Cu)
electrophile in the MAl4– isomers. It is found that in all MAl4– isomers the nucleophilicity
of the Al42– aromatic unit overwhelms its electrophilic trend (i.e. +−
gg ωω ) and therefore
gωΔ is positive, whereas the electrophilicity of M
+ dominates over its nucleophilicity (i.e.
gg ωω ) and therefore
gωΔ is negative as expected. It is important to note that
gωΔ of
Al42– is maximum in the case of most stable C4v isomer of the MAl4– molecule. The order
of the ∓
gωΔ value of Al4
2– nucleophile in MAl4–,
vvv CCC ∞24 , i.e. stabilization of an
MAl4– isomer (except in KAl4–) increases its nucleophilicity and accordingly can be used
as a better molecular cathode. It is also important to note that the nucleophilicity of the
Al42– unit in MAl4– (C4v) increases as K Cu Na Li≺ ≺ ≺ according to the respective
nucleophilicity excess values. Standard expressions1-5 for ∆N and ∆E in terms of group
electronegativity and group hardness will provide additional insights into the electron
transfer process.
Variation of kωΔ along the IRC of three selected reactions,
36 viz., a) a
thermoneutral reaction: Fa– + CH3-Fb → Fa-CH3 + Fb–, b) an endothermic reaction: HNO
→ HON, c) an exothermic reaction: H2OO → HOOH is provided in figures 7 (a) – 7(c).
For the thermoneutral reaction, both the Fa– (bond making) and Fb– (bond breaking) are
nucleophilic. The net nucleophilicity of the Fa– atom is more than that of the Fb– atom along
the IRC from reactant side to TS and the situation is reversed for the IRCs pertaining to the
TS to product side. For the endothermic reaction, the net nucleophilicity of O (bond
making) is higher than that of N (bond breaking) along the IRC. In the case of exothermic
reaction, the O1 (bond making) atom is more electrophilic than its nucleophilic activity.
Moreover, its Fukui function values as calculated through Mulliken Population Analysis
(MPA) scheme become negative in some cases. For the thermoneutral reaction kωΔ is
minimum at the transition state. For other two reactions, kωΔ does not always follow the
trend that the IRC corresponding to the minimum value of kω
± (if not zero) is in
accordance with the Hammond’s postulate.36 Figures 8 (a) – 8 (c) provide the profiles for
the corresponding reaction forces.37
Apart from the important points corresponding to the reactant (R), the transition
state (TS) and the product (P) there exists two other important points associated with the
configurations having the force maximum (Fmax) and the force minimum (Fmin). The
zeroes, maxima and minima of the reaction force define key points along the reaction
coordinate, which divide it into three reaction regions that are identified through vertical
dashed lined in Figure 8. The first stage, in the reactant region, tends to be preparative in
nature with emphasis in structural effects such as rotation, bond stretching, angle bending,
etc., that will facilitate subsequent steps. The transition state region is mostly characterized
by electronic rearrangements whereas the product region is mainly associated to structural
relaxation necessary to reach the products. We have shown that analyzing a chemical
reaction in terms of these regions can provide significant insight into its mechanism and
the roles played by external factors, such as external potentials and solvents.37, 38 Partition
of the activation energies in terms of the work done in going from i) R to Fmin: W1, ii) Fmin
to TS: W2, iii) TS to Fmax: W3 and iv) Fmax to P: W4 gives the activation energy for the
forward reaction (Ef#) as (W1+W2) and that of the reverse reaction (Er#) as -(W3+W4).
Therefore the reaction energy (∆E0) becomes (Ef# – Er# = W1+W2+W3+W4). These values
are provided in Table 11. As expected ∆E0 is zero, negative and positive for the
thermoneutral, exothermic and endothermic reactions respectively. The skew-symmetric
nature of the force profile for the thermoneutral reaction suggests that A=W1+W4 and
B=W2+W3 would be zero. Similarly A, B would be positive (negative) for the
endo(exo)thermic reactions. The transition state at the IRC=0 configuration lies at the
middle between Fmax and Fmin configurations for the thermoneutral reaction whereas it lies
towards the Fmin(Fmax) configurations for the exo(endo)thermic reaction, a signature of the
Hammond postulate via reaction force.
Similar values of W1 and W2 (see Table 11) together with the changes observed in
the nucleophilicity along the reaction coordinate for the thermoneutral SN2 substitution and
for the exothermic reaction H2OO → HOOH indicate that structural and electronic
reordering show up at the very beginning of the reaction, 37,38 through a sharp decrease of
the nucleophilicity, this change practically ceases at the transition state of the exothermic
reaction to reach the product value. It is interesting to note that in both cases the lowering
of nucleophilicity of the key atoms from the reactants (Δω(Fa/Fb) ~ 0.014; Δω(O1) ~ 0.14)
to the transition state (Δω(Fa/Fb) ~ 0.004; Δω(O1) ~ 0.0) requires a similar amount of
energy (9.54 kcal/mol and 7.39 kcal/mol, respectively). It can be observed in Table 11 that
for the thermoneutral reaction W1>W2 indicating that the preparation step requires more
energy than the transition to product step. On the other hand, the W2 values for the
thermoneutral and exothermic reactions are quite close to each other and the work W1
associated to the preparation step in the thermoneutral reaction is larger than that of the
exothermic reaction, this indicates that in the SN2 reaction the structural reordering of the
CH3 group to reach the D3h structure at the transition state is the key transformation that
involve most of the activation energy. In the endothermic HNO → HON reaction the small
changes of nucleophilicity together with large values of W1 and W2 indicates that the
reaction is mainly driven by the structural reordering in the preparation step.
5. Conclusions
A multiphilicity descriptor (Δωk) is proposed and tested in this work. It is shown
that, Δωk helps in identifying the electrophilic/nucleophilic nature of a specific site within
a molecule. A comparison between different local reactivity descriptors is carried out on a
set of carbonyl compounds. Also a selected set of amines is analyzed using Δωk. Further,
we also consider a cope rearrangement of hexa-1,5-diene to test the variation of Δωk along
IRC path. It is seen that Δωk presents a clear distinction between electrophilic and
nucleophilic sites within a molecule in terms of their magnitude and sign. Hence they
reveal the fact that multiphilic descriptor can effectively be used in characterizing the
electrophilic/nucleophilic nature of a given site in a molecule. Also the importance of
nucleophilicity excess ( ∓
gωΔ ) descriptor on the reactivity of all-metal aromatic
compounds, viz., MAl4– (M=Li, Na, K and Cu) is successfully analyzed. Important insight
into three different types of reactions, viz., a) thermoneutral, b) endothermic and c)
exothermic are obtained through the analysis of the multiphilic descriptor profiles within
the reaction regions defined by reaction force along the reaction path.
The results discussed so far clearly show the importance of the selected descriptors,
namely, multiphilic descriptor and nucleophilicity excess in analyzing the overall reactivity
trends in molecular systems.
Acknowledgment:
PKC and DRR thank BRNS, Mumbai for financial assistance. JP and BSK thank the IIT
Kharagpur for providing the facilities required for a summer project. JP also thanks the
UGC for selecting him to carryout his Ph.D. work under FIP. ATL and SGO wish to thank
financial support from FONDECYT, grant N° 1060590, FONDAP through project N°
11980002 (CIMAT) and Programa Bicentenario en Ciencia y Tecnología (PBCT),
Proyecto de Inserción Académica N° 8. ATL is also indebted to the John Simon
Guggenheim Foundation for a fellowship.
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TABLE 1: Calculated Global Reactivity Properties of the Selected Molecules using
B3LYP/6-311+g** and BLYP/DND method.
η μ ω S
η μ ω S
Molecules B3LYP/6-311+g** (eV) BLYP/DND (eV)
HCHO 2.960 -4.707 3.742 0.169 1.942 -4.260 4.673 0.258
CH3CHO 3.115 -4.224 2.864 0.161 2.096 -3.791 3.425 0.238
CH3COCH3 3.144 -3.910 2.432 0.159 2.133 -3.456 2.800 0.234
C2H5COC2H5 3.153 -3.799 2.288 0.159 2.151 -3.367 2.635 0.233
CH2=CHCHO 2.503 -4.904 4.805 0.200 1.545 -4.413 6.303 0.324
CH3CH=CHCHO 2.542 -4.631 4.217 0.197 1.593 -4.132 5.359 0.314
TABLE 2: Calculated Local Reactivity Properties of the Selected Molecules using B3LYP/6-311+g** method for NPA derived
charges.
Molecule fk
- Δfk
+- fk
HCHO C 0.8323 -0.1722 0.1406 -0.0291 -4.8331 3.1146 -0.6444 1.0045 0.1697 3.7591
O 0.0399 0.9409 0.0067 0.1589 0.0424 0.1494 3.5211 -0.9010 -0.1522 -3.3718
CH3CHO C1 0.8178 -0.2416 0.1313 -0.0388 -3.3856 2.3419 -0.6917 1.0593 0.1700 3.0337
O 0.0072 0.9320 0.0012 0.1496 0.0077 0.0206 2.6691 -0.9250 -0.1484 -2.6485
CH3COCH3 C1 0.3142 -0.2916 0.0500 -0.0464 -1.0772 0.7640 -0.7092 0.6058 0.0964 1.4732
O -0.2540 0.9286 -0.0404 0.1477 -0.2734 -0.6170 2.2582 -1.1820 -0.1881 -2.8755
C2H5COC2H5 C1 0.3064 -0.2944 0.0486 -0.0467 -1.0408 0.7011 -0.6736 0.6007 0.0953 1.3746
O -0.2650 0.8751 -0.0420 0.1388 -0.3024 -0.606 2.0025 -1.1400 -0.1807 -2.6080
CH2=CHCHO C6 0.2789 0.2070 0.0557 0.0413 1.3472 1.3402 0.9944 0.0719 0.0144 0.3458
C1 0.4355 -0.2288 0.0870 -0.0457 -1.9033 2.0926 -1.0995 0.6643 0.1327 3.1921
O -0.0560 0.9265 -0.0112 0.1851 -0.0605 -0.2700 4.4518 -0.9830 -0.1963 -4.7213
CH3CH=CHCHO C6 0.3437 0.0926 0.0676 0.0182 3.7143 1.4494 0.3904 0.2511 0.0494 1.0590
C1 0.4408 -0.2365 0.0867 -0.0465 -1.8642 1.8592 -0.9973 0.6773 0.1332 2.8566
O -0.0670 0.9281 -0.0132 0.1825 -0.0721 -0.2820 3.9142 -0.9950 -0.1957 -4.1964
TABLE 3: Calculated Local Reactivity Properties of the Selected Molecules using BLYP/DND method for HPA derived
charges.
Molecule fk
- Δfk
+- fk
HCHO C 0.3973 0.2373 0.1023 0.0611 1.6744 1.8563 1.1088 0.1600 0.0412 0.7476
O 0.3010 0.4232 0.0775 0.1090 0.7113 1.4064 1.9774 -0.1222 -0.0315 -0.5710
CH3CHO C1 0.2998 0.1642 0.0715 0.0391 1.8267 1.0268 0.5624 0.1356 0.0324 0.4644
O 0.2708 0.3782 0.0646 0.0902 0.7165 0.9275 1.2953 -0.1074 -0.0256 -0.3678
CH3COCH3 C1 0.2108 0.1154 0.0494 0.0271 1.8262 0.5902 0.3231 0.0954 0.0223 0.2671
O 0.2359 0.3499 0.0553 0.0820 0.6742 0.6605 0.9797 -0.1140 -0.0267 -0.3192
C2H5COC2H5 C1 0.1346 0.0990 0.0313 0.0230 1.3598 0.3547 0.2609 0.0356 0.0083 0.0938
O 0.1449 0.2873 0.0337 0.0668 0.5045 0.3818 0.7570 -0.1424 -0.0331 -0.3752
CH2=CHCHO C1 0.1780 0.1357 0.0577 0.0440 1.3117 1.1219 0.8553 0.0423 0.0137 0.2666
C6 0.2062 0.1253 0.0668 0.0406 1.6457 1.2997 0.7898 0.0809 0.0262 0.5099
O 0.1797 0.3414 0.0582 0.1106 0.5264 1.1326 2.1518 -0.1620 -0.0524 -1.0191
CH3CH=CHCHO C6 0.1592 0.1114 0.0500 0.0350 1.4291 0.8532 0.5970 0.0478 0.0150 0.2562
C1 0.1741 0.1095 0.0547 0.0344 1.5900 0.9330 0.5868 0.0646 0.0203 0.3462
O 0.1739 0.2450 0.0546 0.0769 0.7098 0.9319 1.3130 -0.0710 -0.0223 -0.3810
TABLE 4: Calculated Global Reactivity Properties of the Selected Molecules using
B3LYP/6-311+g** and BLYP/DND method.
η μ ω S
η μ ω S
Molecules B3LYP/6-311+g** (eV) BLYP/DND (eV)
NH2OH 3.869 -3.553 1.632 0.129 3.411 -1.399 0.287 0.147
CH3ONH2 3.630 -3.738 1.925 0.138 3.549 -3.053 1.313 0.141
CH3NHOH 3.482 -3.392 1.652 0.144 3.229 -1.308 0.265 0.155
OHCH2CH2NH2 3.343 -3.507 1.840 0.150 3.348 -2.689 1.080 0.149
CH3SNH2 3.050 -3.331 1.819 0.164 2.447 -1.750 0.626 0.204
CH3NHSH 3.148 -3.629 2.092 0.159 2.466 -3.596 2.622 0.203
SHCH2CH2NH2 3.135 -3.417 1.862 0.159 2.521 -1.843 0.674 0.198
TABLE 5: Calculated Local Reactivity Properties of the Selected Molecules using B3LYP/6-
311+g** method for NPA derived charges.
Molecule fk
- sk
- Δfk
+- fk
NH2OH N 0.1870 0.4140 0.0274 0.0607 2.2139 0.0536 0.1187 -0.2270 -0.0333 -0.0651
O 0.2390 0.2300 0.0350 0.0337 0.9623 0.0685 0.0659 0.0090 0.0013 0.0026
CH3ONH2 C 0.0870 0.0680 0.1410 1.3130 0.0123 0.0096 0.7816 0.1142 0.0893 0.0190
N 0.1500 0.3510 0.0211 0.0495 2.3400 0.1969 0.4608 -0.2010 -0.0283 -0.2639
O 0.0720 0.1740 0.0101 0.0245 2.4167 0.0945 0.2284 -0.1020 -0.0144 -0.1339
CH3NHOH C 0.0470 0.0740 0.0073 0.0115 1.5745 0.0124 0.0196 -0.0270 -0.0042 -0.0071
N 0.1200 0.3390 0.0186 0.0525 2.8250 0.0318 0.0898 -0.2190 -0.0339 -0.0580
O 0.2100 0.1770 0.0325 0.0274 0.8429 0.0556 0.0469 0.0330 0.0051 0.0087
OHCH2CH2NH2 C1 0.0540 0.0330 0.0081 0.0049 0.6111 0.0583 0.0356 0.0210 0.0031 0.0227
C2 0.0400 0.0610 0.006 0.0091 1.5250 0.0432 0.0659 -0.0210 -0.0031 -0.0227
N 0.0630 0.3470 0.0094 0.0518 5.5079 0.0680 0.3746 -0.2840 -0.0424 -0.3066
O 0.1400 0.1010 0.0209 0.0151 0.7214 0.1511 0.1090 0.0390 0.0058 0.0421
CH3SNH2 C 0.0550 0.0640 0.0112 0.0131 1.1636 0.0344 0.0400 -0.0090 -0.0018 -0.0056
N 0.1490 0.0820 0.0305 0.0168 0.5503 0.0932 0.0513 0.0670 0.0137 0.0419
S 0.3580 0.5510 0.0732 0.1126 1.5391 0.2239 0.3447 -0.1930 -0.0394 -0.1207
CH3NHSH C 0.0530 0.0540 0.0107 0.0110 1.0189 0.1390 0.1416 -0.0010 -0.0002 -0.0026
N 0.1310 0.1740 0.0266 0.0353 1.3282 0.3434 0.4562 -0.0430 -0.0087 -0.1127
S 0.4530 0.4420 0.0919 0.0896 0.9757 1.1876 1.1588 0.0110 0.0022 0.0288
SHCH2CH2NH2 C1 0.0780 0.0410 0.0155 0.0081 0.5256 0.0525 0.0276 0.0370 0.0073 0.0249
C2 0.0290 0.0250 0.0058 0.0050 0.8621 0.0195 0.0168 0.0040 0.0008 0.0027
N 0.0380 0.1270 0.0075 0.0252 3.3421 0.0256 0.0856 -0.0890 -0.0177 -0.0600
S 0.3890 0.4710 0.0772 0.0934 1.2108 0.2621 0.3173 -0.0820 -0.0163 -0.0552
TABLE 6 Calculated Local Reactivity Properties of the Selected Molecules using BLYP/DND
method for HPA derived charges.
Molecule fk
- sk
- Δfk
+- fk
NH2OH N 0.1837 0.9327 0.0237 0.1205 5.0777 0.2997 1.5218 -0.7490 -0.0970 -1.2220
O -0.0770 0.5114 -0.0100 0.0661 -6.6170 -0.1261 0.8344 -0.5890 -0.0760 -0.9610
CH3ONH2 C 0.5410 0.0819 0.0746 0.0113 0.1513 1.0412 0.1576 0.4592 0.0633 0.8837
N -0.1510 0.2534 -0.0210 0.0349 -1.6740 -0.2913 0.4877 -0.4050 -0.0560 -0.7790
O -0.1790 0.9011 -0.0250 0.1242 -5.0267 -0.3450 1.7342 -1.0800 -0.1490 -2.0790
CH3NHOH C 0.4598 0.1677 0.0660 0.0241 0.3647 0.7598 0.2771 0.2921 0.0419 0.4827
N -0.0580 0.7950 -0.0080 0.1142 -13.725 -0.0957 1.3136 -0.8530 -0.1220 -1.4090
O -0.2690 0.4537 -0.0390 0.0651 -1.6855 -0.4448 0.7497 -0.7230 -0.1040 -1.1940
OHCH2CH2NH2 C1 0.1186 0.0254 0.0177 0.0038 0.2140 0.2181 0.0467 0.0932 0.0139 0.1715
C2 0.4003 0.1067 0.0599 0.0160 0.2666 0.7365 0.1964 0.2936 0.0439 0.5401
N -0.3040 0.9520 -0.0450 0.1424 -3.1337 -0.5589 1.7514 -1.2560 -0.1880 -2.3100
O -0.3340 0.5965 -0.0500 0.0892 -1.7842 -0.6151 1.0974 -0.9310 -0.1390 -1.7120
CH3SNH2 C 0.0667 0.3358 0.0100 0.0502 5.0377 0.1226 0.6178 -0.2690 -0.0400 -0.4950
N -0.297 0.4790 -0.044 0.0717 -1.6119 -0.5467 0.8813 -0.7760 -0.1160 -1.4280
S 0.3671 0.6485 0.0549 0.0970 1.7667 0.6753 1.1931 -0.2810 -0.0420 -0.5180
CH3NHSH C 0.1715 0.1732 0.0256 0.0259 1.0100 0.3154 0.3186 -0.0020 -0.0003 -0.0030
N -0.225 0.9064 -0.0340 0.1356 -4.0267 -0.4141 1.6676 -1.1320 -0.1690 -2.0820
S 0.3479 0.2249 0.0520 0.0336 0.6465 0.6400 0.4137 0.1230 0.01840 0.2262
SHCH2CH2NH2 C1 0.0117 0.2268 0.0017 0.0339 19.432 0.0215 0.4172 -0.2150 -0.0320 -0.3960
C2 0.1651 0.0876 0.0247 0.0131 0.5309 0.3037 0.1612 0.0774 0.0116 0.1425
N -0.292 0.7628 -0.0440 0.1141 -2.6164 -0.5364 1.4035 -1.0540 -0.1580 -1.9400
S 0.1064 0.5646 0.0159 0.0845 5.3089 0.1957 1.0388 -0.4580 -0.0690 -0.8430
TABLE 7: Atomic site with maximum value for multiphilic descriptor (∆ωk) for the
selected set of amines.
site with maximum value for ∆ωk molecule
NPA HPA
NH2OH N N
CH3ONH2 O N
CH3NHOH N N
OHCH2CH2NH2 N N
CH3SNH2 N S
CH3NHSH N N
SHCH2CH2NH2 N N
TABLE 8: Global reactivity descriptors calculated at B3LYP/6-31G* level of theory.
Species η
(eV)
(eV)
(eV)
Reactant 3.64 -2.89 1.15
Transition State 2.48 -2.79 1.57
Product 3.64 -2.89 1.15
TABLE 9: Global reactivity descriptors along the intrinsic reaction coordinate
calculated at B3LYP/6-31G* level of theory.
Points along
(Hartrees)
(eV)
(eV)
(eV)
(a.u.)
1 -234.5673091 2.65 -2.7825 1.46 64.94
2 -234.5661087 2.63 -2.7827 1.47 65.21
3 -234.5649450 2.61 -2.7828 1.49 65.47
4 -234.5638273 2.59 -2.7836 1.50 65.74
5 -234.5627655 2.57 -2.7836 1.51 65.98
6 -234.5617681 2.55 -2.7843 1.52 66.22
7 -234.5608445 2.54 -2.7843 1.53 66.42
8 -234.5600030 2.53 -2.7851 1.54 66.63
9 -234.5592516 2.51 -2.7852 1.54 66.80
10 -234.5585980 2.50 -2.7859 1.55 66.96
11 -234.5580104 2.50 -2.7857 1.56 67.07
12 -234.5575677 2.49 -2.7866 1.56 67.20
13 -234.5575677 2.49 -2.7866 1.56 67.20
14 -234.5580104 2.50 -2.7857 1.56 67.07
15 -234.5585980 2.50 -2.7859 1.55 66.96
16 -234.5592516 2.51 -2.7852 1.54 66.80
17 -234.5600030 2.53 -2.7851 1.54 66.63
18 -234.5608445 2.54 -2.7843 1.53 66.42
19 -234.5617681 2.55 -2.7843 1.52 66.22
20 -234.5627655 2.57 -2.7836 1.51 65.98
21 -234.5638273 2.59 -2.7836 1.50 65.74
22 -234.5649450 2.61 -2.7830 1.49 65.47
23 -234.5661087 2.63 -2.7827 1.47 65.21
24 -234.5673092 2.65 -2.7825 1.46 64.94
TABLE 10: Group Philicity ( +
gω ) Values for Nucleophilic and Electrophilic
Attacks Respectively for the Ionic Units of Different Isomers of LiAl4–, NaAl4–,
KAl4– and CuAl4–.
Isomers Ionic Unit
gωΔ
Al42– 0.0070 0.0095 0.0025 LiAl4–
(C∞v) Li+ 0.0063 0.0037 -0.0025
Al42– 1.3E-05 0.0055 0.0055 LiAl4–
(C2v) Li+ 0.0068 0.0013 -0.0055
Al42– -0.0372 0.2965 0.3338 LiAl4–
(C4v) Li+ 0.4055 0.0718 -0.3338
Al42– 0.0070 0.0102 0.0032 NaAl4–
(C∞v) Na+ 0.0074 0.0042 -0.0032
Al42– -0.0001 0.0078 0.0079 NaAl4–
(C2v) Na+ 0.0096 0.0017 -0.0079
Al42– -0.0073 0.1024 0.1097 NaAl4–
(C4v) Na+ 0.1301 0.0204 -0.1097
Al42– 0.0044 0.0095 0.0051 KAl4–
(C∞v) K+ 0.0106 0.0054 -0.0051
Al42– 0.0023 0.0101 0.0078 KAl4–
(C2v) K+ 0.0118 0.0039 -0.0078
Al42– 0.0008 0.0066 0.0057 KAl4–
(C4v) K+ 0.0078 0.0021 -0.0057
Al42– 0.0031 0.0036 0.0006 CuAl4–
(C∞v) Cu+ 0.0014 0.0009 -0.0006
Al42– 0.0036 0.0036 0.0048 CuAl4–
(C2v) Cu+ 0.0008 0.0008 -0.0048
Al42– 0.0178 0.0332 0.0154 CuAl4–
(C4v) Cu+ 0.0131 -0.0023 -0.0154
TABLE 11: Profiles of the forward activation energy ( #fEΔ ), reverse activation
energy ( #rEΔ ) and reaction energy (
0EΔ ) of a thermoneutral reaction (Fa– + CH3-Fb
→ Fa--CH3 + Fb–; an endothermic reaction (HNO → HON) and an exothermic
reaction (H2OO → HOOH).
Reaction #fEΔ
ξ1 ξ2 W1 W2 W3 W4
Thermo-neutral
B3LYP/6-311++G**
9.54 9.54 0.0 -1.33 1.33 5.42 4.12 -4.12 - 5.42
Endothermic
B3LYP/6-311+G**
75.39 34.84 40.55 -0.80 0.60 43.97 31.42 -13,20 - 21.64
Exothermic
B3LYP/6-311+G**
7.39 52.85 -45.46 -0.65 0.87 3.93 3.46 - 19.99 - 32.86
Figure 1. Optimized structures with atom numbering for the selected carbonyl
compounds.
Figure 2. Optimized structures with atom numbering for the selected amine systems.
Reactant Transition State Product
Figure 3:Optimized geometrical structures calculated using B3LYP/6-31G* level of
theory.
-234.568
-234.566
-234.564
-234.562
-234.560
-234.558
-234.556
Energy (Hartree)
Electrophilicity Index (eV)
Intrinsic Reaction Coordinate
lectrophilicity Index (eV
2.66 Chemical Hardness (eV)
Polarizability (au)
Intrinsic Reaction Coordinate
olarizability (au)
Figure 4 (a-b):Variation of global reactivity descriptors along intrinsic reaction
coordinate.
-0.008
-0.006
-0.004
-0.002
0.000
0.002
0.004
0.006
Intrinsic Reaction Coordinate
C1,C3 sites
-0.008
-0.006
-0.004
-0.002
0.000
0.002
0.004
0.006
Intrinsic Reaction Coordinate
C6,C11 sites
Figure 5 (a-b): Variation of multiphilic descriptor along intrinsic reaction coordinate for
the selected atomic sites.
MAl4– [C∞v]
MAl4– [C2v]
MAl4– [C4v]
M=Li, Na, K, Cu
Figure 6. Optimized structures of various isomers of MAl4– (M ≡ Li, Na, K, Cu).
-3 -2 -1 0 1 2 3
-239.704
-239.702
-239.700
-239.698
-239.696
-239.694
-239.692
-239.690
-239.688
-239.686
-3 -2 -1 0 1 2 3
0.002
0.004
0.006
0.008
0.010
0.012
0.014
0.016
0.018
Energy Δω (Fa)
Δω (F
(a)
-2.5 -2.0 -1.5 -1.0 -0.5 0.0 0.5 1.0 1.5 2.0
-130.52
-130.50
-130.48
-130.46
-130.44
-130.42
-130.40
-130.38
-2.5 -2.0 -1.5 -1.0 -0.5 0.0 0.5 1.0 1.5 2.0
-0.07
-0.06
-0.05
-0.04
-0.03
-0.02
-0.01
Energy
-2 -1 0 1 2 3
-151.60
-151.58
-151.56
-151.54
-151.52
-151.50
-2 -1 0 1 2 3
-0.02
Energy
(O2) Δω
(O1)
(c)
Figure 7 (a-c): Profiles of net nucleophilicity (∆ωk) of along the path of the gas phase (a)
thermoneutral SN2 substitution: Fa- + CH3-Fb → Fa-CH3 + Fb-, (b) endothermic reaction:
HNO → HON and (c) exothermic reaction: H2OO → HOOH. Also shown is the profile
of energy.
Figure 8 Reaction force profiles along the reaction coordinate for (a) thermoneutral
reaction: Fa– + CH3-Fb → Fa--CH3 + Fb–; (b) endothermic reaction: HNO → HON; (c) the
exothermic reaction: H2OO → HOOH. The vertical dashed lines define the reaction
regions as follows: reactant (left), transition state (middle) and product (right).
-4 -2 0 2 4
-2 -1 0 1 2 3
-2 -1 0 1 2
max(a)
|
0704.0335 | Approximation of the distribution of a stationary Markov process with
application to option pricing | Approximation of the distribution of a stationary Markov process with application to option pricing
Bernoulli 15(1), 2009, 146–177
DOI: 10.3150/08-BEJ142
Approximation of the distribution of a
stationary Markov process with application to
option pricing
GILLES PAGÈS1 and FABIEN PANLOUP2
Laboratoire de Probabilités et Modèles Aléatoires, UMR 7599, Université Paris 6, Case 188, 4
pl. Jussieu, F-75252 Paris Cedex 5. E-mail: [email protected]
Laboratoire de Statistiques et Probabilités, Université Paul Sabatier & INSA Toulouse, 135,
Avenue de Rangueil, 31077 Toulouse Cedex 4. E-mail: [email protected]
We build a sequence of empirical measures on the space D(R+,R
d) of Rd-valued cadlag functions
on R+ in order to approximate the law of a stationary R
d-valued Markov and Feller process (Xt).
We obtain some general results on the convergence of this sequence. We then apply them to
Brownian diffusions and solutions to Lévy-driven SDE’s under some Lyapunov-type stability
assumptions. As a numerical application of this work, we show that this procedure provides an
efficient means of option pricing in stochastic volatility models.
Keywords: Euler scheme; Lévy process; numerical approximation; option pricing; stationary
process; stochastic volatility model; tempered stable process
1. Introduction
1.1. Objectives and motivations
In this paper, we deal with an Rd-valued Feller Markov process (Xt) with semigroup
(Pt)t≥0 and assume that (Xt) admits an invariant distribution ν0. The aim of this work is
to propose a way to approximate the whole stationary distribution Pν0 of (Xt). More pre-
cisely, we want to construct a sequence of weighted occupation measures (ν(n)(ω,dα))n≥1
on the Skorokhod space D(R+,R
d) such that ν(n)(ω,F )
n→+∞−→
F (α)Pν0(dα) a.s. for a
class of functionals F :D(R+,R
d) which includes bounded continuous functionals for the
Skorokhod topology.
One of our motivations is to develop a new numerical method for option pricing in sta-
tionary stochastic volatility models which are slight modifications of the classical stochas-
tic volatility models, where we suppose that the volatility evolves under its stationary
regime.
This is an electronic reprint of the original article published by the ISI/BS in Bernoulli,
2009, Vol. 15, No. 1, 146–177. This reprint differs from the original in pagination and
typographic detail.
1350-7265 c© 2009 ISI/BS
http://arxiv.org/abs/0704.0335v3
http://isi.cbs.nl/bernoulli/
http://dx.doi.org/10.3150/08-BEJ142
mailto:[email protected]
mailto:[email protected]
http://isi.cbs.nl/BS/bshome.htm
http://isi.cbs.nl/bernoulli/
http://dx.doi.org/10.3150/08-BEJ142
Approximation of the distribution of a stationary Markov process 147
1.2. Background and construction of the procedure
This work follows on from a series of recent papers due to Lamberton and Pagès ([12, 13]),
Lemaire ([14, 15]) and Panloup ([18, 19, 20]), where the problem of the approximation
of the invariant distribution is investigated for Brownian diffusions and for Lévy-driven
SDE’s.1 In these papers, the algorithm is based on an adapted Euler scheme with de-
creasing step (γk)k≥1. To be precise, let (Γn) be the sequence of discretization times:
Γ0 = 0, Γn =
k=1 γk for every n≥ 1, and assume that Γn → +∞ when n→+∞. Let
(X̄Γn)n≥0 be the Euler scheme obtained by “freezing” the coefficients between the Γn’s
and let (ηn)n≥1 be a sequence of positive weights such that Hn :=
k=1 ηk →+∞ when
k→+∞. Then, under some Lyapunov-type stability assumptions adapted to the stochas-
tic processes of interest, one shows that for a large class of steps and weights (ηn, γn)n≥1,
ν̄n(ω, f) :=
ηkf(X̄Γk−1)
n→+∞−→
f(x)ν0(dx) a.s., (1)
(at least)2 for every bounded continuous function f .
Since the problem of the approximation of the invariant distribution has been deeply
studied for a wide class of Markov processes (Brownian diffusions and Lévy-driven SDE’s)
and since the proof of (1) can be adapted to other classes of Markov processes under some
specific Lyapunov assumptions, we choose in this paper to consider a general Markov pro-
cess and to assume the existence of a time discretization scheme (X̄Γk)k≥0 such that (1)
holds for the class of bounded continuous functions. The aim of this paper is then to inves-
tigate the convergence properties of a functional version of the sequence (ν̄n(ω,dα))n≥1.
Let (Xt) be a Markov and Feller process and let (X̄t)t≥0 be a stepwise constant time
discretization scheme of (Xt) with non-increasing step sequence (γn)n≥1 satisfying
γn = 0, Γn :=
n→+∞−→ +∞. (2)
Letting Γ0 := 0 and X̄0 = x0 ∈Rd, we assume that
X̄t = X̄Γn ∀t ∈ [Γn,Γn+1[ (3)
and that (X̄Γn)n≥0 can be simulated recursively.
We denote by (Ft)t≥0 and (F̄t)t≥0 the usual augmentations of the natural filtrations
(σ(Xs,0≤ s≤ t))t≥0 and (σ(X̄s,0≤ s≤ t))t≥0, respectively.
1Note that computing the invariant distribution is equivalent to computing the marginal laws of the
stationary process (Xt) since ν0Pt = ν0 for every t≥ 0.
2The class of functions for which (1) holds depends on the stability of the dynamical system. In
particular, in the Brownian diffusion case, the convergence may hold for continuous functions with
subexponential growth, whereas the class of functions strongly depends on the moments of the Lévy
process when the stochastic process is a Lévy-driven SDE.
148 G. Pagès and F. Panloup
For k ≥ 0, we denote by (X̄(k)t )t≥0 the shifted process defined by
t := X̄Γk+t.
In particular, X̄
t = X̄t. We define a sequence of random probabilities (ν
(n)(ω,dα))n≥1
on D(R+,R
d) by
ν(n)(ω,dα) =
ηk1{X̄(k−1)(ω)∈dα},
where (ηk)k≥1 is a sequence of weights. For t ≥ 0, (ν(n)t (ω,dx))n≥1 will denote the se-
quence of “marginal” empirical measures on Rd defined by
t (ω,dx) =
ηk1{X̄(k−1)
(ω)∈dx}
1.3. Simulation of (ν(n)(ω,F ))
For every functional F :D(R+,R
d)→R, the following recurrence relation holds for every
n≥ 1:
ν(n+1)(ω,F ) = ν(n)(ω,F ) +
(F (X(n)(ω))− ν(n)(ω,F )). (4)
Then, if T is a positive number and F :D(R+,R
d) → R is a functional depending only
on the trajectory between 0 and T , (ν(n)(ω,F ))n≥1 can be simulated by the following
procedure.
Step 0. (i) Simulate (X̄
t )t≥0 on [0, T ], that is, simulate (X̄Γk)k≥0 for k =
0, . . . ,N(0, T ), where
N(n,T ) := inf{k ≥ n,Γk+1 − Γn > T }
= max{k ≥ 0,Γk − Γn ≤ T }, n≥ 0, T > 0.
Note that n 7→N(n, t) is an increasing sequence since (γn) is non-increasing, and that
ΓN(n,T ) − Γn ≤ T < ΓN(n,T )+1 − Γn.
(ii) Compute F ((X̄
t )t≥0) and ν
(1)(ω,F ). Store the values of (X̄Γk) for k =
1, . . . ,N(0, T ).
Step n (n≥ 1). (i) Since the values (X̄Γk)k≥0 are stored for k = n, . . . ,N(n− 1, T ),
simulate (X̄Γk)k≥0 for k =N(n−1, T )+1, . . . ,N(n,T ) in order to obtain a path of (X̄
on [0, T ].
(ii) Compute F ((X̄
t )t≥0) and use (4) to compute ν
(n+1)(ω,F ). Store the values of
(X̄Γk) for k = n+ 1, . . . ,N(n,T ).
Approximation of the distribution of a stationary Markov process 149
Remark 1. As shown in the description of the procedure, one generally has to store
the vector [X̄Γn , . . . , X̄ΓN(n,T) ] at time n. Since (γn) is a sequence with infinite sum that
decreases to 0, it follows that the size of this vector increases “slowly” to +∞. For
instance, if γn = Cn
−ρ with ρ ∈ (0,1), its size is of order nρ. However, it is important
to remark that even though the number of values to be stored tends to +∞, that is
not always the case for the number of operations at each step. Indeed, since X̄(n+1)
is obtained by shifting X̄(n), it is usually possible to use, at step n+ 1, the preceding
computations and to simulate the sequence (F (X̄(n)))n≥0 in a “quasi-recursive” way.
For instance, such remark holds for Asian options because the associated pay-off can be
expressed as a function of an additive functional (see Section 5 for simulations).
Before outlining the sequel of the paper, we list some notation linked to the spaces
D(R+,R
d) and D([0, T ],Rd) of cadlag Rd-valued functions on R+ and [0, T ], respectively,
endowed with the Skorokhod topology. First, we denote by d1 the Skorokhod distance
on D([0,1],Rd) defined for every α, β ∈D([0,1],Rd) by
d1(α,β) = inf
t∈[0,1]
|α(t)− β(λ(t))|, sup
0≤s<t≤1
λ(t)− λ(s)
where Λ1 denotes the set of increasing homeomorphisms of [0,1]. Second, for T > 0,
φT :D(R+,R
d) 7→D([0,1],Rd) is the function defined by (φT (α))(s) = α(sT ) for every s ∈
[0,1]. We then denote by d the distance on D(R+,R
d) defined for every α,β ∈D(R+,Rd)
d(α,β) =
e−t(1∧ d1(φt(α), φt(β))) dt. (6)
We recall that (D(R+,R
d), d) is a Polish space and that the induced topology is the usual
Skorokhod topology on D(R+,R
d) (see, e.g., Pagès [16]). For every T > 0, we set
σ(πu,0≤ u≤ s),
where πs :D(R+,R
d)→Rd is defined by πs(α) = α(s). For a functional F :D(R+,Rd)→
R, FT denotes the functional defined for every α ∈D(R+,Rd) by
FT (α) = F (α
T ) with αT (t) = α(t ∧ T ) ∀t≥ 0. (7)
Finally, we will say that a functional F :D(R+,R
d)→R is Sk-continuous if F is contin-
uous for the Skorokhod topology on D(R+,R
d) and the notation “
=⇒” will denote the
weak convergence on D(R+,R
In Section 2, we state our main results for a general Rd-valued Feller Markov process.
Then, in Section 3, we apply them to Brownian diffusions and Lévy-driven SDE’s. Section
4 is devoted to the proofs of the main general results. Finally, in Section 5, we complete
this paper with an application to option pricing in stationary stochastic volatility models.
150 G. Pagès and F. Panloup
2. General results
In this section, we state the results on convergence of the sequence (ν(n)(ω,dα))n≥1 when
(Xt) is a general Feller Markov process.
2.1. Weak convergence to the stationary regime
As explained in the Introduction, since the a.s. convergence of (ν
0 (ω,dx))n≥1 to the
invariant distribution ν0 has already been deeply studied for a large class of Markov
processes (Brownian diffusions and Lévy driven SDE’s), our approach will be to derive
the convergence of (ν(n)(ω,dα))n≥1 toward Pν0 from that of (ν
0 (ω,dx))n≥1 to the
invariant distribution ν0. More precisely, we will assume in Theorem 1 that
(C0,1): (Xt) admits a unique invariant distribution ν0 and
0 (ω,dx)
=⇒ ν0(dx) a.s.,
whereas in Theorem 2, we will only assume that
(C0,2): (ν
0 (ω,dx))n≥1 is a.s. tight on R
We also introduce three other assumptions, (C1), (C2) and (C3,ε), regarding the conti-
nuity in probability of the flow x 7→ (Xxt ), the asymptotic convergence of the shifted time
discretization scheme to the true process (Xt) and the steps and weights, respectively.
(C1): For every x0 ∈Rd, ǫ > 0 and T > 0,
limsup
0≤t≤T
|Xxt −Xx0t | ≥ ǫ
= 0. (8)
(C2): (X̄t) is a non-homogeneous Markov process and for every n≥ 0, it is possible to
construct a family of stochastic processes (Y
(n,x)
t )x∈Rd such that
(i) L(Y (n,x)) D(R+,R
= L(X̄(n)|X̄(n)0 = x);
(ii) for every compact set K of Rd, for every T ≥ 0,
0≤t≤T
|Y (n,x)t −Xxt |
n→+∞−→ 0 in probability. (9)
(C3,ε): For every n≥ 1, ηn ≤CγnHεn.
Remark 2. Assumption (C2) implies, in particular, that asymptotically and uniformly
on compact sets of Rd, the law of the approximate process (X̄(n)), given its initial value,
is close to that of the true process.
If there exists a unique invariant distribution ν0, the second part of (C2) can be relaxed
to the following, less stringent, assertion: for all ǫ > 0, there exists a compact set Aǫ ⊂Rd
such that ν0(A
ǫ)≤ ǫ and such that
0≤t≤T
|Y (n,x)t −Xxt |
n→+∞−→ 0 in probability. (10)
Approximation of the distribution of a stationary Markov process 151
This weaker assumption can some times be needed in stochastic volatility models like
the Heston model (see Section 5 for details).
The preceding assumptions are all that we require for the convergence of (ν(n)(ω,dα))n≥1
to Pν0 along the bounded Sk-continuous functionals, that is, for the a.s. weak conver-
gence on D(R+,R
d). However, the integration of non-bounded continuous functionals
F :D([0, T ],Rd)→ R will need some additional assumptions, depending on the stability
of the time discretization scheme and on the steps and weights sequences. We will sup-
pose that F is dominated (in a sense to be specified later) by a function V : Rd → R+
that satisfies the following assumptions for some s≥ 2 and ε < 1.
H(s, ε): For every T > 0,
(i) sup
0≤t≤T
Vs(Y (n,x)t )
≤CTVs(x),
(ii) sup
0 (V)<+∞,
(iii)
E[V2(X̄Γk−1)]<+∞,
∆N(k,T )
E[Vs(1−ε)(X̄Γk−1)]<+∞,
where T 7→CT is locally bounded on R+ and ∆N(k,T ) =N(k,T )−N(k− 1, T ).
For every ε < 1, we then set
K(ε) = {V ∈ C(Rd,R+),H(s, ε) holds for some s≥ 2}.
Remark 3. Apart from assumption (i), which is a classical condition on the finite time
horizon control, the assumptions in H(s, ε) strongly rely on the stability of the time
discretization scheme (and then, to that of the true process). More precisely, we will see
when we apply our general results to SDE’s that these properties are some consequences
of the Lyapunov assumptions needed for the tightness of (ν
0 (ω,dx))n≥1.
We can now state our first main result.
Theorem 1. Assume (C0,1), (C1), (C2) and (C3,ε) with ε ∈ (−∞,1). Then, a.s., for
every bounded Sk-continuous functional F :D(R+,R
d)→R,
ν(n)(ω,F )
n→+∞−→
F (α)Pν0 (dα), (11)
where Pν0 denotes the stationary distribution of (Xt) (with initial law ν0).
Furthermore, for every T > 0, for every non-bounded Sk-continuous functional
F :D(R+,R
d)→ R, (11) holds a.s. for FT (defined by (7)) if there exists V ∈ K(ε) and
152 G. Pagès and F. Panloup
ρ ∈ [0,1) such that
|FT (α)| ≤C sup
0≤t≤T
Vρ(αt) ∀α ∈D(R+,Rd). (12)
In the second result, the uniqueness of the invariant distribution is not required and
the sequence (ν
0 (ω,dx))n≥1 is only supposed to be tight.
Theorem 2. Assume (C0,2), (C1), (C2) and (C3,ε) with ε ∈ (−∞,1). Assume that
0 (ω,dx))n≥1 is a.s. tight on R
d. We then have the following.
(i) The sequence (ν(n)(ω,dα))n≥1 is a.s. tight on D(R+,R
d) and a.s., for ev-
ery convergent subsequence (nk(ω))n≥1, for every bounded Sk-continuous functional
F :D(R+,R
d)→R,
ν(nk(ω))(ω,F )
n→+∞−→
F (α)Pν∞(dα), (13)
where Pν∞ is the law of (Xt) with initial law ν∞ being a weak limits for (ν
0 (ω,dx))n≥1.
Furthermore, for every T > 0, for every non-bounded Sk-continuous functional
F :D(R+,R
d)→R, (13) holds a.s. for FT if (12) is satisfied with V ∈K(ε) and ρ ∈ [0,1).
(ii) If, moreover,
l≥k+1
|∆ηℓ|
n→+∞−→ 0, (14)
then ν∞ is necessarily an invariant distribution for the Markov process (Xt).
Remark 4. Condition (14) holds for a large class of steps and weights. For instance,
if ηn = C1n
−ρ1 and γn = C2n
−ρ2 with ρ1 ∈ [0,1] and ρ2 ∈ (0,1], then (14) is satisfied if
ρ1 = 0 or if ρ1 ∈ (max(0,2ρ2 − 1),1).
2.2. Extension to the non-stationary case
Even though the main interest of this algorithm is the weak approximation of the pro-
cess when stationary, we observe that when ν0 is known, the algorithm can be used to
approximate Pµ0 if µ0 is a probability on R
d that is absolutely continuous with respect
to ν0.
Indeed, assume that µ0(dx) = φ(x)ν0(dx), where φ :R
d → R is a continuous non-
negative function. For a functional F :D(R+,R
d)→ R, denote by Fφ the functional de-
fined on D(R+,R
d) by Fφ(α) = F (α)φ(α(0)).
Then, if ν(n)(ω,dα)
(Sk)⇒ Pν0(dα) a.s., we also have the following convergence: a.s., for
every bounded Sk-continuous functional F :D(R+,R
d)→R,
ν(n)(ω,Fφ)
n→+∞−→
Fφ(α)Pν0 (dα) =
F (α)Pµ0 (dα).
Approximation of the distribution of a stationary Markov process 153
3. Application to Brownian diffusions and
Lévy-driven SDE’s
Let (Xt)t≥0 be a cadlag stochastic process solution to the SDE
dXt = b(Xt−) dt+ σ(Xt−) dWt + κ(Xt−) dZt, (15)
where b :Rd → Rd, σ :Rd 7→Md,ℓ (set of d× ℓ real matrices) and κ :Rd 7→Md,ℓ are con-
tinuous functions with sublinear growth, (Wt)t≥0 is an ℓ-dimensional Brownian motion
and (Zt)t≥0 is an integrable purely discontinuous R
ℓ-valued Lévy process independent of
(Wt)t≥0 with Lévy measure π and characteristic function given for every t≥ 0 by
E[ei〈u,Zt〉] = exp
ei〈u,y〉 − 1− i〈u, y〉π(dy)
Let (γn)n≥1 be a non-increasing step sequence satisfying (2). Let (Un)n≥1 be a sequence
of i.i.d. random variables such that U1
=N (0, Iℓ) and let ξ := (ξn)n≥1 be a sequence of
independent Rℓ-valued random variables, independent of (Un)n≥1. We then denote by
(X̄t)t≥0 the stepwise constant Euler scheme of (Xt) for which (X̄Γn)n≥0 is recursively
defined by X̄0 = x ∈Rd and
X̄Γn+1 = X̄Γn + γn+1b(X̄Γn) +
γn+1σ(X̄Γn)Un+1 + κ(X̄Γn)ξn+1. (16)
We recall that the increments of (Zt) cannot be simulated in general. That is why we
generally need to construct the sequence (ξn) with some approximations of the true
increments. We will come back to this construction in Section 3.2.
As in the general case, we denote by (X̄(k))k≥0 and (ν
(n)(ω,dα))n≥1 the sequences of
associated shifted Euler schemes and empirical measures, respectively.
Let us now introduce some Lyapunov assumptions for the SDE. Let EQ(Rd) denote
the set of essentially quadratic C2-functions V :Rd → R∗+ such that limV (x) = +∞ as
|x| →+∞, |∇V | ≤C
V and D2V is bounded. Let a ∈ (0,1] denote the mean reversion
intensity. The Lyapunov (or mean reversion) assumption is the following.
(Sa): There exists a function V ∈ EQ(Rd) such that:
(i) |b|2 ≤CV a, Tr(σσ∗(x)) + ‖κ(x)‖2 |x|→+∞= o(V a(x));
(ii) there exist β ∈R and ρ > 0 such that 〈∇V, b〉 ≤ β − ρV a.
From now on, we separate the Brownian diffusions and Lévy-driven SDE cases.
3.1. Application to Brownian diffusions
In this part, we assume that κ= 0. We recall a result by Lamberton and Pagès [13].
Proposition 1. Let a ∈ (0,1] such that (Sa) holds. Assume that the sequence (ηn/γn)n≥1
is non-increasing.
154 G. Pagès and F. Panloup
(a) Let (θn)n≥1 be a sequence of positive numbers such that
n≥1 θnγn < +∞ and
that there exists n0 ∈N such that (θn)n≥n0 is non-increasing. Then, for every positive r,
θnγnE[V
r(X̄Γn−1)]<+∞.
(b) For every r > 0,
0 (ω,V
r)<+∞ a.s. (17)
Hence, the sequence (ν
0 (ω,dx))n≥1 is a.s. tight.
(c) Moreover, every weak limit of this sequence is an invariant probability for the SDE
(15). In particular, if (Xt)t≥0 admits a unique invariant probability ν0, then for every
continuous function f such that f ≤CV r with r > 0, limn→∞ ν(n)0 (ω, f) = ν0(f) a.s.
Remark 5. For instance, if V (x) = 1 + |x|2, then the preceding convergence holds for
every continuous function with polynomial growth. According to Theorem 3.2 in Lemaire
[14], it is possible to extend these results to continuous functions with exponential growth,
but it then strongly depends on σ. Further the conditions on steps and weights can be
less restrictive and may contain the case ηn = 1, for instance (see Remark 4 of Lamberton
and Pagès [13] and Lemaire [14]).
We then derive the following result from the preceding proposition and from Theorems
1 and 2.
Theorem 3. Assume that b and σ are locally Lipschitz functions and that κ = 0. Let
a ∈ (0,1] such that (Sa) holds and assume that (ηn/γn) is non-increasing.
(a) The sequence (ν(n)(ω,dα))n≥1 is a.s. tight on C(R+,Rd)3 and every weak limit
of (ν(n)(ω,dα))n≥1 is the distribution of a stationary process solution to (15). In par-
ticular, when uniqueness holds for the invariant distribution ν0, a.s., for every bounded
continuous functional F :C(R+,Rd)→R,
ν(n)(ω,F )
n→+∞−→
F (x)Pν0 (dx). (18)
(b) Furthermore, if there exists s ∈ (2,+∞) and n0 ∈N such that
∆N(k,T )
is non-increasing and
∆N(k,T )
<+∞, (19)
3C(R+,R
d) denotes the space of continuous functions on R+ with values in R
d endowed with the
topology of uniform convergence on compact sets.
Approximation of the distribution of a stationary Markov process 155
then, for every T > 0, for every non-bounded continuous functional F :C(R+,Rd)→ R,
(18) holds for FT if the following condition is satisfied:
∃r > 0 such that |FT (α)| ≤C sup
0≤t≤T
V r(αt) ∀α ∈ C(R+,Rd).
Remark 6. If ηn =C1n
−ρ1 and γn =C2n
−ρ2 with 0< ρ2 ≤ ρ1 ≤ 1, then for s ∈ (1,+∞),
(19) is fulfilled if and only if s > 1/(1− ρ1). It follows that there exists s ∈ (2,+∞) such
that (19) holds as soon as ρ1 < 1.
Proof of Theorem 3. We want to apply Theorem 2. First, by Proposition 1, assumption
(C0,2) is fulfilled and every weak limit of (ν
0 (ω,dx)) is an invariant distribution. Second,
it is well known that (C1) and (C2) are fulfilled when b and σ are locally Lispchitz
sublinear functions. Then, since (C3,ε) holds with ε = 0, (18) holds for every bounded
continuous functional F . Finally, one checks that H(s,0) holds with V := V r (r > 0).
It is classical that assumption (a) is true when b and σ are sublinear. Assumption (b)
follows from Proposition 1(b). Let θn,1 = ηn/(γnH
n) and θn,2 =∆N(n,T )/(γnH
n). Using
(19) and the fact that (ηn/γn) is non-increasing yields that (θn,1) and (θn,2) satisfy the
conditions of Proposition 1 (see (35) for details). Then, (iii) and (iv) of H(s,0) are
consequences of Proposition 1(a). This completes the proof. �
3.2. Application to Lévy-driven SDE’s
When we want to extend the results obtained for Brownian SDE’s to Lévy-driven SDE’s,
one of the main difficulties comes from the moments of the jump component (see Panloup
[18] for details). For simplification, we assume here that (Zt) has a moment of order
2p≥ 2, that is, that its Lévy measure π satisfies the following assumption with p≥ 1:
(H1p) :
|y|>1
π(dy)|y|2p <+∞.
We also introduce an assumption about the behavior of the moments of the Lévy measure
at 0:
(H2q) :
|y|≤1
π(dy)|y|2q <+∞, q ∈ [0,1].
This assumption ensures that (Zt) has finite 2q-variations. Since
|y|≤1
|y|2π(dy) is finite,
this is always satisfied for q = 1.
Let us now specify the law of (ξn) introduced in (16). When the increments of (Zt) can
be exactly simulated, we denote by (E) the Euler scheme and by (ξn,E) the associated
sequence
= Zγn ∀n≥ 1.
156 G. Pagès and F. Panloup
When the increments of (Zt) cannot be simulated, we introduce some approximated Euler
schemes (P) and (W) built with some sequences (ξn,P ) and (ξn,W ) of approximations of
the true increment (see Panloup [19] for more detailed presentations of these schemes).
In scheme (P),
=Zγn,n,
where (Z·,n)n≥1 a sequence of compensated compound Poisson processes obtained by
truncating the small jumps of (Zt)t≥0:
Zt,n :=
0<s≤t
∆Zs1{|∆Zs|>un} − t
|y|>un
yπ(dy) ∀t≥ 0, (20)
where (un)n≥1 is a sequence of positive numbers such that un → 0. We recall that
n→+∞−→ Z locally uniformly in L2 (see, e.g., Protter [21]).
As shown in Panloup [19], the error induced by this approximation is very large when
the local behavior of the small jumps component is irregular. However, it is possible to
refine this approximation by a Wienerization of the small jumps, that is, by replacing
the small jumps by a linear transform of a Brownian motion instead of discarding them
(see Asmussen and Rosinski [2]). The corresponding scheme is denoted by (W) with ξn,W
satisfying
= ξn,P +
γnQnΛn ∀n≥ 1,
where (Λn)n≥1 is a sequence of i.i.d. random variables, independent of (ξn,P )n≥1 and
(Un)n≥1, such that Λ1
=N (0, Iℓ) and (Qn) is a sequence of ℓ× ℓ matrices such that
n)i,j =
|y|≤uk
yiyjπ(dy).
We recall the following result obtained in Panloup [18] in our slightly simplified frame-
work.
Proposition 2. Let a ∈ (0,1], p≥ 1 and q ∈ [0,1] such that (H1p), (H2q) and (Sa) hold.
Assume that the sequence (ηn/γn)n≥1 is non-increasing. Then, the following assertions
hold for schemes (E), (P) and (W).
(a) Let (θn) satisfy the conditions of Proposition 1. Then,
n≥1 θnγnE[V
p+a−1(X̄Γn−1)]<
(b) We have
0 (ω,V
p/2+a−1)<+∞ a.s. (21)
Hence, the sequence (ν
0 (ω,dx))n≥1 is a.s. tight as soon as p/2+ a− 1> 0.
Approximation of the distribution of a stationary Markov process 157
(c) Moreover, if Tr(σσ∗)+ ‖κ‖2q ≤CV p/2+a−1, then every weak limit of this sequence
is an invariant probability for the SDE (15). In particular, if (Xt)t≥0 admits a unique
invariant probability ν0, for every continuous function f such that f = o(V
p/2+a−1),
limn→∞ ν
0 (ω, f) = ν0(f) a.s.
Remark 7. For schemes (E) and (P), the above proposition is a direct consequence of
Theorem 2 and Proposition 2 of Panloup [18]. As concerns scheme (W), a straightforward
adaptation of the proof yields the result.
Our main functional result for Lévy-driven SDE’s is then the following.
Theorem 4. Let a ∈ (0,1] and p≥ 1 such that p/2+ a− 1> 0 and let q ∈ [0,1]. Assume
(H1p), (H
q) and (Sa). Assume that b, σ and κ are locally Lipschitz functions. If, more-
over, (ηn/γn)n≥1 is non-increasing, then the following result holds for schemes (E), (P)
and (W).
(a) The sequence (ν(n)(ω,dα))n≥1 is a.s. tight on D(R+,R
d). Moreover, if
Tr(σσ∗) + ‖κ‖2q ≤CV p/2+a−1 or 1
l≥k+1
|∆ηℓ|
n→+∞−→ 0, (22)
then every weak limit of (ν(n)(ω,dα))n≥1 is the distribution of a stationary process solu-
tion to (15).
(b) Assume that the invariant distribution is unique. Let ε≤ 0 such that (C3,ε) holds.
Then, a.s., for every T > 0, for every Sk-continuous functional F :D(R+,R
d)→R, (18)
holds for FT if there exist ρ ∈ [0,1) and s≥ 2, such that
|FT (α)| ≤C sup
0≤t≤T
V (ρ(p+a−1))/s(αt) ∀α ∈D(R+,Rd)
and if
∆N(k,T )
s(1−ε)
is non-increasing and
∆N(k,T )
s(1−ε)
<+∞. (23)
Remark 8. In (22), both assumptions imply the invariance of every weak limit of
0 (ω,dx)). These two assumptions are very different. The first is needed in Proposition
2 for using the Echeverria–Weiss invariance criteria (see Ethier and Kurtz [7], page 238,
Lamberton and Pagès [12] and Lemaire [14]), whereas the second appears in Theorem
2, where our functional approach shows that under some mild additional conditions on
steps and weights, every weak limit is always invariant.
For (23), we refer to Remark 6 for simple sufficient conditions when (γn) and (ηn) are
some polynomial steps and weights.
158 G. Pagès and F. Panloup
4. Proofs of Theorems 1 and 2
We begin the proof with some technical lemmas. In Lemma 1, we show that the a.s
weak convergence of the random measures (ν(n)(ω,dα))n≥1 can be characterized by the
convergence (11) along the set of bounded Lipschitz functionals F for the distance d.
Then, in Lemma 2, we show with some martingale arguments that if the functional
F depends only on the restriction of the trajectory to [0, T ], then the convergence of
(ν(n)(ω,F ))n≥1 is equivalent to that of a more regular sequence. This step is fundamental
for the sequel of the proof.
Finally, Lemma 4 is needed for the proof of Theorem 2. We show that under some mild
conditions on the step and weight sequences, any Markovian weak limit of the sequence
(ν(n)(ω,dα))n≥1 is stationary.
4.1. Preliminary lemmas
Lemma 1. Let (E,d) be a Polish space and let P(E) denote the set of probability
measures on the Borel σ-field B(E), endowed with the weak convergence topology. Let
(µ(n)(ω,dα))n≥1 be a sequence of random probabilities defined on Ω×B(E).
(a) Assume that there exists µ(∞) ∈ P(E) such that for every bounded Lipschitz func-
tion F :E→R,
µ(n)(ω,F )
n→+∞−→ µ(∞)(F ) a.s. (24)
Then, a.s., (µ(n)(ω,dα))n≥1 converges weakly to µ
(∞) on P(E).
(b) Let U be a subset of P(E). Assume that for every sequence (Fk)k≥1 of Lipschitz
and bounded functions, a.s., for every subsequence (µ(φω(n))(ω,dα)), there exists a sub-
sequence (µ(φω◦ψω(n))(ω,dα)) and a U -valued random probability µ(∞)(ω,dα) such that
for every k ≥ 1,
µ(ψω◦φω(n))(ω,Fk)
n→+∞−→ µ(∞)(ω,Fk) a.s. (25)
Then, (µ(n)(ω,dα))n≥1 is a.s. tight with weak limits in U .
Proof. We do not give a detailed proof of the next lemma, which is essentially based
on the fact that in a separable metric space (E,d), one can build a sequence of bounded
Lipschitz functions (gk)k≥1 such that for any sequence (µn)n≥1 of probability measures
on B(E), (µn)n≥1 weakly converges to a probability µ if and only if the convergence
holds along the functions gk, k ≥ 1 (see Parthasarathy [22], Theorem 6.6, page 47 for a
very similar result). �
For every n≥ 0, for every T > 0, we introduce τ(n,T ) defined by
τ(n,T ) := min{k ≥ 0,N(k,T )≥ n}=min{k ≤ n,Γk + T ≥ Γn}. (26)
Approximation of the distribution of a stationary Markov process 159
Note that for k ∈ {0, . . . , τ(n,T )− 1}, {X̄(k)t ,0≤ t≤ T } is �FΓn -measurable and
T − γτ(n,T )−1 ≤ Γn − Γτ(n,T ) ≤ T.
Lemma 2. Assume (C3,ε) with ε < 1. Let F :D(R+,R
d)→R be a Sk-continuous func-
tional. Let (Gk) be a filtration such that F̄Γk ⊂ Gk for every k ≥ 1. Then, for any T > 0:
(a) if FT (defined by (7)) is bounded,
ηk(FT (X̄
(k−1))−E[FT (X̄(k−1))/Gk−1])
n→+∞−→ 0 a.s.; (27)
(b) if FT is not bounded, (27) holds if there exists V :Rd→R+, satisfying H(s, ε) for
some s≥ 2, such that |FT (α)| ≤C sup0≤t≤T V(αt) for every α ∈D(R+,Rd); furthermore,
ν(n)(ω,FT )<+∞ a.s. (28)
Proof. We prove (a) and (b) simultaneously. Let Υ(k) be defined by Υ(k) = FT (X̄
(k)).
We have
(k−1) −E[Υ(k−1)/Gk−1])
(k−1) −E[Υ(k−1)/Gn]) (29)
ηk(E[Υ
(k−1)/Gn]−E[Υ(k−1)/Gk−1]). (30)
We have to prove that the right-hand side of (29) and (30) tend to 0 a.s. when n→+∞.
We first focus on the right-hand side of (29). From the very definition of τ(n,T ), we
have that {X̄(k)t ,0≤ t≤ T } is F̄Γn -measurable for k ∈ {0, . . . , τ(n,T )− 1}. Hence, since
FT is σ(πs,0≤ s≤ T )-measurable and F̄Γn ⊂ Gn, it follows that Υ(k) is Gn-measurable
and that Υ(k) = E[Υ(k)/Gn] for every k ≤ τ(n,T )− 1. Then, if FT is bounded, we derive
from (C3,ε) that
(k−1) −E[Υ(k−1)/Gn])
≤ 2‖FT ‖sup
k=τ(n,T )+1
k=τ(n,T )+1
H1−εn
(Γn − Γτ(n,T ))
160 G. Pagès and F. Panloup
≤ C(T )
H1−εn
n→+∞−→ 0 a.s.,
where we used the fact that (Hn)n≥1 and (γn)n≥1 are non-decreasing and non-increasing
sequences, respectively.
Assume, now, that the assumptions of (b) are fulfilled with V satisfying H(s, ε) for
some s≥ 2 and ε < 1. By the Borel–Cantelli-like argument, it suffices to show that
k=τ(n,T )+1
(k−1) −E[Υ(k−1)/Gn])
<+∞. (31)
Let us prove (31). Let ak := η
(s−1)/s
k and bk(ω) := η
(k−1) − E[Υ(k−1)/Gn]). The
Hölder inequality applied with p̄= s/(s− 1) and q̄ = s yields
k=τ(n,T )+1
akbk(ω)
k=τ(n,T )+1
)s−1( n
k=τ(n,T )+1
ηk|Υ(k−1) −E[Υ(k−1)/Gn]|s
Now, since FT (α) ≤ sup0≤t≤T V(α), it follows from the Markov property and from
H(s, ε)(i) that
E[|FT (X̄(k))|s/F̄Γk ]≤CE
0≤t≤T
Vs(X̄(k)t )/F̄Γk
≤CTVs(X̄Γk).
Then, using the two preceding inequalities and (C3,ε) yields
k=τ(n,T )+1
(k−1) −E[Υ(k−1)/Gn])
k=τ(n,T )+1
)s−1( n
k=τ(n,T )+1
ηkE[Vs(X̄Γk−1)]
k=τ(n,T )+1
k=τ(n,T )+1
Vs(X̄Γk−1)
k=τ(n,T )+1
t∈[0,S(n,T )]
Vs(X̄τ(n,T )t )
where S(n,T ) = Γn−1 − Γτ(n,T ) and C does not depend n. By the definition of τ(n,T ),
S(n,T )≤ T . Then, again using H(s, ε)(i) yields
k=τ(n,T )
(k−1) −E[Υ(k−1)/Gn])
s(1−ε)
E[Vs(X̄(τ(n,T )))].
Approximation of the distribution of a stationary Markov process 161
Since n 7→ N(n,T ) is an increasing function, n 7→ τ(n,T ) is a non-decreasing function
and Card{n, τ(n,T ) = k}=∆N(k+1, T ) :=N(k+1, T )−N(k,T ). Then, since n 7→Hn
increases, a change of variable yields
k=τ(n,T )+1
(k−1) −E[Υ(k−1)/Gn])
∆N(k,T )
s(1−ε)
E[Vs(X̄Γk−1)]<+∞,
by H(s, ε)(iv).
Second, we prove that (30) tends to 0. For every n≥ 1, we let
(E[Υ(k−1)/Gn]−E[Υ(k−1)/Gk−1]). (32)
The process (Mn)n≥1 is a (Gn)-martingale and we want to prove that this process is
L2-bounded. Set Φ(k,n) = E[FT (X̄
(k))/Gn]− E[FT (X̄(k))/Gk]. Since FT is σ(πs,0 ≤ s ≤
T )-measurable, the random variable Φ(k,n) is F̄ΓN(k,T) -measurable. Then, for every i ∈
{N(k,T ), . . . , n}, Φ(k,n) is Gi-measurable so that
E[Φ(i,n)Φ(k,n)] =E[Φ(k,n)E[Φ(i,n)/Gi]] = 0.
It follows that
E[M2n] =
E[(Φ(k−1,n))
] + 2
N(k−1,T )∧n
i=k+1
E[Φ(i−1,n)Φ(k−1,n)]. (33)
Then,
E[M2n] ≤
E[(Φ(k−1,n))
] + 2
N(k−1,T )
i=k+1
E[Φ(i−1,n)Φ(k−1,n)]
H2−εk
E[(Φ(k−1,n))
] (34)
H2−εk
N(k−1,T )
i=k+1
γi sup
E[Φ(i−1,n)Φ(k−1,n)]
162 G. Pagès and F. Panloup
where, in the second inequality, we used assumption (C3,ε) and the decrease of i 7→
1/H1−εi . Hence, if FT is bounded, using the fact that
∑N(k−1,T )
i=k+1 γi ≤ T yields
E[M2n]≤C
H2−εk
H2−ε1
<+∞ (35)
since ε < 1. Assume, now, that the assumptions of (b) hold and let FT be dominated
by a function V satisfying H(s, ε). By the Markov property, the Jensen inequality and
H(s, ε)(i),
E[(Φ(k,n))
0≤t≤T
V2(X̄(k)t )/F̄Γk
≤CTE[V2(X̄Γk)].
We then derive from the Cauchy–Schwarz inequality that for every n, k ≥ 1, for every
i ∈ {k, . . . ,N(k,T )},
|E[Φ(i,n)Φ(k,n)]| ≤C
E[V2(X̄Γi)]
E[V2(X̄Γk)]≤C sup
t∈[0,T ]
E[V2(X̄(k)t )]≤CE[V2(X̄Γk)],
where, in the last inequality, we once again used H(s, ε)(i). It follows that
E[M2n]≤C
H2−εk
E[V2(X̄Γk−1)]<+∞,
by H(s, ε)(iii). Therefore, (34) is finite and (Mn) is bounded in L
2. Finally, we derive
from the Kronecker lemma that
ηk(E[FT (X̄
(k−1))/Gn]−E[FT (X̄(k−1))/Gk−1])
n→+∞−→ 0 a.s.
As a consequence, supn≥1 ν
(n)(ω,FT )<+∞ a.s. if and only if
E[FT (X̄
(k−1))/Fk−1]<+∞ a.s.
This last property is easily derived from H(s, ε)(i) and (ii). This completes the proof. �
Lemma 3. (a) Assume (C1) and let x0 ∈Rd. We then have limx→x0 E[d(Xx,Xx0)] = 0.
In particular, for every bounded Lispchitz (w.r.t. the distance d) functional F :D(R+,R
R, the function ΦF defined by ΦF (x) = E[F (Xx)] is a (bounded) continuous function on
(b) Assume (C2). For every compact set K ⊂Rd,
E[d(Y n,x,Xx)]
n→+∞−→ 0. (36)
Approximation of the distribution of a stationary Markov process 163
Set ΦFn (x) = E[F (Y
n,x)]. Then, for every bounded Lispchitz functional F :D(R+,R
d)→R,
|ΦF (x)−ΦFn (x)|
n→+∞−→ 0 for every compact set K ⊂Rd. (37)
Proof. (a) By the definition of d, for every α, β ∈D(R+,Rd) and for every T > 0,
d(α,β)≤
1∧ sup
0≤t≤T
|α(t)− β(t)|
+ e−T . (38)
It easily follows from assumption (C1) and from the dominated convergence theorem
limsup
E[d(Xx,Xx0)]≤ e−T for every T > 0.
Letting T →+∞ implies that limx→x0 E[d(Xx,Xx0)] = 0.
(b) We deduce from (38) and from assumption (C2) that for every compact setK ⊂Rd,
for every T > 0,
limsup
E[d(Y n,x,Xx)]≤ e−T .
Letting T →+∞ yields (36). �
Lemma 4. Assume that (ηn)n≥1 and (γn) satisfy (C3,ε) with ε < 1 and (14). Then:
(i) for every t≥ 0, for every bounded continuous function f :Rd→R,
t (ω, f)− ν
0 (ω, f)
n→+∞−→ 0 a.s.;
(ii) if, moreover, a.s., every weak limit ν(∞)(ω,dα) of (ν(n)(ω,dα))n≥1 is the dis-
tribution of a Markov process with semigroup (Qωt )t≥0, then, a.s., ν
(∞)(ω,dα) is the
distribution of a stationary process.
Proof. (i) Let f :Rd →R be a bounded continuous function. Since X̄(k)t = X̄ΓN(k,t) , we
t (ω, f)− ν
0 (ω, f) =
ηk(f(X̄ΓN(k−1,t))− f(X̄Γk−1)).
From the very definition of N(n,T ) and τ(n,T ), one checks that N(k − 1, T )≤ n− 1 if
and only if τ(n,T )≥ k. Then,
ηkf(X̄Γk−1) =
τ(n,t)
ηN(k−1,t)+1f(X̄ΓN(k−1,t))
ηkf(X̄Γk−1)1{k−1/∈N({0,...,n},t)}.
164 G. Pagès and F. Panloup
It follows that
t (ω, f)− ν
0 (ω, f) =
τ(n,t)
(ηk − ηN(k−1,t)+1)f(X̄ΓN(k−1,t))
τ(n,t)+1
ηkf(X̄ΓN(k−1,t))
ηkf(X̄Γk−1)1{k−1/∈N({0,...,n},t)}.
Then, since f is bounded and since
ηk1{k−1/∈N({0,...,n},t)} =
τ(n,t)
ηN(k−1,t)+1
τ(n,t)
|ηk − ηN(k−1,t)+1|+
k=τ(n,t)+1
we deduce that
|ν(n)t (ω, f)− ν
0 (ω, f)| ≤ 2‖f‖∞
τ(n,t)
|ηk − ηN(k−1,t)+1|+
k=τ(n,t)+1
Hence, we have to show that the sequences of the right-hand side of the preceding in-
equality tend to 0. On the one hand, we observe that
|ηk − ηN(k−1,t)+1| ≤
N(k−1,T )+1
ℓ=k+1
|ηℓ − ηℓ−1| ≤ max
ℓ≥k+1
|∆ηℓ|
N(k−1,T )+1
Using the fact that
∑N(k−1,T )+1
ℓ=k γℓ ≤ T + γ1 and condition (14) yields
τ(n,t)
|ηk − ηN(k−1,t)+1|
n→+∞−→ 0.
On the other hand, by (C3,ε), we have
k=τ(n,T )+1
H1−εn
k=τ(n,T )+1
H1−εn
n→+∞−→ 0 a.s.,
which completes the proof of (i).
Approximation of the distribution of a stationary Markov process 165
(ii) Let Q+ denote the set of non-negative rational numbers. Let (fℓ)ℓ≥1 be an every-
where dense sequence in CK(Rd) endowed with the topology of uniform convergence on
compact sets. Since Q+ and (fℓ)ℓ≥1 are countable, we derive from (i) that there exists
Ω̃⊂Ω such that P(Ω̃) = 1 and such that for every ω ∈ Ω̃, every t ∈Q+ and every ℓ≥ 1,
t (ω, fℓ)− ν
0 (ω, fℓ)
n→+∞−→ 0.
Let ω ∈ Ω̃ and let ν(∞)(ω,dα) denote a weak limit of (ν(n)(ω,dα))n≥1. We have
t (ω, fℓ) = ν
0 (ω, fℓ) ∀t ∈Q+ ∀ℓ≥ 1
and we easily deduce that
t (ω, f) = ν
0 (ω, f) ∀t ∈R+ ∀f ∈ CK(Rd).
Hence, if ν(∞)(ω,dα) is the distribution of a Markov process (Yt) with semigroup (Q
t )t≥0,
we have, for all f ∈ CK(Rd),
Qωt f(x)ν
0 (ω,dx) =
f(x)ν
0 (ω,dx) ∀t≥ 0.
0 (ω,dx) is then an invariant distribution for (Yt). This completes the proof. �
4.2. Proof of Theorem 1
Thanks to Lemma 1(a) applied with E =D(R+,R
d) and d defined by (6),
ν(n)(ω,dα)
=⇒ Pν0(dα) a.s.⇐⇒ ν(n)(ω,F )
n→+∞−→
F (x)Pν0 (dx) a.s. (39)
for every bounded Lipschitz functional F :D(R+,R
d)→ R. Now, consider such a func-
tional. By the assumptions of Theorem 1, we know that a.s., (ν
0 (ω,dx))n≥1 converges
weakly to ν0. Set Φ
F (x) := E[F (Xx)], x ∈Rd. By Lemma 3(a), ΦF is a bounded contin-
uous function on Rd. It then follows from (C0,1) that
F (X̄
(k−1)
n→+∞−→
ΦF (x)ν0(dx) =
F (x)Pν0 (dx) a.s.
Hence, the right-hand side of (39) holds for F as soon as
ηk(F (X̄
(k−1))−ΦF (X̄(k−1)0 ))
n→+∞−→ 0 a.s. (40)
166 G. Pagès and F. Panloup
Let us prove (40). First, let T > 0 and let FT be defined by (7). By Lemma 2,
ηkFT (X̄
(k−1))− 1
ηkE[FT (X̄
(k−1))/F̄Γk−1 ]
n→+∞−→ 0 a.s. (41)
With the notation of Lemma 3(b), we derive from assumption (C2)(i) that
E[FT (X̄
(k−1))/F̄Γk−1 ] = Φ
k (X̄
(k−1)
Let N ∈N. On one hand, by Lemma 3(b),
k (X̄
(k−1)
0 )−ΦFT (X̄
(k−1)
0 ))1{|X̄(k−1)
n→+∞−→ 0 a.s. (42)
On the other hand, the tightness of (ν
0 (ω,dx))n≥1 on R
d yields
ψ(ω,N) := sup
0 (ω, (B(0,N)
N→+∞−→ 0 a.s.
It follows that, a.s.,
ηk|ΦFTk (X̄
(k−1)
0 )−ΦFT (X̄
(k−1)
0 )|1{|X̄(k−1)
≤ 2‖F‖∞ψ(ω,N)
N→+∞−→ 0.
Hence, a combination of (42) and (43) yields
∀T > 0 1
k (X̄
(k−1)
0 )−ΦFT (X̄
(k−1)
n→+∞−→ 0 a.s. (44)
Finally, let (Tℓ)ℓ≥1 be a sequence of positive numbers such that, Tℓ→+∞ when ℓ→+∞.
Combining (44) and (41), we obtain that, a.s., for every ℓ≥ 1,
limsup
ηk(F (X̄
(k−1))−ΦF (X̄(k−1)))
≤ lim sup
ηk(F (X̄
(k−1))−FTℓ(X̄(k−1)))
+ limsup
FTℓ (X̄
(k−1)
0 )−ΦF (X̄
(k−1)
Approximation of the distribution of a stationary Markov process 167
By the definition of d, |F − FTℓ | ≤ e−Tℓ . Then, a.s.,
limsup
ηk(F (X̄
(k−1))−ΦF (X̄(k−1)0 ))
≤ 2e−Tℓ ∀ℓ≥ 1.
Letting ℓ→+∞ implies (40).
The generalization to non-bounded functionals in Theorem 1 is then derived from (28)
and from a uniform integrability argument.
4.3. Proof of Theorem 2
(i) We want to prove that the conditions of Lemma 1(b) are fulfilled. Since (ν
0 (ω,dx))n≥1
is supposed to be a.s. tight, one can check that for every bounded Lipschitz functional
F :D(R+,R
d)→R, (40) is still valid. Then, let (Fℓ)ℓ≥1 be a sequence of bounded Lipschitz
functionals. There exists Ω̃⊂Ω with P(Ω̃) = 1 such that for every ω ∈ Ω̃, (ν(n)0 (ω,dx))n≥1
is tight and
ηk(Fℓ(X̄
(k−1)(ω))−ΦFℓ(X̄(k−1)0 (ω)))
n→+∞−→ 0 ∀ℓ≥ 1. (45)
Let ω ∈ Ω̃ and let φω :N 7→N be an increasing function. As (ν(φω(n))0 (ω,dx))n≥1 is tight,
there exists a convergent subsequence (ν
(φω◦ψω(n))
0 (ω,dx))n≥1. We denote its weak limit
by ν∞. Since Φ
Fℓ is continuous for every ℓ≥ 1 (see Lemma 3(a)),
(φω◦ψω(n))
0 (ω,Φ
n→+∞−→ ν∞(ΦFℓ) =
Fℓ(α)Pν∞(dα) ∀ℓ≥ 1.
We then derive from (45) that for every ℓ≥ 1
ν(φω◦ψω(n))(ω,Fℓ)
n→+∞−→
Fℓ(α)Pν∞(dα).
It follows that the conditions of Lemma 1(b) are fulfilled with U = {Pµ, µ ∈ I}, where
µ ∈P(Rd),∃ω ∈ Ω̃ and an increasing function φ :N 7→N, µ= lim
ν(φ(n))(ω,dα)
Hence, by Lemma 1(b), we deduce that (ν(n)(ω,dα))n≥1 is a.s. tight with U -valued limits.
Finally, Theorem 2(ii) is a consequence of condition (14) and Lemma 4(ii).
168 G. Pagès and F. Panloup
5. Path-dependent option pricing in stationary
stochastic volatility models
In this section, we propose a simple and efficient method to price options in stationary
stochastic volatility (SSV) models. In most stochastic volatility (SV) models, the volatil-
ity is a mean reverting process. These processes are generally ergodic with a unique
invariant distribution (the Heston model or the BNS model for instance (see below) but
also the SABR model (see Hagan et al. [8]), . . .). However, they are usually considered
in SV models under a non-stationary regime, starting from a deterministic value (which
usually turns out to be the mean of their invariant distribution). However, the instanta-
neous volatility is not easy to observe on the market since it is not a traded asset. Hence,
it seems to be more natural to assume that it evolves under its stationary regime than
to give it a deterministic value at time 0.4
From a purely calibration viewpoint, considering an SV model in its SSV regime will
not modify the set of parameters used to generate the implied volatility surface, although
it will modify its shape, mainly for short maturities. This effect can in fact be an asset
of the SSV approach since it may correct some observed drawbacks of some models (see,
e.g., the Heston model below).
From a numerical point of view, considering SSV models is no longer an obstacle, es-
pecially when considering multi-asset models (in the unidimensional case, the stationary
distribution can be made more or less explicit like in the Heston model; see below) since
our algorithm is precisely devised to compute by simulation some expectations of func-
tionals of processes under their stationary regime, even if this stationary regime cannot
be directly simulated.
As a first illustration (and a benchmark) of the method, we will describe in detail
the algorithm for the pricing of Asian options in a Heston model. We will then show
in our numerical results to what extent it differs, in terms of smile and skew, from the
usual SV Heston model for short maturities. Finally, we will complete this section with
a numerical test on Asian options in the BNS model where the volatility is driven by a
tempered stable subordinator. Let us also mention that this method can be applied to
other fields of finance like interest rates, and commodities and energy derivatives where
mean-reverting processes play an important role.
4When one has sufficiently close observations of the stock price, it is in fact possible to derive a rough
idea of the size of the volatility from the variations of the stock price (see, e.g., Jacod [10]). Then, using
this information, a good compromise between a deterministic initial value and the stationary case may
be to assume that the distribution µ0 of the volatility at time 0 is concentrated around the estimated
value (see Section 2.2 for application of our algorithm in this case).
Approximation of the distribution of a stationary Markov process 169
5.1. Option pricing in the Heston SSV model
We consider a Heston stochastic volatility model. The dynamic of the asset price process
(St)t≥0 is given by S0 = s0 and
dSt = St(rdt+
(1− ρ2)vt dW 1t + ρ
vt dW
dvt = k(θ− vt) dt+ ς
vt dW
where r denotes the interest rate, (W 1,W 2) is a standard two-dimensional Brownian
motion, ρ ∈ [−1,1] and k, θ and ς are some non-negative numbers. This model was
introduced by Heston in 1993 (see Heston [9]). The equation for (vt) has a unique (strong)
pathwise continuous solution living in R+. If, moreover, 2kθ > ς
2, then (vt) is a positive
process (see Lamberton and Lapeyre [11]). In this case, (vt) has a unique invariant
probability ν0. Moreover, ν0 = γ(a, b) with a= (2k)/ς
2 and b= (2kθ)/ς2. In the following,
we will assume that (vt) is in its stationary regime, that is, that
L(v0) = ν0.
5.1.1. Option price and stationary processes
Using our procedure to price options in this model naturally needs to express the option
price as the expectation of a functional of a stationary stochastic process.
Näıve method. (may work) Since (vt)t≥0 is stationary, the first idea is to express the
option price as the expectation of a functional of (vt)t≥0: by Itô calculus, we have
St = s0 exp
rt− 1
vs ds
vs dW
1− ρ2
vs dW
. (46)
Since
vs dW
s =Λ(t, (vt)) :=
vt − v0 − kθt+ k
vs ds
it follows by setting Mt =
vs dW
s that
St =Ψ(t, (vs), (Ms)), (47)
where Ψ is given for every t≥ 0, u and w ∈ C(R+,R) by
Ψ(t, u,w) = s0 exp
rt− 1
u(s) ds
+ ρΛ(t, u) +
1− ρ2w(t)
Then, let F :C(R+,R) → R be a non-negative measurable functional. Conditioning by
FW 2T yields
E[FT ((St)t≥0)] = E[F̃T ((vt)t≥0)],
170 G. Pagès and F. Panloup
where, for every u ∈ C(R+,R),
F̃T (u) = E
t, u,
u(s) dW 1s
For some particular options such as the European call or put (thanks to the Black–
Scholes formula), the functional F̃ is explicit. In those cases, this method seems to be
very efficient (see Panloup [20] for numerical results). However, in the general case, the
computation of F̃ will need some Monte Carlo methods at each step. This approach is
then very time-consuming in general – that is why we are going to introduce another
representation of the option as a functional of a stationary process.
General method. (always works) We express the option premium as the expectation
of a functional of a two-dimensional stationary stochastic process. This method is based
on the following idea. Even though (vt,Mt) is not stationary, (St) can be expressed as a
functional of a stationary process (vt, yt). Indeed, consider the following SDE given by
dyt =−yt dt+
vt dW
dvt = k(θ− vt) dt+ ς
vt dW
First, one checks that the SDE has a unique strong solution and that assumption (S1) is
fulfilled with V (x1, x2) = 1+ x
2. This ensures the existence of an invariant distribu-
tion ν̃0 for the SDE (see, e.g., Pagès [17]). Then, since (vt) is positive and has a unique
invariant distribution, the uniqueness of the invariant distribution follows. Then, assume
that L(y0, v0) = ν̃0. Since (vt,Mt) = (vt, yt − y0 +
ys ds), we have, for every positive
measurable functional F :C(R+,R)→R,
E[FT ((St)t≥0)] = E[FT ((ψ(t, vt,Mt))t≥0)]
= Eν̃0
t, vt, yt − y0 +
ys ds
where Pν̃0 is the stationary distribution of the process (vt, yt). Every option price can
then be expressed as the expectation of an explicit functional of a stationary process. We
will develop this second general approach in the numerical tests below.
Remark 9. The idea of the second method holds for every stochastic volatility model
for which (St) can be written as follows:
St =Φ
t, vt,
hi(|vs|) dY is
, (50)
where, for every i ∈ {1, . . . , p}, hi :R+ →R is a positive function such that hi(x) = o(|x|)
as |x| → +∞, (Y it ) is a square-integrable centered Lévy process and (vt) is a mean
reverting stochastic process solution to a Lévy driven SDE.
Approximation of the distribution of a stationary Markov process 171
In some complex models, showing the uniqueness of the invariant distribution may be
difficult. In fact, it is important to note at this stage that the uniqueness of the invariant
distribution for the couple (vt, yt) is not required. Indeed, by construction, the local
martingale (Mt) does not depend on the choice of y0. It follows that if L(y0, v0) = µ̃,
with µ̃ constructed such that L(v0) = ν0, (49) still holds. This implies that it is only
necessary that uniqueness holds for the invariant distribution of the stochastic volatility
process.
5.1.2. Numerical tests on Asian options
We recall that (vt) is a Cox–Ingersoll–Ross process. For this type of processes, it is well
known that the genuine Euler scheme cannot be implemented since it does not preserve
the non-negativity of the (vt). That is why some specific discretization schemes have
been studied by several authors (Alfonsi [1], Deelstra and Delbaen [5] and Berkaoui et al.
[4, 6]). In this paper, we consider the scheme studied by the last authors in a decreasing
step framework. We denote it by (v̄t). We set v̄0 = x > 0 and
v̄Γn+1 = |v̄Γn + kγn+1(θ− v̄Γn) + ς
v̄Γn(W
−W 2Γn)|.
We also introduce the stepwise constant Euler scheme (ȳt) of (yt)t≥0 defined by
ȳΓn+1 = ȳΓn − γn+1ȳΓn +
v̄Γn(W̃
− W̃ 1Γn), ȳ0 = y ∈R
Denote by (v̄
t ) and (ȳ
t ) the shifted processes defined by v̄
t := v̄Γk+t and ȳ
ȳΓk+t, and let (ν
(n)(ω,dα))n≥1 be the sequence of empirical measures defined by
ν(n)(ω,dα) =
ηk1{(v̄(k−1),ȳ(k−1))∈dα}.
The specificity of both the model and the Euler scheme implies that Theorems 1 and 2
cannot be directly applied here. However, a specific study using the fact that (9) holds
for every compact set of R∗+ ×R when 2kθ/ς2 > 1+ 2
6/ς (see Theorem 2.2 of Berkaoui
et al. [4] and Remark 9) shows that
ν(n)(ω,dα)
=⇒ Pν̃0(dα) a.s.
when 2kθ/ς2 > 1+ 2
6/ς . Details are left to the reader.
Let us now state our numerical results obtained for the pricing of Asian options with
this discretization. We denote by Cas(ν0,K,T ) and Pas(ν0,K,T ) the Asian call and put
prices in the SSV Heston model. We have
Cas(ν0,K,T ) = e
Ss ds−K
172 G. Pagès and F. Panloup
Pas(ν0,K,T ) = e
K − 1
Ss ds
With the notation of (49), approximating Cas(ν0,K,T ) and Pas(ν0,K,T ) by our proce-
dure needs to simulate the sequences (Cnas)n≥1 and (P
as)n≥1 defined by
Cnas =
Ψ(s, v̄(k−1), M̄ (k−1)) ds−K
Pnas =
K − 1
Ψ(s, v̄(k−1), M̄ (k−1)) ds
These sequences can be computed by the method developed in Section 1.3. Note that
the specific properties of the exponential function and the linearity of the integral imply
that (
Ψ(t, v̄(n−1), M̄ (n−1)) ds) can be computed quasi-recursively.
Let us state our numerical results for the Asian call with parameters
s0 = 50, r = 0.05, T = 1, ρ= 0.5,
θ = 0.01, ς = 0.1, k = 2.
We also assume that K ∈ {44, . . . ,56} and choose the following steps and weights: γn =
ηn = n
−1/3. In Table 1, we first state the reference value for the Asian call price obtained
for N = 108 iterations. In the two following lines, we state our results for N = 5.104 and
N = 5.105 iterations. Then, in the last lines, we present the numerical results obtained
Table 1. Approximation of the Asian call price
K 44 45 46 47 48 49 50
Asian call (ref.) 6.92 5.97 5.04 4.12 3.25 2.46 1.78
N = 5 · 104 6.89 6.07 5.07 4.13 3.18 2.49 1.77
N = 5 · 105 6.90 6.02 5.00 4.11 3.24 2.46 1.79
N = 5 · 104 (CP parity) 6.92 5.96 5.04 4.13 3.26 2.46 1.78
N = 5 · 105 (CP parity) 6.92 5.97 5.04 4.12 3.25 2.47 1.78
K 51 52 53 54 55 56
Asian call (ref.) 1.23 0.82 0.53 0.33 0.21 0.12
N = 5 · 104 1.21 0.81 0.51 0.34 0.22 0.11
N = 5 · 105 1.23 0.82 0.53 0.33 0.21 0.13
N = 5 · 104 (CP parity) 1.23 0.82 0.53 0.31 0.21 0.12
N = 5 · 105(CP parity) 1.23 0.82 0.53 0.33 0.21 0.13
Approximation of the distribution of a stationary Markov process 173
using the call-put parity
Cas(ν0,K,T )− Pas(ν0, S0,K,T ) =
(1− e−rT )−Ke−rT (52)
as a means of variance reduction. The computation times for N = 5.104 and N = 5.105
(using MATLAB with a Xeon 2.4 GHz processor) are about 5 s and 51 s, respectively. In
particular, the complexity is quasi-linear and the additional computations needed when
we use the call-put parity are negligible.
5.2. Implied volatility surfaces of Heston SSV and SV models
Given a particular pricing model (with initial value s0 and interest rate r) and its asso-
ciated European call prices denoted by Ceur(K,T ), we recall that the implied volatility
surface is the graph of the function (K,T ) 7→ σimp(K,T ), where σimp(K,T ) is defined for
every maturity T > 0 and strike K as the unique solution of
CBS(s0,K,T, r, σimp(K,T )) =Ceur(K,T ),
where CBS(s0,K,T, r, σ) is the price of the European call in the Black–Scholes model
with parameters s0, r and σ. When Ceur(K,T ) is known, the value of σimp(K,T ) can be
numerically computed using the Newton method or by dichotomy if the first method is
not convergent.
In this last part, we compare the implied volatility surfaces induced by the SSV and SV
Heston models where we suppose that the initial value of (vt) in the SV Heston model is
the mean of the invariant distribution, that is, we suppose that v0 = θ.
5 We also assume
that the parameters are those of (51), except the correlation coefficient ρ.
In Figures 1 and 2, the volatility curves obtained when T = 1 are depicted, whereas in
Figures 3 and 4, we set the strikeK atK = 50 and let the time vary. These representations
show that when the maturity is long, the differences between the SSV and SV Heston
models vanish. This is a consequence of the convergence of the stochastic volatility to its
stationary regime when T →+∞.
The main differences between these models then appear for short maturities. That is
why we complete this part by a representation of the volatility curve when T = 0.1 for
ρ= 0 and ρ= 0.5 in Figures 5 and 6, respectively. We observe that for short maturities,
the volatility smile is more curved and the skew is steeper. These phenomena seem
interesting for calibration since one well-known drawback of the standard Heston model
is that it can have overly flat volatility curves for short maturities.
5.3. Numerical tests on Asian options in the BNS SSV model
The BNS model introduced in Barndorff-Nielsen and Shephard [3] is a stochastic volatility
model where the volatility process is a Lévy-driven positive Ornstein–Uhlenbeck process.
5This choice is the most usual in practice.
174 G. Pagès and F. Panloup
Figure 1. ρ= 0, K 7→ σimp(K,1).
The dynamic of the asset price (St) is given by St = S0 exp(Xt),
dXt = (r− 12vt) dt+
vt dWt + ρdZt, ρ≤ 0,
dvt = −µvt dt+dZt, µ > 0,
Figure 2. ρ= 0.5, K 7→ σimp(K,1).
Approximation of the distribution of a stationary Markov process 175
Figure 3. ρ= 0, T 7→ σimp(50, T ).
where (Zt) is a subordinator without drift term and Lévy measure π. In the following,
we assume that (Zt) is a tempered stable subordinator, that is, that
π(dy) = 1{y>0}
c exp(−λy)
dy, c > 0, λ > 0, α∈ (0,1).
As in the Heston model, we want to use our algorithm as a way of option pricing when
the stochastic volatility evolves under its stationary regime and test it on Asian options
using the method described in detail in Section 5.1. This model does not require a specific
Figure 4. ρ= 0.5, T 7→ σimp(50, T ).
176 G. Pagès and F. Panloup
Figure 5. ρ= 0, T 7→ σimp(50, T ).
discretization and the approximate Euler scheme (P) (see Section 3.2) relative to (vt)
can be implemented using the rejection method. In Table 2, we present our numerical
results obtained for the following choices of parameters, steps and weights:
ρ=−1, λ= µ= 1, c= 0.01, α= 1
, γn = ηn = n
−1/3.
The computation times forN = 5.104 andN = 5.105 are about 8.5 s and 93 s, respectively.
Note that for this model, the convergence seems to be slower because of the approximation
of the jump component.
Figure 6. ρ= 0.5, T 7→ σimp(50, T ).
Approximation of the distribution of a stationary Markov process 177
Table 2. Approximation of the Asian call price in the BNS model
K 44 45 46 47 48 49 50
Asian call (ref.) 6.75 5.83 4.93 4.05 3.18 2.35 1.57
N = 5 · 104 6.83 5.91 5.01 4.10 3.22 2.35 1.51
N = 5 · 105 6.78 5.86 4.96 4.06 3.19 2.34 1.52
N = 5 · 104 (CP parity) 6.76 5.85 4.94 4.07 3.20 2.29 1.51
N = 5 · 105 (CP parity) 6.75 5.83 4.93 4.04 3.17 2.32 1.54
K 51 52 53 54 55 56
Asian call (ref.) 0.91 0.55 0.39 0.29 0.23 0.18
N = 5 · 104 0.77 0.46 0.33 0.27 0.22 0.19
N = 5 · 105 0.79 0.48 0.34 0.27 0.21 0.17
N = 5 · 104 (CP parity) 0.79 0.47 0.37 0.27 0.23 0.19
N = 5 · 105(CP parity) 0.83 0.50 0.36 0.28 0.22 0.17
Acknowledgement
The authors would like to thank Vlad Bally for interesting comments on the paper.
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Introduction
Objectives and motivations
Background and construction of the procedure
Simulation of ((n)(,F))n1
General results
Weak convergence to the stationary regime
Extension to the non-stationary case
Application to Brownian diffusions and Lévy-driven SDE's
Application to Brownian diffusions
Application to Lévy-driven SDE's
Proofs of Theorems 1 and 2
Preliminary lemmas
Proof of Theorem 1
Proof of Theorem 2
Path-dependent option pricing in stationary stochastic volatility models
Option pricing in the Heston SSV model
Option price and stationary processes
Numerical tests on Asian options
Implied volatility surfaces of Heston SSV and SV models
Numerical tests on Asian options in the BNS SSV model
Acknowledgement
References
|
0704.0336 | Influence of Phonon dimensionality on Electron Energy Relaxation | Influence of Phonon dimensionality on Electron Energy Relaxation
J. T. Karvonen and I. J. Maasilta
Nanoscience Center, Department of Physics, P.O. Box 35, FIN-40014 University of Jyväskylä, Finland.
We studied experimentally the role of phonon dimensionality on electron-phonon (e-p) interaction
in thin copper wires evaporated either on suspended silicon nitride membranes or on bulk substrates,
at sub-Kelvin temperatures. The power emitted from electrons to phonons was measured using
sensitive normal metal-insulator-superconductor (NIS) tunnel junction thermometers. Membrane
thicknesses ranging from 30 nm to 750 nm were used to clearly see the onset of the effects of two-
dimensional (2D) phonon system. We observed for the first time that a 2D phonon spectrum clearly
changes the temperature dependence and strength of the e-p scattering rate, with the interaction
becoming stronger at the lowest temperatures below ∼ 0.5 K for the 30 nm membranes.
PACS numbers: 63.22.+m, 63.20.Kr, 85.85.+j
It is an established fact that at sub-Kelvin tempera-
tures the thermal coupling between conduction electrons
and the lattice becomes very weak [1]. This has signifi-
cant implications for the operation of low-temperature
detectors and coolers [2], or for any solid-state sys-
tems where dissipation and cooling are relevant. Low-
temperature electron-phonon (e-p) interaction has been
studied widely during the past decades, but mostly only
for the case in which the phonons are fully three di-
mensional (3D) [3, 4, 5, 6]. However, due to signifi-
cant advances in fabrication of thin suspended structures,
many practical devices and detectors exist in which the
phonons are expected to move freely only within the
plane of a membrane, forming a quasi-2D system [7].
The question how the two-dimensionality of the phonon
modes influences e-p interaction has been addressed the-
oretically for certain cases [8, 9, 10], but no clear exper-
imental observation of the effect has been reported to
date, although several attempts have been made [11, 12].
In this paper, we show for the first time experimen-
tally that the electron-phonon interaction clearly changes
depending on the dimensionality of the phonons, as ex-
pected from theory. E-p coupling was measured with the
help of sensitive NIS tunnel junction thermometry [13],
for thin Cu wires on suspended silicon nitride (SiNx)
membranes with thickness varying from 30 nm to 750
nm, which spans the transition from 2D to 3D phonons.
In addition, samples with identical Cu wires on bulk
substrates were also measured for comparison. For the
thinnest membranes, the e-p interaction was strengthened
in comparison with the bulk samples, and its tempera-
ture dependence changed significantly, as is predicted by
the theory [8, 9, 10]. The change was large enough to
give indirect evidence that the dispersive (ω ∼ k2), flex-
ural modes of the membrane likely play a major role in
the e-p interaction.
In the presence of stress-free boundaries, the bulk
transversal and longitudinal phonon modes (with sound
velocities ct and cl, respectively) couple to each other
and form a new set of eigenmodes, which in the case
of a suspended membrane are known as the horizontal
shear modes (h), and symmetric (s) and antisymmet-
ric (a) Lamb modes [14]. The frequencies ω for the h
modes are simply ω = ct
+ (mπ/d)2, where k‖ is the
wave vector component parallel to the membrane sur-
faces, d is the membrane thickness and the integer m is
the branch number. However, the dispersion relations of
the s and a Lamb modes cannot be given in a closed an-
alytical form, but have to be calculated numerically. The
lowest three branches, dominant for thin membranes at
low temperatures, have low frequency analytical expres-
sions: ωh = ctk‖, ωs = csk‖, and ωa =
k2‖, where
cs = 2ct
− c2t )/c
is the effective sound velocity of
the s mode, and m⋆ = ~
− c2t )/3c
is an
effective mass for the a-mode ”particle”. This lowest a-
mode with its quadratic dispersion is mostly responsible
for the non-trivial behavior of the e-p interaction [9, 10].
Note that already a single free surface affects the modes
[15] and the e-p interaction [16], as the bulk modes cou-
ple and form another new set of eigenstates, including
the surface localized Rayleigh-mode. Thus, the widely
observed result for e-p power flow P = ΣV (T 5e − T
from a metal volume V with Te the electron and Tp the
phonon temperature, is not expected to hold even for
thin enough films on bulk substrates.
A schematic of the Cu wire samples on suspended sil-
icon nitride membranes and the used measuring circuit
is shown in Fig. 1. 17 samples were made on either sus-
pended membranes or bulk substrates, where nitridized
(100) Si wafers with 30, 200 and 750 nm thick low-stress
SiNx top layers were used as the substrate for both cases.
The suspension of the SiNx membranes (size 600×300
µm2) was achieved by anisotropic backside wet etching of
the silicon substate in KOH, and the metallic structures
were fabricated using standard e-beam lithography and
multi-angle shadow mask evaporation techniques. As the
e-p interaction strength is sensitive to the thickness and
disorder level of the metal [17], we minimized its effect
by evaporating the Cu wires of a specific thickness on
all the different substrates simultaneously. Ultrathin Cu
layers (t=14-30 nm) were used to strengthen the effect of
the thin membranes. The oxide layer forming the tun-
nel junction barriers was produced by thermal oxidation
of Al. Table I presents the essential dimensions of the
http://arxiv.org/abs/0704.0336v2
samples discussed in this paper, measured by scanning
electron (SEM) and atomic force (AFM) microscopies.
The electron mean free path l was determined from the
resistance of the wire at base temperature 60 mK, using
the accurately measured dimensions of the wire.
TABLE I: Parameters for samples. M= suspended SiNx mem-
brane and B= bulk substrate. B6 had an oxidized Si sub-
strate.
Sample SiNx d Cu t V l τ (0.2K) τ (0.8K)
(nm) (nm) [(µm)3] (nm) (µs) (µs)
M1 30 14 2.71 5.7 2.6 0.16
B1 30 14 2.46 4.9 7.1 0.030
M2 200 14 2.44 4.6 15.0 0.11
B2 200 18 3.67 4.1 6.4 0.045
M3 30 19 5.50 11.2 2.2 0.30
B3 30 19 4.62 9.8 4.3 0.034
M4 750 22 6.09 10.3 3.1 0.030
B4 750 22 5.87 8.7 3.9 0.013
M5 30 32 6.09 22 1.8 0.31
B5 30 32 5.09 19 2.7 0.038
B6 - 32 7.10 22 1.6 0.031
CuAl Nb/Al
FIG. 1: (Color online) A Schematic of the suspended samples
and the measuring circuit. Red lines are the normal metal
Cu, light gray Al for SINIS-junctions and dark gray Al or Nb
for SN-junctions.
We used the hot-electron technique [3] to measure the
e-p interaction by overheating the electrons by Joule heat
power P and measuring the resulting electron tempera-
ture Te. All the samples had two electrically isolated
Cu normal metal wires next to each other (Fig. 1). The
longer wire (L = 500µm) was heated by applying a slowly
ramping voltage across the pair of superconducting Nb
(or Al) leads in direct metallic contact to Cu, forming
SN junctions. These junctions provide excellent electri-
cal, but very poor thermal conductance due to Andreev
reflection, as the junctions are biased within the super-
conducting gap ∆. Thus, due to the lack of outdiffu-
sion of electrons and the long length of the wire, input
heat is distributed uniformly in the interior of the wire
and the electron gas cools dominantly by phonons, in-
stead of diffusively [18] or by thermal photons [19]. Since
L >> Le−e, the electron-electron scattering length, elec-
tron temperature is also well defined without complica-
tions from non-equilibrium [20]. In our sample geometry
the electron temperature is measured with two additional
Al leads forming a NIS tunnel junctions pair (SINIS) in
the middle of heated wire, as a function of input Joule
power P = IV measured in a four probe configuration.
The purpose of the short Cu wire, with additional SI-
NIS thermometer on it, is to give an estimate of the local
phonon temperature Tp, as the e-p power flow depends
on both Te and Tp.
The current-biased Al SINIS thermometer is ideally
suited to measure temperature below a few Kelvins, [2]
due to its high sensitivity (in our DC measurement ∼ 0.1
mK at 0.1 K) and low power dissipation. In addition,
for all the data here, the SINIS voltage vs. temperature
response follows the BCS theory without fitting param-
eters very accurately at least down to ∼ 0.2 K, where
typically saturation sets in. This saturation depends on
the strength of the e-p interaction (size of thermometer
and type of substrate) and the amount of filtering, and
thus we conclude that it is most likely caused by external
noise heating. For this reason we take the most conser-
vative approach and assume that all saturation is caused
by it, in which case we can use BCS theory to convert
the measured voltage data for all temperatures.
Even if the electrons lose their energy overwhelmingly
to the phonons in our sample geometry, it is still pos-
sible that the measured temperature is not only deter-
mined by the e-p interaction. This is because the emitted
phonons could be removed so ineffectively from the mem-
brane that the phonon transmission becomes a bottleneck
for the energy flow. Bulk scattering of phonons at low
temperatures is very weak [7], even for thin disordered
membranes [21], as is boundary resistance for thin films
on bulk substrates [22, 23]. In contrast, almost noth-
ing quantitative is known about the boundary resistance
between a thin metal film and a thin 2D membrane, or
between a thin 2D membrane and a bulk substrate. How-
ever, it seems clear that if the combined metal film and
membrane thickness is below the thermal wavelength of
the phonons, the phonon modes in the two materials are
strongly coupled, leading to an effectively non-existent
boundary resistance. Hence, if we check that the mem-
brane temperature Tp is not too high compared to Te
(effective enough hot phonon removal), we can be confi-
dent that the measured Te reflects the e-p interaction.
Figure 2 shows the main result of the measurements,
with Te and Tp plotted vs. the heating power density
p = P/V for all membrane thicknesses (30 nm, 200 nm
and 750 nm). In addition, data from a few represen-
tative bulk samples are shown. Compared to the cor-
responding bulk substrate sample (B4), Te of the 750
nm membrane (M4) shows no difference at all, and it
effectively behaves as bulk. This is reasonable, because
for the 750 nm membrane the estimated dimensional-
ity cross-over temperature [24, 25] Tcr = ~ct/(2kBd) is
∼ 30 mK, with ct = 6200 m/s for SiN. The phonon
temperatures Tp, however, show a big difference: The
0.1 1 10 100 1000
of M3 T
of M2
of M4 T
of B1-B6
of M1
of M2
of M4
of B1 and B2
of B4
Heating power density [pW / ( m)3]
FIG. 2: (Color online) Measured electron and phonon tem-
peratures Te and Tp versus the applied heating power density
in log-log-scale.
bulk samples show almost no response from the satura-
tion value of the thermometer ∼ 190 mK, whereas the
membrane phonons heat up measurably, most likely due
to the boundary resistance between the membrane and
the bulk. Nevertheless, this increase in Tp for all sam-
ples is small enough not to influence the e-p interaction.
For the 200 nm thick membrane (M2) (Tcr ∼ 110 mK),
at low heating power densities [p < 40 pW/(µm)3] the
temperature dependence follows the behavior of the bulk
sample (B2), although with a difference in the absolute
value. This shows that the strength of the e-p coupling
weakens compared to the bulk. At higher powers and
temperatures (p > 40 pW/(µm)3, where Te > 0.6 K),
Te starts to increase more rapidly in the membrane sam-
ple, most likely due to the boundary resistance effects.
The phonons in the 30 nm thick membrane sample (M1)
are expected to be in the 2D limit at low temperatures
(Tcr ∼ 0.5K), and a clear sign of this can be seen in
Fig. 2 as a strongly different behavior of the measured
Te vs. p curve with respect to all other samples. Below
∼ 6 pW/(µm)3 the e-p coupling is notably stronger (Te
lower) than in the corresponding bulk (B1) or any other
sample, but again at highest temperatures the influence
of other effects starts to dominate over the e-p coupling.
To study the temperature dependence of the data in
Fig. 2 more accurately, we plot the logarithmic deriva-
tives d(log p)/d(logTe) in Fig. 3 (a)-(c). For low heat-
ing powers (T ne >> T
p ) Pe−p ≈ T
e , where n is the
power law of the e-p interaction, thus in that regime
d(log p)/d(logTe) = n. Typically this exponent is n ≈ 5
for thicker (t > 30 nm) metal films on bulk substrates
[3, 4, 17], if the disorder in the film is not too strong
[26, 27, 28]. From Fig. 3 (a) we first of all see that
for the 30 nm membrane sample M1, the difference to
the bulk sample B1 is very clear. The M1 data has a
0.1 1 10 100 1000
M1 B1
M2 B2
M4 B4
Heating power density [pW / ( m)3]
FIG. 3: (Color online) Numerical logarithmic derivatives of
the measured data in Fig. 2. (a) Te data for M1 and B1, (b)
Te data for M2 and B2, (c) Te data for M4 and B4.
plateau of n ∼ 4.5 between p = 0.1 - 6 pW/(µm)3, while
for B1, n continuously decreases from much higher val-
ues. Note that the strong increase of d(log p)/d(logTe)
below p ∼ 0.1 pW/(µm)3 is caused by the saturation
of the Te measurement, and not by the e-p interaction.
The point where n starts deviating from n = 4.5 cor-
responds to Te ≈ 0.4 K, which is surprisingly consistent
with the estimated Tcr ∼ 0.5 K. In contrast, the tempera-
ture dependence of the 200 nm membrane (M2) and bulk
(B2) samples [Fig. 3 (b)] are identical with each other
and with the 30 nm bulk sample (B1), as long as the
e-p interaction is dominant (up to 40 pW/(µm)3). The
750 nm membrane (M4) and bulk (B4) samples also give
identical values of n [Fig. 3 (c)]. The difference between
sample pairs M4,B4 and M2,B2 is caused by the Cu wire
thickness, which is expected to influence the temperature
dependence strongly [16, 27].
Finally, we discuss the effect of the Cu wire thickness
on the measured e-p interaction. The results for the
thinnest 30 nm membrane samples, with Cu thickness
t = 14,19 and 32 nm are shown in Figs 4 (a) and (c).
It is apparent that the metal film thickness has only a
minor effect on the e-p interaction on thin membranes,
and only influences the boundary resistance in the 3D
limit, by increasing its effect for thicker t, as expected.
However, for wires on bulk substrates, Figs 4 (b) and (d),
the effect of the Cu wire thickness on e-p interaction is
more profound. The thinner the Cu film, the more its
temperature dependence deviates from n = 5, which, for
comparison, is observed for a more typical t = 32 nm
Cu wire on oxidized Si (B6). This behavior is qualita-
0.1 1 10 100 1000
0.1 1 10 100 1000
(a) (b)
Heating power density [pW /( m)3]
FIG. 4: (Color online) (a) Te versus p = P/V for 30 nm mem-
brane samples M1,M3,M5. (b) Te versus p for bulk samples,
from top to bottom B1 (top), B3, B5 and B6 (bottom). (c)
d(log p)/d(log T ) of the data in (a). (d) d(log p)/d(log T ) of
the data in (b). From top to bottom: green line B1 (top),
magenta B3, blue B5, Red B6 (bottom). In (d) noise has
been filtered to help the eye.
tively consistent with the predicted effect of the surface
phonon modes [16], but could also depend on the disor-
der, as the thickening of the film increases the mean free
path l (Table I) and pushes the sample closer to the clean
limit. An apparent exponent as high as ∼ 7 could pos-
sibly be explained by the combination of strong disorder
and surface modes, but again, detailed theory is lacking.
In conclusion, we have obtained the first clear evidence
that the electron-phonon interaction at low tempera-
tures changes quite significantly when the phonon modes
become two-dimensional. To quantify the effects, the
electron thermal relaxation times τ = γV Te/(dP/dTe),
where γ = 100 J/K2m3 for Cu, are presented in Table I
for all the samples at two temperatures Te = 0.2 and 0.8
K. At Te < 0.5 K, the thinnest membranes can have a
a factor 2-3 strengthening effect, whereas at higher tem-
peratures the thermal relaxation from membranes can be
an order of magnitude weaker compared to bulk samples.
The membrane close to transition region (d=200 nm) was
shown to have a weaker (∼ factor of two) e-p interaction
strength than the bulk samples. Thinning the metal film
on bulk substrates also leads to a sizeable weakening of
the e-p interaction. The observed power law exponent
for the 2D limit is consistent with n ≈ 4.5, and is much
smaller than the corresponding bulk exponent n = 6..7.
A reduction by more than a factor one gives indirect evi-
dence of the importance of the flexural, dispersive Lamb-
modes for the membrane electron-phonon interaction, in
agreement with theory [9, 10].
Discussions with T. Kühn and A. Sergeev and tech-
nical assistance by H. Niiranen are acknowledged. This
work was supported by the Academy of Finland project
Nos. 118665 and 118231, and by the Finnish Academy
of Sciences and Letters (J.T.K.).
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http://arxiv.org/abs/0705.1936
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|
0704.0337 | Bursting Dynamics of the 3D Euler Equations in Cylindrical Domains | Bursting Dynamics of the 3D Euler Equations
in Cylindrical Domains
François Golse ∗ †
Ecole Polytechnique, CMLS
91128 Palaiseau Cedex, France
Alex Mahalov ‡and Basil Nicolaenko §
Department of Mathematics and Statistics
Arizona State University
Tempe, AZ 85287-1804, USA
Abstract
A class of three-dimensional initial data characterized by uniformly
large vorticity is considered for the 3D incompressible Euler equations in
bounded cylindrical domains. The fast singular oscillating limits of the
3D Euler equations are investigated for parametrically resonant cylin-
ders. Resonances of fast oscillating swirling Beltrami waves deplete the
Euler nonlinearity. These waves are exact solutions of the 3D Euler
equations. We construct the 3D resonant Euler systems; the latter are
countable uncoupled and coupled SO(3;C) and SO(3;R) rigid body
systems. They conserve both energy and helicity. The 3D resonant Eu-
ler systems are vested with bursting dynamics, where the ratio of the
enstrophy at time t = t∗ to the enstrophy at t = 0 of some remarkable
orbits becomes very large for very small times t∗; similarly for higher
norms Hs, s ≥ 2. These orbits are topologically close to homoclinic
cycles. For the time intervals where Hs norms, s ≥ 7/2 of the limit
resonant orbits do not blow up, we prove that the full 3D Euler equa-
tions possess smooth solutions close to the resonant orbits uniformly in
strong norms.
Key-Words: Incompressible Euler Equations, Rotating Fluids, Rigid Body
Dynamics, Enstrophy Bursts
MSC: 35Q35, 76B03, 76U05
∗[email protected]
†and Laboratoire J.-L. Lions, Université Paris Diderot-Paris 7
‡[email protected]
§[email protected]
http://arxiv.org/abs/0704.0337v1
1 Introduction
The issues of blowup of smooth solutions and finite time singularities of
the vorticity field for 3D incompressible Euler equations are still a major
open problem. The Cauchy problem in 3D bounded axisymmetric cylindrical
domains is attracting considerable attention: with bounded, smooth, non-
axisymmetric 3D initial data, under the constraints of conservation of bounded
energy, can the vorticity field blow up in finite time? Outstanding numeri-
cal claims for this have recently been disproven [Ke], [Hou1], [Hou2]. The
classical analytical criterion of Beale-Kato-Majda [B-K-M] for non-blow up
in finite time requires the time integrability of the L∞ norm of the vorticity.
DiPerna and Lions [Li] have given examples of global weak solutions of the
3D Euler equations which are smooth (hence unique) if the initial conditions
are smooth (specifically in W1,p(D), p > 1). However, these flows are really
2-Dimensional in x1, x2, 3-components flows, independent from the third co-
ordinate x3. Their examples [DiPe-Li] show that solutions (even smooth ones)
of the 3D Euler equations cannot be estimated in W1,p for 1 < p <∞ on any
time interval (0, T ) if the initial data are only assumed to be bounded inW1,p.
Classical local existence theorems in 3D bounded or periodic domains by Kato
[Ka], Bourguignon-Brézis [Bou-Br] and Yudovich [Yu1], [Yu2] require some
minimal smoothness for the initial conditions (IC), e.g., in Hs(D), s > 5
The classical formulation for the Euler equations is
∂tV+ (V · ∇)V = −∇p, ∇ ·V = 0, (1.1)
V ·N = 0 on ∂D, (1.2)
where ∂D is the boundary of a bounded, connected domain D, N the normal
to ∂D, V(t, y) = (V1, V2, V3) the velocity field, y = (y1, y2, y3), and p is the
pressure.
The equivalent Lamé form [Ar-Khe]
∂tV + curlV ×V +∇
= 0, (1.3)
∇ ·V = 0, (1.4)
∂tω + curl(ω ×V) = 0, (1.5a)
ω = curlV, (1.5b)
implies conservation of Energy:
E(t) =
|V(t, y)|2 dy. (1.6)
The helicity Hel(t) [Ar-Khe], [Mof], is conserved:
Hel(t) =
V · ω dy, (1.7)
for D = R3 and when D is a periodic lattice. Helicity is also conserved for
cylindrical domains, provided that ω·N = 0 on the cylinder’s lateral boundary
at t = 0 (see [M-N-B-G]).
From the theoretical point of view, the principal difficulty in the analysis
of 3D Euler equations is due to the presence of the vortex stretching term
(ω · ∇)V in the vorticity equation (1.5a). The equations (1.3) and (1.5a) are
equivalent to:
∂tω + [ω,V] = 0, (1.8)
where [a, b] = curl (a × b) is the commutator in the infinite dimensional Lie
algebra of divergence-free vector fields [Ar-Khe]. This point of view has led to
celebrated developments in Topological Methods in Hydrodynamics [Ar-Khe],
[Mof]. The striking analogy between the Euler equations for hydrodynamics
and the Euler equations for a rigid body (the latter associated to the Lie
Algebra of the Lie group SO(3,R)) had already been pointed out by Moreau
[Mor1]; Moreau was the first to demonstrate conservation of Helicity (1961)
[Mor2]. This has led to extensive speculations to what extent/in what cases
are the solutions of the 3D Euler equations “close” to those of coupled 3D rigid
body equations in some asymptotic sense. Recall that the Euler equations for
a rigid body in R3 is:
mt + ω ×m = 0, m = Aω, (1.9a)
mt + [ω,m] = 0, (1.9b)
where m is the vector of angular momentum relative to the body, ω the
angular velocity in the body and A the inertia operator [Ar1], [Ar-Khe].
The Russian school of Gledzer, Dolzhansky, Obukhov [G-D-O] and Vishik
[Vish] has extensively investigated dynamical systems of hydrodynamic type
and their applications. They have considered hydrodynamical models built
upon generalized rigid body systems in SO(n,R), following Manakhov [Man].
Inspired by turbulence physics, they have investigated “shell” dynamical sys-
tems modeling turbulence cascades; albeit such systems are flawed as they
only preserve energy, not helicity. To address this, they have constructed and
studied in depth n-dimensional dynamical systems with quadratic homoge-
neous nonlinearities and two quadratic first integrals F1, F2. Such systems
can be written using sums of Poisson brackets:
i2,...,in
ǫi1i2...inpi4...in
− ∂F1
, (1.10)
where constants pi4...in are antisymmetric in i4, ..., in.
A simple version of such a quadratic hydrodynamic system was introduced
by Gledzer [Gl1] in 1973. A deep open issue of the work by the Gledzer-
Obukhov school is whether there exist indeed classes of I.C. for the 3D Cauchy
Euler problem (1.1) for which solutions are actually asymptotically close in
strong norm, on arbitrary large time intervals to solutions of such hydro-
dynamic systems, with conservation of both energy and helicity. Another
unresolved issue is the blowup or global regularity for the “enstrophy” of such
systems when their dimension n→ ∞.
This article reviews some current new results of a research program in
the spirit of the Gledzer-Obukhov school; this program builds-up on the re-
sults of [M-N-B-G] for 3D Euler in bounded cylindrical domains. Following
the original approach of [B-M-N1]-[B-M-N4] in periodic domains, [M-N-B-G]
prove the non blowup of the 3D incompressible Euler equations for a class
of three-dimensional initial data characterized by uniformly large vorticity in
bounded cylindrical domains. There are no conditional assumptions on the
properties of solutions at later times, nor are the global solutions close to
some 2D manifold. The initial vortex stretching is large. The approach of
proving regularity is based on investigation of fast singular oscillating limits
and nonlinear averaging methods in the context of almost periodic functions
[Bo-Mi], [Bes], [Cor]. Harmonic analysis tools based on curl eigenfunctions
and eigenvalues are crucial. One establishes the global regularity of the 3D
limit resonant Euler equations without any restriction on the size of 3D initial
data. The resonant Euler equations are characterized by a depleted nonlin-
earity. After establishing strong convergence to the limit resonant equations,
one bootstraps this into the regularity on arbitrary large time intervals of the
solutions of 3D Euler Equations with weakly aligned uniformly large vorticity
at t = 0. [M-N-B-G] theorems hold for generic cylindrical domains, for a set
of height/radius ratios of full Lebesgue measure. For such cylinders, the 3D
limit resonant Euler equations are restricted to two-wave resonances of the
vorticity waves and are vested with an infinite countable number of new con-
servation laws. The latter are adiabatic invariants for the original 3D Euler
equations.
Three-wave resonances exist for a nonempty countable set of h/R (h
height, R radius of the cylinder) and moreover accumulate in the limit of
vanishingly small vertical (axial) scales. This is akin to Arnold tongues [Ar2]
for the Mathieu-Hill equations and raises nontrivial issues of possible sin-
gularities/lack thereof for dynamics ruled by infinitely many resonant triads
at vanishingly small axial scales. In such a context, the 3D resonant Euler
equations do conserve the energy and helicity of the field.
In this review, we consider cylindrical domains with parametric resonances
in h/R and investigate in depth the structure and dynamics of 3D resonant
Euler systems. These parametric resonances in h/R are proven to be non-
empty. Solutions to Euler equations with uniformly large initial vorticity are
expanded along a full complete basis of elementary swirling waves (T2 in time).
Each such quasiperiodic, dispersive vorticity wave is a quasiperiodic Beltrami
flow; these are exact solutions of 3D Euler equations with vorticity parallel
to velocity. There are no Galerkin-like truncations in the decomposition of
the full 3D Euler field. The Euler equations, restricted to resonant triplets of
these dispersive Beltrami waves, determine the “resonant Euler systems”. The
basic “building block” of these (a priori ∞-dimensional) systems are proven
to be SO(3;C) and SO(3;R) rigid body systems:
U̇k + (λm − λn)UmUn = 0
U̇m + (λn − λk)UnUk = 0
U̇n + (λk − λm)UkUm = 0
(1.11)
These λ’s are eigenvalues of the curl operator in the cylinder, curlΦ±n =
±λnΦ±n ; the curl eigenfunctions are steady elementary Beltrami flows, and
the dispersive Beltrami waves oscillate with the frequencies ± h
, n3 ver-
tical wave number (vertical shear), 0 < ǫ < 1. Physicists [Ch-Ch-Ey-H] have
computationally demonstrated the physical impact of the polarization of Bel-
trami modes Φ± on intermittency in the joint cascade of energy and helicity
in turbulence.
Another “building block” for resonant Euler systems is a pair of SO(3;C)
or SO(3;R) rigid bodies coupled via a common principal axis of inertia/mo-
ment of inertia:
ȧk = (λm − λn)Γaman (1.12a)
ȧm = (λn − λk)Γanak (1.12b)
ȧn = (λk − λm)Γakam + (λk̃ − λm̃)Γ̃ak̃am̃ (1.12c)
ȧm̃ = (λn − λk̃)Γ̃anak̃ (1.12d)
ȧk̃ = (λm̃ − λn)Γ̃am̃an, (1.12e)
where Γ and Γ̃ are parameters in R defined in Theorem 4.10. Both reso-
nant systems (1.11) and (1.12) conserve energy and helicity. We prove that
the dynamics of these resonant systems admit equivariant families of homo-
clinic cycles connecting hyperbolic critical points. We demonstrate bursting
dynamics: the ratio
||u(t)||2Hs/||u(0)||2Hs , s ≥ 1
can burst arbitrarily large on arbitrarily small times, for properly chosen para-
metric domain resonances h/R. Here
||u(t)||2Hs =
2s|un(t)|2 . (1.13)
The case s = 1 is the enstrophy. The “bursting” orbits are topologically close
to the homoclinic cycles.
Are such dynamics for the resonant systems relevant to the full 3D Euler
equations (1.1)-(1.8)? The answer lies in the following crucial “shadowing”
Theorem 2.10. Given the same initial conditions, given the maximal time
interval 0 ≤ t < Tm where the resonant orbits of the resonant Euler equations
do not blow up, then the strong norm Hs of the difference between the exact
Euler orbit and the resonant orbit is uniformly small on 0 ≤ t < Tm, provided
that the vorticity of the I.C. is large enough. Paradoxically, the larger the
vortex streching of the I.C., the better the uniform approximation. This deep
result is based on cancellation of fast oscillations in strong norms, in the
context of almost periodic functions of time with values in Banach spaces
(Section 4 of [M-N-B-G]). It includes uniform approximation in the spaces
Hs, s > 5/2. For instance, given a quasiperiodic orbit on some time torus Tl
for the resonant Euler systems, the exact solutions to the Euler equations will
remain ǫ-close to the resonant quasiperiodic orbit on a time interval 0 ≤ t ≤
maxTi, 1 ≤ i ≤ l, Ti elementary periods, for large enough initial vorticity. If
orbits of the resonant Euler systems admit bursting dynamics in the strong
norms Hs, s ≥ 7/2, so do some exact solutions of the full 3D Euler equations,
for properly chosen parametrically resonant cylinders.
2 Vorticity waves and resonances of elemen-
tary swirling flows
We study initial value problem for the three-dimensional Euler equations with
initial data characterized by uniformly large vorticity:
∂tV+ (V · ∇)V = −∇p, ∇ ·V = 0, (2.1)
V(t, y)|t=0 = V(0) = Ṽ0(y) +
e3 × y (2.2)
where y = (y1, y2, y3), V(t, y) = (V1, V2, V3) is the velocity field and p is the
pressure. In Eqs. (1.1) e3 denotes the vertical unit vector and Ω is a constant
parameter. The field Ṽ0(y) depends on three variables y1, y2 and y3. Since
curl(Ω
e3 × y) = Ωe3, the vorticity vector at initial time t = 0 is
curlV(0, y) = curlṼ0(y) + Ωe3, (2.3)
and the initial vorticity has a large component weakly aligned along e3, when
Ω >> 1. These are fully three-dimensional large initial data with large initial
3D vortex stretching. We denote by Hsσ the usual Sobolev space of solenoidal
vector fields.
The base flow
Vs(y) =
e3 × y, curlVs(y) = Ωe3 (2.4)
is called a steady swirling flow and is a steady state solution (1.1)-(1.4), as
curl(Ωe3×Vs(y)) = 0. In (2.2) and (2.3), we consider I.C. which are an arbi-
trary (not small) perturbation of the base swirling flow Vs(y) and introduce
V(t, y) =
e3 × y + Ṽ(t, y), (2.5)
curlV(t, y) = Ωe3 + curlṼ(t, y), (2.6)
∂tṼ + curlṼ× Ṽ +Ωe3 × Ṽ + curlṼ×Vs(y) +∇p′ = 0, ∇ · Ṽ = 0, (2.7)
Ṽ(t, y)|t=0 = Ṽ0(y). (2.8)
Eqs. (2.1) and (2.7) are studied in cylindrical domains
C = {(y1, y2, y3) ∈ R3 : 0 < y3 < 2π/α, y21 + y22 < R2} (2.9)
where α and R are positive real numbers. If h is the height of the cylinder,
α = 2π/h. Let
Γ = {(y1, y2, y3) ∈ R3 : 0 < y3 < 2π/α, y21 + y22 = R2}. (2.10)
Without loss of generality, we can assume that R = 1. Eqs. (2.1) are consid-
ered with periodic boundary conditions in y3
V(y1, y2, y3) = V(y1, y2, y3 + 2π/α) (2.11)
and vanishing normal component of velocity on Γ
V ·N = Ṽ ·N = 0 on Γ; (2.12)
where N is the normal vector to Γ. From the invariance of 3D Euler equations
under the symmetry y3 → −y3, V1 → V1, V2 → V2, V3 → −V3, all results in
this article extend to cylindrical domains bounded by two horizontal plates.
Then the boundary conditions in the vertical direction are zero flux on the
vertical boundaries (zero vertical velocity on the plates). One only needs to
restrict vector fields to be even in y3 for V1, V2 and odd in y3 for V3, and
double the cylindrical domain to −h ≤ y3 ≤ +h.
We choose Ṽ0(y) in H
s(C), s > 5/2. In [M-N-B-G], for the case of “non-
resonant cylinders”, that is, non-resonant α = 2π/h, we have established
regularity for arbitrarily large finite times for the 3D Euler solutions for Ω
large, but finite. Our solutions are not close in any sense to those of the 2D
or “quasi 2D” Euler and they are characterized by fast oscillations in the e3
direction, together with a large vortex stretching term
ω(t, y) · ∇V(t, y) = ω1
, t ≥ 0
with leading component
V(t, y)
≫ 1. There are no assumptions on
oscillations in y1, y2 for our solutions (nor for the initial condition Ṽ0(y)).
Our approach is entirely based on sturying fast singular oscillating limits
of Eqs. (1.1)-(1.5a), nonlinear averaging and cancelation of oscillations in the
nonlinear interactions for the vorticity field for large Ω. This has been devel-
oped in [B-M-N2], [B-M-N3], and [B-M-N4] for the cases of periodic lattice
domains and the infinite space R3.
It is well known that fully three-dimensional initial conditions with uni-
formly large vorticity excite fast Poincaré vorticity waves [B-M-N2], [B-M-N3],
[B-M-N4], [Poi]. Since individual Poincaré wave modes are related to the
eigenfunctions of the curl operator, they are exact time-dependent solutions
of the full nonlinear 3D Euler equations. Of course, their linear superposition
does not preserve this property. Expanding solutions of (2.1)-(2.8) along such
vorticity waves demonstrates potential nonlinear resonances of such waves.
First recall spectral properties of the curl operator in bounded, connected
domains:
Proposition 2.1 ([M-N-B-G]) The curl operator admits a self-adjoint ex-
tension under the zero flux boundary conditions, with a discrete real spectrum
λn = ±|λn|, |λn| > 0 for every n and |λn| → +∞ as |n| → ∞. The corre-
sponding eigenfunctions Φ±n
curlΦ±n = ±|λn|Φ
n (2.13)
are complete in the space
U ∈ L2(D) : ∇ ·U = 0 and U ·N|∂D = 0 and
U dz = 0
(2.14)
Remark 2.2 In cylindrical domains, with cylindrical coordinates (r, θ, z), the
eigenfunctions admit the representation:
Φn1,n2,n3 = (Φr,n1,n2,n3(r),Φθ,n1,n2,n3(r),Φz,n1,n2,n3(r)) e
in2θeiαn3z, (2.15)
with n2 = 0,±1,±2, ..., n3 = ±1,±2, ... and n1 = 0, 1, 2, .... Here n1 indexes
the eigenvalues of the equivalent Sturm-Liouville problem in the radial coor-
dinates, and n = (n1, n2, n3). See [M-N-B-G] for technical details. From now
on, we use the generic variable z for any vertical (axial) coordinate y3 or x3.
For n3 = 0 (vertical averaging along the axis of the cylinder), 2-Dimensional,
3-component solenoidal fields must be expanded along a complete basis for
fields derived from 2D stream functions:
curl(φne3), φne3
, φn = φn(r, θ),
−△φn = µnφn, φn|∂Γ = 0, and
curlΦn =
curl(φne3), µnφne3
a, be3
denotes a 3-component vector whose horizontal projection is
a and vertical projection is be3.
Let us explicit elementary swirling wave flows which are exact solutions to
(2.1) and (2.7):
Lemma 2.3 For every n = (n1, n2, n3), the following quasiperiodic (T
time) solenoidal fields are exact solution of the full 3D nonlinear Euler equa-
tions (2.1):
V(t, y) =
e3 × y + exp(
Jt)Φn(exp(−
Jt)y) exp(±i
αΩt), (2.16)
n3 is the vertical wave number of Φn and exp(
Jt) the unitary group of rigid
body rotations:
0 −1 0
1 0 0
0 0 0
, eΩJt/2 =
cos(Ωt
) − sin(Ωt
sin(Ωt
) cos(Ωt
0 0 1
. (2.17)
Remark 2.4 These fields are exact quasiperiodic, nonaxisymmetric swirling
flow solutions of the 3D Euler equations. For n3 6= 0, their second components
Ṽ(t, y) = exp(
Jt)Φn(exp(−
Jt)y) exp(± in3
αΩt) (2.18)
are Beltrami flows (curlṼ×Ṽ ≡ 0) exact solutions of (2.7) with Ṽ(t = 0, y) =
Φn(y).
Ṽ(t, y) in Eq. (2.18) are dispersive waves with frequencies Ω
and n3α|λn|Ω,
where α = 2π
. Moreover, each Ṽ(t, y) is a traveling wave along the cylinder’s
axis, since it contains the factor
iαn3(±z ±
Note that n3 large corresponds to small axial (vertical) scales, albeit 0 ≤
α|n3/λn| ≤ 1.
Proof of Lemma 2.3. Through the canonical rigid body transformation for
both the field V(t, y) and the space coordinates y = (y1, y2, y3):
V(t, y) = e+ΩJt/2U(t, e−ΩJt/2y) +
Jy, x = e−ΩJt/2y, (2.19)
the 3D Euler equations (2.1), (2.2) transform into:
∂tU+ (curlU+Ωe3)×U = −∇
(|x1|2 + |x2|2) +
, (2.20)
∇ ·U = 0, U(t, x)|t=0 = U(0) = Ṽ0(x), (2.21)
For Beltrami flows such that curlU×U ≡ 0, these Euler equations (2.20)-
(2.21) in a rotating frame reduce to:
∂tU+Ωe3 ×U+∇π = 0, ∇ ·U = 0,
which are identical to the Poincaré-Sobolev nonlocal wave equations in the
cylinder [M-N-B-G], [Poi], [Sob], [Ar-Khe]:
∂tΨ+Ωe3 ×Ψ+∇π = 0, ∇ ·Ψ = 0, (2.22)
curl2Ψ−Ω2 ∂
Ψ = 0, Ψ ·N|∂D = 0. (2.23)
It suffices to verify that the Beltrami flows Ψn(t, x) = Φn(x) exp
±iαn3|λn|Ωt
where Φ±n (x) and ±|λn| are curl eigenfunctions and eigenvalues, are exact
solutions to the Poincaré-Sobolev wave equation, in such a rotating frame of
reference.
Remark 2.5 The frequency spectrum of the Poincaré vorticity waves (solu-
tions to (2.22)) is exactly ±iαn3|λn|Ω, n = (n1, n2, n3) indexing the spectrum
of curl. Note that n3 = 0 (zero frequency of rotating waves) corresponds to
2-Dimensional, 3-Components solenoidal vector fields.
We now transform the Cauchy problem for the 3D Euler equations (2.1)-
(2.2) into an infinite dimensional nonlinear dynamical system by expanding
V(t, y) along the swirling wave flows (2.16)-(2.18):
V(t, y) =
e3 × y
(2.24a)
+ exp
n=(n1,n2,n3)
un(t) exp
(2.24b)
V(t = 0, y) =
e3 × y + Ṽ0(y)
(2.24c)
Ṽ0(y) =
n=(n1,n2,n3)
un(0)Φn(y),
(2.24d)
where Φn denotes the curl eigenfunctions of Proposition 2.1 if n3 6= 0, and
curl(φne3), φne3
if n3 = 0 (2D case, Remark 2.2).
As we focus on the case where helicity is conserved for (2.1)-(2.2), we
consider the class of initial data Ṽ0 such that [M-N-B-G]:
curlṼ0 ·N = 0 on Γ,
where Γ is the lateral boundary of the cylinder.
The infinite dimensional dynamical system is then equivalent to the 3D
Euler equations (2.1)-(2.2) in the cylinder, with n = (n1, n2, n3) ranging over
the whole spectrum of curl, e.g.:
k3+m3=n3
k2+m2=n2
× < curlΦk ×Φm,Φn > uk(t)um(t)
(2.25)
curlΦ±k = ±λkΦ
k if k3 6= 0,
curlΦk =
curl(φke3), µkφke3
if k3 = 0
(2D, 3-components, Remark 2.2), similarly for m3 = 0 and n3 = 0. The inner
product < , > denotes the L2 complex-valued inner product in D.
This is an infinite dimensional system of coupled equations with quadratic
nonlinearities, which conserve both the energy
E(t) =
|un(t)|2
and the helicity
Hel(t) =
±|λn| |u±n (t)|2.
The quadratic nonlinearities split into resonant terms where the exponential
oscillating phase factor in (2.25) reduces to unity and fast oscillating non-
resonant terms (Ω >> 1). The resonant set K is defined in terms of vertical
wavenumbers k3,m3, n3 and eigenvalues ±λk, ±λm, ±λn of curl:
K = {± k3
= 0, n3 = k3 +m3, n2 = k2 +m2}. (2.27)
Here k2,m2, n2 are azimuthal wavenumbers.
We shall call the “resonant Euler equations” the following ∞-dimensional
dynamical system restricted to (k,m, n) ∈ K:
(k,m,n)∈K
< curlΦk ×Φm,Φn > ukum = 0, (2.28a)
un(0) ≡< Ṽ0,Φn >, (2.28b)
here curlΦ±k = ±λkΦ
k if k3 6= 0, curlΦk =
curl(φke3), µkφke3
if k3 = 0;
similarly for m3 = 0 and n3 = 0 (2D components, Remark 2.2). If there
are no terms in (2.28a) satisfying the resonance conditions, then there will be
some modes for which
Lemma 2.6 The resonant 3D Euler equations (2.28) conserve both energy
E(t) and helicity Hel(t). The energy and helicity are identical to that of the
full exact 3D Euler equations (2.1)-(2.2).
The set of resonances K is studied in depth in [M-N-B-G]. To summarize,
K splits into:
(i ) 0-wave resonances, with n3 = k3 = m3 = 0; the corresponding reso-
nant equations are identical to the 2-Dimensional, 3-Components Euler
equations, with I.C.
Ṽ0(y1, y2, y3) dy3.
(ii) Two-Wave resonances, with k3m3n3 = 0, but two of them are not null;
the corresponding resonant equations (called “catalytic equations”) are
proven to possess an infinite, countable set of new conservation laws
[M-N-B-G].
(iii) Strict three-wave resonances for a subset K∗ ⊂ K.
Definition 2.7 The set K∗ of strict 3 wave resonances is:
= 0, k3m3n3 6= 0, n3 = k3 +m3, n2 = k2 +m2
(2.29)
Note that K∗ is parameterized by h/R, since α = 2π
parameterizes the eigen-
values λn, λk, λm of the curl operator.
Proposition 2.8 There exist a countable, non-empty set of parameters h
which K∗ 6= ∅.
Proof. The technical details, together with a more precise statement, are
postponed to the proof of Lemma 3.7. Concrete examples of resonant ax-
isymmetric and helical waves are discussed in [Mah] ( cf. Figure 2 in the
article).
Corollary 2.9 Let
Ṽ0(y1, y2, y3) dy3 = 0, i.e. zero vertical mean for the
I.C. Ṽ0(y) in (2.2), (2.8), (2.24d) and (2.28b). Then the resonant 3D Euler
equations are invariant on K∗:
(k,m,n)∈K∗
λk < Φk ×Φm,Φn > ukum = 0, k3m3n3 6= 0, (2.30a)
un(0) =< Ṽ0,Φn > (2.30b)
(where Ṽ0 has spectrum restricted to n3 6= 0).
Proof. This is an immediate corollary of the “operator splitting” Theorem
3.2 in [M-N-B-G]. �
We shall call the above dynamical systems the “strictly resonant Euler
system”. This is an ∞-dimensional Riccati system which conserves Energy
and Helicity. It corresponds to nonlinear interactions depleted on K∗.
How do dynamics of the resonant Euler equations (2.28) or (2.30) approx-
imate exact solutions of the Cauchy problem for the full Euler equations in
strong norms? This is answered by the following theorem, proven in Section
4 of [M-N-B-G]:
Theorem 2.10 Consider the initial value problem
V(t = 0, y) =
e3 × y + Ṽ0(y), Ṽ0 ∈ Hsσ, s > 7/2
for the full 3D Euler equations, with ||Ṽ0||Hs
≤M0s and curlṼ0 ·N = 0 on Γ.
• Let V(t, y) = Ω
e3 × y + Ṽ(t, y) denote the solution to the exact Euler
equations.
• Let w(t, x) denote the solution to the resonant 3D Euler equations with
Initial Condition w(0, x) ≡ w(0, y) = Ṽ0(y).
• Let ||w(t, y)||Hsσ ≤Ms(TM ,M
s ) on 0 ≤ t ≤ TM , s > 7/2.
Then, ∀ǫ > 0, ∃ Ω∗(TM ,M0s , ǫ) such that, ∀Ω ≥ Ω∗:
Ṽ(t, y)− exp
un(t)e
−i n3
on 0 ≤ t ≤ TM , ∀β ≥ 1, β ≤ s− 2. Here || · ||Hβ is defined in (1.13).
The 3D Euler flow preserves the condition curlṼ0 · N = 0 on Γ, that is
curlV(t, y) · N = 0 on Γ, for every t ≥ 0 [M-N-B-G]. The proof of this
“error-shadowing” theorem is delicate, beyond the usual Gronwall differential
inequalities and involves estimates of oscillating integrals of almost periodic
functions of time with values in Banach spaces. Its importance lies in that
solutions of the resonant Euler equations (2.28) and/or (2.30) are uniformly
close in strong norms to those of the exact Euler equations (2.1)-(2.2), on
any time interval of existence of smooth solutions of the resonant system.
The infinite dimensional Riccati systems (2.28) and (2.30) are not just hydro-
dynamic models, but exact asymptotic limit systems for Ω ≫ 1. This is in
contrast to all previous literature on conservative 3D hydrodynamic models,
such as in [G-D-O].
3 Strictly resonant Euler systems: the SO(3)
We investigate the structure and the dynamics of the “strictly resonant Euler
systems” (2.30). Recall that the set of 3-wave resonances is:
(k,m, n) : ± k3
= 0, k3m3n3 6= 0,
n3 = k3 +m3, n2 = k2 +m2
(3.1)
From the symmetries of the curl eigenfunctions Φn and eigenvalues λn in the
cylinder, the following identities hold under the transformation n2 → −n2,
n3 → −n3
Φ(n1,−n2,−n3) = Φ∗(n1, n2, n3) ,
λ(n1,−n2,−n3) = λ(n1, n2, n3) .
(3.2)
where ∗ designates the complex conjugate (see Section 3, [M-N-B-G] for
details). The eigenfunctions Φ(n1, n2, n3) involve the radial functions
Jn2(β(n1, n2, αn3)r) and J
(β(n1, n2, αn3)r), with
λ2(n1, n2, n3) = β
2(n1, n2, αn3) + α
2n23;
β(n1, n2, αn3) are discrete, countable roots of equation (3.30) in [M-N-B-G],
obtained via an equivalent Sturm-Liouville radial problem. Since the curl
eigenfunctions are even in r → −r, n1 → −n1, we will extend the indices
n1 = 1, 2, ...,+∞ to −n1 = −1,−2, ... with the above radial symmetry in
mind.
Corollary 3.1 The 3-wave resonance set K∗ is invariant under the symme-
tries σj , j = 0, 1, 2, 3, where
σ0(n1, n2, n3) = (n1, n2, n3),
σ1(n1, n2, n3) = (−n1, n2, n3),
σ2(n1, n2, n3) = (n1,−n2, n3)
σ3(n1, n2, n3) = (n1, n2,−n3) .
Remark 3.2 For 0 < i ≤ 3, 0 < j ≤ 3, 0 < l ≤ 3 σ2j = Id, σiσj = −σl if
i 6= j and σiσjσl = −Id, for i 6= j 6= l. The σj do preserve the convolution
conditions in K∗.
We choose an α for which the set K∗ is not empty. We further take the
hypothesis of a single triple wave resonance (k,m, n), modulo the symmetries
Hypothesis 3.3 K∗ is such that there exists a single triple wave number
resonance (n, k,m), modulo the symmetries σj , j = 1, 2, 3 and σj(k) 6=
k, σj(m) 6= m, σj(n) 6= n for j = 2 and j = 3.
Under the above hypothesis, one can demonstrate that the strictly resonant
Euler system splits into three uncoupled systems in C3:
Theorem 3.4 Under hypothesis 3.3, the resonant Euler system reduces to
three uncoupled rigid body systems in C3:
+ i(λk − λm)CkmnUkUm = 0 (3.3a)
− i(λm − λn)CkmnUnU∗m = 0 (3.3b)
− i(λn − λk)CkmnUnU∗k = 0 (3.3c)
where Ckmn = i < Φk ×Φm,Φ∗n >, Ckmn real and the other two uncoupled
systems obtained with the symmetries σ2(k,m, n) and σ3(k,m, n). The energy
and the helicity of each subsystem are conserved:
k + UmU
m + UnU
n) = 0,
(λkUkU
k + λmUmU
m + λnUnU
n) = 0.
Proof. It follows from U−k = U
k , λ(−k) = λ(+k), similarly for m and
n; and in a very essential way from the antisymmetry of < Φk ×Φm,Φ∗n >,
together with curlΦk = λkΦk. That Ckmn is real follows from the eigenfunc-
tions explicited in Section 3 of [M-N-B-G]. �
Remark 3.5 This deep structure, i.e. SO(3;C) rigid body systems in C3 is
a direct consequence of the Lamé form of the full 3D Euler equations, cf. Eqs.
(1.3) and (2.7), and the nonlinearity curlV ×V.
The system (3.3) is equivariant with respect to the symmetry operators
(z1, z2, z3) → (z∗1 , z∗2 , z∗3), (z1, z2, z3) → (exp(iχ1)z1, exp(iχ2)z2, exp(iχ3)z3) ,
provided χ1 = χ2 + χ3. It admits other integrals known as the Manley-
Rowe relations (see, for instance [We-Wil]). It differs from the usual 3-
wave resonance systems investigated in the literature, such as in [Zak-Man1],
[Zak-Man2], [Gu-Ma] in that
(1) helicity is conserved,
(2) dynamics of these resonant systems rigorously “shadow” those of the
exact 3D Euler equations, see Theorem 2.10.
Real forms of the system (3.3) are found in Gledzer et al. [G-D-O], corre-
sponding to the exact invariant manifold Uk ∈ iR, Um ∈ R, Un ∈ R, albeit
without any rigorous asymptotic justification. The C3 systems (3.3) with
helicity conservation laws are not discussed in [G-D-O].
The only nontrivial Manley-Rowe conservation laws for the resonant sys-
tem (3.3), rigid body SO(3;C), which are independent from energy and he-
licity, are:
(rkrmrn sin(θn − θk − θm)) = 0,
where Uj = rj exp(iθj), j = k,m, n, and
E1 = (λk − λm)r2n − (λm − λn)r2k,
E2 = (λm − λn)r2k − (λn − λk)r2m.
The resonant system (3.3) is well known to possess hyperbolic equilibria
and heteroclinic/homoclinic orbits on the energy surface. We are interested
in rigorously proving arbitrary large bursts of enstrophy and higher norms
on arbitrarily small time intervals, for properly chosen h/R. To simplify
the presentation, we establish the results for the simpler invariant manifold
Uk ∈ iR, and Um, Un ∈ R.
Rescale time as:
t→ t/Ckmn.
Start from the system
U̇n + i(λk − λm)UkUm = 0
U̇k − i(λm − λn)UnU∗m = 0
U̇m − i(λn − λk)UnU∗k = 0
(3.4)
Assume that Uk ∈ iR and that Um, Un ∈ R: set p = iUk, q = Um and r = Un,
as well as λk = λ, λm = µ and λn = ν: then
ṗ+ (µ− ν)qr = 0
q̇ + (ν − λ)rp = 0
ṙ + (λ− µ)pq = 0
(3.5)
This system admits two first integrals:
E = p2 + q2 + r2 (energy)
H = λp2 + µq2 + νr2 (helicity)
(3.6)
System (3.5) is exactly the SO(3,R) rigid body dynamics Euler equations,
with inertia momenta Ij =
|λj | , j = k,m, n [Ar1].
Lemma 3.6 ([Ar1], [G-D-O]) With the ordering λk > λm > λn, i.e. λ > µ >
ν, the equilibria (0,±1, 0) are hyperbolic saddles on the unit energy sphere,
and the equilibria (±1, 0, 0), (0, 0,±1) are centers. There exist equivariant
families of heteroclinic connections between (0,+1, 0) and (0,−1, 0). Each
pair of such connections correspond to equivariant homoclinic cycles at (0, 1, 0)
and (0,−1, 0).
We investigate bursting dynamics along orbits with large periods, with
initial conditions close to the hyperbolic point (0, E(0), 0) on the energy sphere
E. We choose resonant triads such that λk > 0, λn < 0, λk ∼ |λn|, |λm| ≪ λk,
equivalently:
λ > µ > ν, λν < 0, |µ| ≪ λ and λ ∼ |ν|. (3.7)
Lemma 3.7 There exist h/R with K∗ 6= ∅, such that
λk > λm > λn, λkλn < 0, |λm| ≪ λk and λk ∼ |λn|.
Remark 3.8 Together with the polarity ± of the curl eigenvalues, these are
3-wave resonances where two of the eigenvalues are much larger in mod-
uli than the third one. In the limit |k|, |m|, |n| ≫ 1, λk ∼ ±|k|, λm ∼
±|m|, λn ∼ ±|n|, the eigenfunctions Φ have leading asymptotic terms which
involve cosines and sines periodic in r, cf. Section 3 [M-N-B-G]. In the
strictly resonant equations (2.30), the summation over the quadratic terms
becomes an asymptotic convolution in n1 = k1+n1. The resonant three waves
in Lemma 3.7 are equivalent to Fourier triads k + m = n, with |k| ∼ |n|
and |m| ≪ |k|, |n|, in periodic lattices. In the physics of spectral theory of
turbulence [Fri], [Les], these are exactly the triads responsible from transfer of
energy between large scales and small scales. These are the triads which have
hampered mathematical efforts at proving the global regularity of the Cauchy
problem for 3D Navier-Stokes equations in periodic lattices [Fe].
Proof of Lemma 3.7 ([M-N-B-G]) The transcendental dispersion law for
3-waves in K∗ for cylindrical domains, is a polynomial of degree four in ϑ3 =
1/h2:
P̃ (ϑ3) = P̃4ϑ
3 + P̃3ϑ
3 + P̃2ϑ
3 + P̃1ϑ3 + P̃0 = 0, (3.8)
with n2 = k2 +m2 and n3 = k3 +m3.
Then with hk =
β2(k1,k2,αk3)
, hm =
β2(m1,m2,αm3)
, hn =
β2(n1,n2,αn3)
, cf.
the radial Sturm-Liouville problem in Section 3, [M-N-B-G], the coefficients
of P̃ (ϑ3) are given by:
P̃4 = −3,
P̃3 = −4(hk + hm + hn),
P̃2 = −6(hkhm + hkhn + hmhn),
P̃1 = −12hkhmhn,
P̃0 = h
n + h
n + h
k − 2(hkhmh2n + hkhnh2m + hmhnh2k).
Similar formulas for the periodic lattice domain were first derived in [B-M-N2],
[B-M-N3], [B-M-N4]. In cylindrical domains the resonance condition for K∗
is identical to
ϑ3 + hk
ϑ3 + hm
ϑ3 + hn
with ϑ3 =
, hk = β
2(k)/k23 , hm = β
2(m)/m23, hn = β
2(n)/n23; Eq. (3.8) is
the equivalent rational form.
From the asymptotic formula (3.44) in [M-N-B-G], for large β:
β(n1, n2, n3) ∼ n1π + n2
+ ψ, (3.9)
where ψ = 0 if lim m2
= 0 (e.g. h fixed, m2/m3 → 0) and ψ = ±π2 if
lim m2
= ±∞ (e.g. m2
fixed, h → ∞). The proof is completed by taking
leading terms P̃0+ϑ3P̃1 in (3.8), ϑ3 =
≪ 1, and m2 = 0, k2 = O(1), n2 =
O(1). �
We now state a theorem for bursting of the H3 norm in arbitrarily small
times, for initial data close to the hyperbolic point (0, E(0), 0):
Theorem 3.9 (Bursting dynamics in H3). Let λ > µ > ν, λν < 0, |µ| ≪ λ
and λ ∼ |ν|. Let W (t) = λ6p(t)2 + µ6q(t)2 + ν6r(t)2 the H3-norm squared of
an orbit of (3.5). Choose initial data such that: W (0) = λ6p(0)2 + µ6q(0)2
with λ6p(0)2 ∼ 1
W (0) and µ6q(0)2 ∼ 1
W (0). Then there exists t∗ > 0, such
W (t) ≥
W (0)
where t∗ ≤ 6√
W (0)
µ2Ln(λ/|µ|)(λ/|µ|)−1.
Remark 3.10 Under the conditions of Lemma 3.7,
≫ 1, whereas
µ2(Ln(λ/|µ|))(λ/|µ|)−1 ≪ 1. Therefore, over a small time interval of length
O(µ2(Ln(λ/|µ|))(λ/|µ|)−1) ≪ 1, the ratio ||U(t)||H3/||U(0)||H3 grows up to
a maximal value O
(λ/|µ|)3
≫ 1. Since the orbit is periodic, the H3 semi-
norm eventually relaxes to its initial state after some time (this being a mani-
festation of the time-reversibility of the Euler flow on the energy sphere). The
“shadowing” theorem 2.10 with s > 7/2 ensures that the full, original 3D Eu-
ler dynamics, with the same initial conditions, will undergo the same type of
burst. Notice that, with the definition (1.13) of ‖ · ‖Hs , one has
||Ωe3 × y||H3 = ||curl3(Ωe3 × y)||L2 = 0 .
Hence the solid rotation part of the original 3D Euler solution does not con-
tribute to the ratio ||V(t)||H3/||V(0)||H3 .
Theorem 3.11 (Bursting dynamics of the enstrophy). Under the same con-
ditions for the 3-wave resonance, let Ξ(t) = λ2p(t)2 + µ2q(t)2 + ν2r(t)2 the
enstrophy. Choose initial data such that Ξ(0) = λ2p(0)2 + µ2q(0)2 + ν2r(0)2
with λ2p(0)2 ∼ 1
Ξ(0), µ2q(0)2 ∼ 1
Ξ(0). Then there exists t∗∗ > 0, such that
Ξ(t∗∗) ≥
where t∗∗ ≤ 1√
Ln (λ/|µ|) (λ/|µ|)−1 .
Remark 3.12 It is interesting to compare this mechanism for bursts with ear-
lier results in the same direction obtained by DiPerna and Lions. Indeed, for
each p ∈ (1,∞), each δ ∈ (0, 1) and each t > 0, Di Perna and Lions [DiPe-Li]
constructed examples of 2D-3 components solutions to Euler equations such
||V(0)||W 1,p ≤ ǫ while ||V(t)||W 1,p ≥ 1/δ .
Their examples essentially correspond to shear flows of the form
V(t, x1, x2) =
u(x2)
w(x1 − tu(x2), x2)
where u ∈W 1,px2 while w ∈W
. Obviously
curlV(t, x1, x2) =
(∂2 − tu′(x2)∂1)w(x1 − tu(x2), x2)
−∂1w(x1 − tu(x2), x2)
−u′(x2)
Thus, all components in curlV(t, x1, x2) belong to L
loc, except for the term
−tu′(x2)∂1w(x1 − tu(x2), x2) .
For each t > 0, this term belongs to Lp for all choices of the functions u ∈
W 1,px2 and w ∈ W
x1,x2
if and only if p = ∞. Whenever p < ∞, DiPerna and
Lions construct their examples as some smooth approximation of the situation
above in the strong W 1,p topology.
In other words, the DiPerna-Lions construction works only in cases where
the initial vorticity does not belong to an algebra — specifically to Lp, which
is not an algebra unless p = ∞.
The type of burst obtained in our construction above is different: in that
case, the original vorticity belongs to the Sobolev space H2, which is an algebra
in space dimension 3. Similar phenomena are observed in all Sobolev spaces
Hβ with β ≥ 2 — which are also algebras in space dimension 3.
In other words, our results complement those of DiPerna-Lions on bursts
in higher order Sobolev spaces, however at the expense of using more intricate
dynamics.
We proceed to the proofs of Theorem 3.9 and 3.11. We are interested in
the evolution of
Ξ = λ2p2 + µ2q2 + ν2r2 (enstrophy) (3.10)
Compute
Ξ̇ = −2
λ2(µ− ν) + µ2(ν − λ) + ν2(λ− µ)
pqr (3.11)
˙(pqr) = −(µ− ν)q2r2 − (ν − λ)r2p2 − (λ− µ)p2q2 (3.12)
Using the first integrals above, one has
(3.13)
where V an is the Vandermonde matrix
V an =
1 1 1
λ µ ν
λ2 µ2 ν2
For λ 6= µ 6= ν 6= λ, this matrix is invertible and
V an−1 =
(λ−µ)(λ−ν)
−(µ+ν)
(λ−µ)(λ−ν)
(λ−µ)(λ−ν)
(µ−ν)(µ−λ)
−(ν+λ)
(µ−ν)(µ−λ)
(µ−ν)(µ−λ)
(ν−λ)(ν−µ)
−(λ+µ)
(ν−λ)(ν−µ)
(ν−λ)(ν−µ)
Hence
(λ− µ)(λ − ν)
(Ξ− (µ+ ν)H + µνE)
(µ− ν)(µ− λ)
(Ξ− (ν + λ)H + νλE)
(ν − λ)(ν − µ)
(Ξ− (λ+ µ)H + λµE)
(3.14)
so that
(µ− ν)q2r2 = − (Ξ− (ν + λ)H + νλE) (Ξ− (λ + µ)H + λµE)
(λ− µ)(λ − ν)(µ− ν)
(ν − λ)r2p2 = − (Ξ− (λ + µ)H + λµE) (Ξ− (µ+ ν)H + µνE)
(λ− µ)(λ − ν)(µ− ν)
(λ− µ)p2q2 = − (Ξ− (µ+ ν)H + µνE) (Ξ− (ν + λ)H + νλE)
(λ− µ)(λ − ν)(µ− ν)
Later on, we shall use the notations
x−(λ, µ, ν) = (µ+ ν)H − µνE
x0 (λ, µ, ν) = (µ+ λ)H − µλE
x+(λ, µ, ν) = (λ+ ν)H − λνE
(3.15)
Therefore, we find that Ξ satisfies the second order ODE
Ξ̈ =− 2Kλ,µ,ν ((Ξ− x−(λ, µ, ν))(Ξ − x0(λ, µ, ν))
+(Ξ− x0(λ, µ, ν))(Ξ − x+(λ, µ, ν)) + (Ξ− x+(λ, µ, ν))(Ξ − x0(λ, µ, ν)))
which can be put in the form
Ξ̈ = −2Kλ,µ,νP ′λ,µ,ν(Ξ) (3.16)
where Pλ,µ,ν is the cubic
Pλ,µ,ν(X) = (X − x−(λ, µ, ν))(X − x0(λ, µ, ν))(X − x+(λ, µ, ν)) (3.17)
Kλ,µ,ν =
λ2(µ− ν) + µ2(ν − λ) + ν2(λ− µ)
(λ− µ)(λ− ν)(µ − ν)
(3.18)
In the sequel, we assume that the initial data for (p, q, r) is such that
r(0) = 0 , p(0)(q(0) 6= 0
Let us compute
x−(λ, µ, ν) = λνp(0)
2 + µ2q(0)2 + µ(λ− ν)p(0)2
x0 (λ, µ, ν) = λ
2p(0)2 + µ2q(0)2
x+(λ, µ, ν) = λ
2p(0)2 +
ν + λ
µ2q(0)2
(3.19)
We shall also assume that
λ > µ > ν , λν < 0 , |µ| ≪ λ and λ ∼ |ν| (3.20)
Then Kλ,µ,ν > 0 — in fact Kλ,µ,ν ∼ 2, and Ξ is a periodic function of t such
Ξ(t) = x0(λ, µ, ν) , sup
Ξ(t) = x+(λ, µ, ν) (3.21)
with half-period
Tλ,µ,ν =
Kλ,µ,ν
∫ x+(λ,µ,ν)
x0(λ,µ,ν)
−Pλ,µ,ν(x)
(3.22)
We are interested in the growth of the (squared) H3 norm
W (t) = λ6p(t)2 + µ6q(t)2 + ν6r(t)2 (3.23)
Expressing p2, q2 and r2 in terms of E, H and Ξ, it is found that
λ6(Ξ− x−(λ, µ, ν))
(λ− µ)(λ − ν)
µ6(Ξ− x+(λ, µ, ν))
(µ− ν)(µ− λ)
ν6(Ξ− x0(λ, µ, ν))
(ν − λ)(ν − µ)
(3.24)
Hence, when Ξ = x+(λ, µ, ν), then
λ6(x+(λ, µ ν) − x−(λ, µ, ν))
(λ− µ)(λ− ν)
ν6(x+(λ, µ ν)− x0(λ, µ, ν))
(ν − λ)(ν − µ)
λ6(x+(λ, µ ν) − x−(λ, µ, ν))
(λ− µ)(λ− ν)
Let us compute
x+(λ, µ ν)− x−(λ, µ, ν) = (λ − µ)(λ− ν)p(0)2 +
ν + λ
µ2q(0)2
& −νλq(0)2 ∼ λ2q(0)2
(3.25)
We shall pick the initial data such that
W (0) = λ6p(0)6 + µ6q(0)6 with λ6p(0)2 ∼ 1
W (0) and µ6q(0)2 ∼ 1
W (0)
(3.26)
Hence, when Ξ reaches x+(λ, µ, ν), one has
λ8q(0)2
(λ− µ)(λ− ν)
µ6(λ − µ)(λ− ν)
W (0) ∼ 1
W (0) . (3.27)
Hence W jumps from W (0) to a quantity ∼ 1
W (0) in an interval of time
that does not exceed one period of the Ξ motion, i.e. 2Tλ,µ,ν . Let us estimate
this interval of time. We recall the asymptotic equivalent for the period of an
elliptic integral in the modulus 1 limit.
Lemma 3.13 Assume that x− < x0 < x+. Then
(x− x−)(x− x0)(x+ − x)
x+ − x−
x+−x0
x+−x−
uniformly in x−, x0, and x+ as
x+−x0
x+−x− → 1.
x+(λ, µ, ν) − x−(λ, µ, ν)
λ2q(0)2
∼ |µ|
W (0)
x0(λ, µ, ν) − x−(λ, µ, ν) = (λ− µ)(λ − ν)p(0)2 (3.28)
so that
x+−x0
x+−x−
1− (λ−µ)(λ−ν)p(0)
(λ−µ)(λ−ν)p(0)2+(µ(ν+λ)−νλ−µ2)q(0)2
∼ (λ− µ)(λ− ν)p(0)
2 + (µ(ν + λ)− νλ− µ2)q(0)2
2(λ− µ)(λ − ν)p(0)2
∼ q(0)
2p(0)2
W (0)/2µ6
W (0)/2λ6
Hence
2Tλ,µ,ν .
W (0)
≤ 12√
W (0)
(3.29)
Conclusion: collecting (3.26), (3.27) and (3.29), we see that the squared H3
norm W varies from W (0) to a quantity ∼ ρ6W (0) in an interval of time
. 12√
W (0)
µ2 ln ρ
. (Here ρ = λ/µ).
We now proceed to obtain similar bursting estimates for the enstrophy.
We return to (3.21) and (3.22). Pick the initial data so that
Ξ(0) = λ2p(0)2 + µ2q(0)2 with λ2p(0)2 ∼ 1
Ξ(0) and µ2q(0)2 ∼ 1
Ξ(0).
x+(λ, µ, ν) − x−(λ, µ, ν)
= (λ− µ)(λ − ν)p(0)2 +
ν + λ
µ2q(0)2
∼ 2λ2p(0)2 + λ2q(0)2 ∼
while
x0(λ, µ, ν)− x−(λ, µ, ν) = (λ− µ)(λ− ν)p(0)2 ∼ 2λ2p(0)2 ∼ Ξ(0).
Hence, in the limit as ρ = λ/|µ| → +∞, one has
2Tλ,µ,ν ∼
ρ2Ξ(0)
1− Ξ(0)1
ρ2Ξ(0)
2Ξ(0)
1− 2ρ−2
2Ξ(0)
And Ξ varies from
x0(λ, µ, ν) = Ξ(0) to x+(λ, µ, ν) ∼ ρ2Ξ(0)
on an interval of time of length Tλ,µ,ν . �
4 Strictly resonant Euler systems: the case of
3-waves resonances on small-scales
4.1 Infinite dimensional uncoupled SO(3) systems
In this section, we consider the 3-wave resonant set K∗ when
|k|2, |m|2, |n|2 ≥
, 0 < η ≪ 1,
i.e. 3-wave resonances on small scales; here |k|2 = k21 + k22 + k23 , where
(k1, k2, k3) index the curl eigenvalues, and similarly for |m|2, |n|2. Recall
that k2 + m2 = n2, k3 + m3 = n3 (exact convolutions), but that the sum-
mation on k1, m1 on the right hand side of Eqs. (2.30) is not a convolution.
However, for |k|2, |m|2, |n|2 ≥ 1
, the summation in k1, m1 becomes an
asymptotic convolution. First:
Proposition 4.1 The set K∗ restricted to |k|2, |m|2, |n|2 ≥ 1
, ∀η, 0 < η ≪ 1
is not empty: there exist at least one h/R with resonant three waves satisfying
the above small scales condition.
Proof. We follow the algebra of the exact transcendental dispersion law
(3.8) derived in the proof of Lemma 3.7. Note that P̃ (ϑ3) < 0 for ϑ3 =
large enough. We can choose hm =
β2(m1,m2,αm3)
= 0, say in the specific limit
→ 0, and β(m1,m2, αm3) ∼ m1π+m2 π2 +
. Then P̃0 = h
n > 0 and
P̃ (ϑ3) must possess at least one (transcendental) root ϑ3 =
In the above context, the radial components of the curl eigenfunctions in-
volve cosines and sines in βr
(cf. Section 3, [M-N-B-G]) and the summation in
k1, m1 on the right hand side of the resonant Euler equations (2.30) becomes
an asymptotic convolution. The rigorous asymptotic convolution estimates are
highly technical and detailed in [Fro-M-N]. The 3-wave resonant systems for
|k|2, |m|2, |n|2 ≥ 1
are equivalent to those of an equivalent periodic lattice
[0, 2π]× [0, 2π]× [0, 2πh], ϑ3 = 1h2 ; the resonant three wave relation becomes:
ϑ3 + ϑ1
ϑ3 + ϑ1
ϑ3 + ϑ1
= 0, (4.1a)
k +m = n, k3m3n3 6= 0. (4.1b)
The algebraic geometry of these rational 3-wave resonance equations has been
investigated in depth in [B-M-N3] and [B-M-N4]. Here ϑ1, ϑ2, ϑ3 are periodic
lattice parameters; in the small-scales cylindrical case, ϑ1 = ϑ2 = 1 (after
rescaling of n2, k2, m2), ϑ3 = 1/h
2, h height. Based on the algebraic
geometry of “resonance curves” in [B-M-N3], [B-M-N4], we investigate the
resonant 3D Euler equations (2.30) in the equivalent periodic lattices.
First, triplets (k,m, n) solution of (4.1) are invariant under the reflec-
tion symmetries σ0, σ1, σ2, σ3 defined in Corollary 3.1 and Remark 3.2: σ0 =
Id, σj(k) = (ǫi,jki), 1 ≤ i ≤ 3, ǫi,j = +1 if i 6= j, ǫi,j = −1 if i = j, 1 ≤ j ≤ 3.
Second the set K∗ in (4.1) is invariant under the homothetic transformations:
(k,m, n) → (γk, γm, γn), γ rational. (4.2)
The resonant triplets lie on projective lines in the wavenumber space, with
equivariance under σj , 0 ≤ j ≤ 3 and γ-rescaling. For every given equivariant
family of such projective lines, the resonant curve is the graph of ϑ3
versus
, for parametric domain resonances in ϑ1, ϑ2, ϑ3.
Lemma 4.2 (p.17, [B-M-N4]). For every equivariant (k,m, n), the resonant
curve in the quadrant ϑ1 > 0, ϑ2 > 0, ϑ3 > 0 is the graph of a smooth function
ϑ3/ϑ1 ≡ F (ϑ2/ϑ1) intersected with the quadrant.
Theorem 4.3 (p.19, [B-M-N4]). A resonant curve in the quadrant ϑ3/ϑ1 > 0,
ϑ2/ϑ1 > 0 is called irreducible if:
k23 k
m23 m
n23 n
6= 0. (4.3)
An irreducible resonant curve is uniquely characterized by six non-negative
algebraic invariants P1, P2, R1, R2, S1, S2, such that
P21 ,P22
R21,R22
S21 ,S22
and permutations thereof.
Lemma 4.4 (p. 25, [B-M-N4]). For resonant triplets (k,m, n) associated to
a given irreducible resonant curve, that is verifying Eq. (4.3), consider the
convolution equation n = k +m. Let σi(n) 6= n, ∀i, 1 ≤ i ≤ 3. Then there
are no more that two solutions (k,m) and (m, k), for a given n, provided
the six non-degeneracy conditions (3.39)-(3.44) in [B-M-N4] for the algebraic
invariants of the irreducible curve are verified.
For more details on the technical non-degeneracy conditions, see the Ap-
pendix. An exhaustive algebraic geometric investigation of all solutions to
n = k +m on irreducible resonant curves is found in [B-M-N4]. The essence
of the above lemma lies in that given such an irreducible, “non-degenerate”
triplet (k,m, n) on K∗, all other triplets on the same irreducible resonant
curves are exhaustively given by the equivariant projective lines:
(k,m, n) → (γk, γm, γn), for some γ rational , (4.4)
(k,m, n) → (σjk, σjm,σjn), j = 1, 2, 3, (4.5)
and permutations of k and m in the above. Of course the homothety γ and
the σj symmetries preserve the convolution. This context of irreducible, “non-
degenerate” resonant curves yields an infinite dimensional, uncoupled system
of rigid body SO(3;R) and SO(3;C) dynamics for the 3D resonant Euler
equations (2.30).
Theorem 4.5 For any irreducible triplet (k,m, n) which satisfy Theorem 4.3,
and under the “non-degeneracy” conditions of Lemma 4.4 (cf. Appendix), the
resonant Euler equations split into the infinite, countable sequence of uncou-
pled SO(3;R) systems:
ȧk = Γkmn(λm − λn)aman, (4.6a)
ȧm = Γkmn(λn − λk)anak, (4.6b)
ȧn = Γkmn(λk − λm)akam, (4.6c)
for all (k,m, n) = γ(σj(k
∗), σj(m
∗), σj(n
∗)), γ = ±1,±2,±3..., 0 ≤ j ≤ 3.
(4.7)
k∗,m∗, n∗ are some relatively prime integer vectors in Z3 characterizing the
equivariant family of projective lines (k,m, n); Γkmn = i < Φk ×Φm,Φ∗n >,
Γkmn real.
Proof. Theorem 4.5 is a simpler version for invariant manifolds of more
general SO(3;C) systems. It is a straightforward corollary of Proposition 3.2,
Proposition 3.3, Theorem 3.3, Theorem 3.4 and Theorem 3.5 in [B-M-N4].
The latter article did not explicit the resonant equations and did not use the
curl-helicity algebra fundamentally underlying this present work. Rigorously
asymptotic infinite countable sequences of uncoupled SO(3;R), SO(3;C) sys-
tems are not derived via the usual harmonic analysis tools of Fourier modes,
in the 3D Euler context. Polarization of curl eigenvalues and eigenfunctions
and helicity play an essential role.
Corollary 4.6 Under the conditions λn∗ − λk∗ > 0, λk∗ − λm∗ > 0, the
resonant Euler systems (4.6) admit a disjoint, countable family of homoclinic
cycles. Moreover, under the conditions λn∗ ≫ +1, λm∗ ≪ −1, |λk∗ | ≪ λn∗ ,
each subsystem (4.6) possesses orbits whose Hs norms, s ≥ 1, burst arbitrarily
large in arbitrarily small times.
Remark 4.7 One can prove that there exists some Γmax, 0 < Γmax < ∞,
such that |Γkmn| < Γmax, for all (k,m, n) on the equivariant projective lines
defined by (4.7). Systems (4.6) “freeze” cascades of energy; their total enstro-
phy Ξ(t) =
(k,m,n)(λ
k(t) + λ
m(t) + λ
n(t)) remains bounded, albeit
with large bursts of Ξ(t)/Ξ(0), on the reversible orbits topologically close to
the homoclinic cycles.
4.2 Coupled SO(3) rigid body resonant systems
We now derive a new resonant Euler system which couples two SO(3;R)
rigid bodies via a common principle axis of inertia and a common moment
of inertia. This 5-dimensional system conserves energy, helicity, and is rather
interesting in that dynamics on its homoclinic manifolds show bursting cas-
cades of enstrophy to the smallest scale in the resonant set. We consider the
equivalent periodic lattice geometry under the conditions of Proposition 4.1.
In Appendix, we prove that for an “irreducible” 3-wave resonant set which
now satisfies the algebraic “degeneracy” (A-4), there exist exactly two “prim-
itive” resonant triplets (k,m, n) and (k̃, m̃, n), where k, m, k̃, m̃ are relative
prime integer valued vectors in Z3:
Lemma 4.8 Under the algebraic degeneracy condition (A-4) the irreducible
equivariant family of projective lines in K∗ is exactly generated by the follow-
ing two “primitive” triplets:
n = k +m, k = ak, m = bm, (4.8a)
n = k̃ + m̃, k̃ = a′σi(k) + b
′σj(m), (4.8b)
that is,
n = ak + bm, (4.8c)
n = a′σi(k) + b
′σj(m), (4.8d)
where σi 6= σj are some reflection symmetries, a, b, a′, b′ are relatively
prime integers, positive or negative, and k, m are relatively prime integer
valued vectors in Z3, that is:
(a, a′) = (b, b′) = (a, b) = (a′, b′) = 1, (k,m) = 1,
where ( , ) denotes the Greatest Common Denominator of two integers. All
other resonant wave number triplets are generated by the group actions σl,
l = 1, 2, 3, and homothetic rescalings (k,m, n) → γ(k,m, n), (k̃, m̃, n) →
γ(k̃, m̃, n), (γ ∈ Z) of the “primitive” triplets.
Remark 4.9 It can be proven that the set of such coupled “primitive” triplets
is not empty on the periodic lattice. The algebraic irreducibility condition of
Lemma 4.2 implies that ±k3/|k| = ±k̃3/|k̃| and ±m3/|m| = ±m̃3/|m̃|, which
is obviously verified in equations (4.8).
Theorem 4.10 Under conditions of Lemma 4.8 the resonant Euler system
reduces to a system of two rigid bodies coupled via an(t):
ȧk = (λm − λn)Γaman (4.9a)
ȧm = (λn − λk)Γanak (4.9b)
ȧn = (λk − λm)Γakam + (λk̃ − λm̃)Γ̃ak̃am̃ (4.9c)
ȧm̃ = (λn − λk̃)Γ̃anak̃ (4.9d)
ȧk̃ = (λm̃ − λn)Γ̃am̃an, (4.9e)
where Γ = i < Φk×Φm,Φ∗n >, Γ̃ = i < Φk̃×Φm̃,Φ∗n >. Energy and Helicity
are conserved.
Theorem 4.11 The resonant system (4.9) possesses three independent con-
servation laws:
E1 = a2k + (1− α)a2m, (4.10a)
E2 = a2n + αa2m + (1 − α̃)a2m̃, (4.10b)
E3 = a2k̃ + α̃a
m̃, (4.10c)
where
α = (λm − λk)/(λn − λk), (4.11a)
α̃ = (λm̃ − λn)/(λk̃ − λn). (4.11b)
Theorem 4.12 Under the conditions
λm < λk < λn, (4.12a)
λm̃ < λn < λk̃, (4.12b)
which imply α < 0, α̃ < 0, the equilibria (±ak(0), 0, 0, 0,±ak̃(0)) are hyper-
bolic for |ak̃(0)| small enough with respect to |ak(0)|. The unstable manifolds
of these equilibria are one dimensional, and the nonlinear dynamics of system
(4.9) are constrained on the ellipse E1 (4.10a) for ak(t), am(t), the hyperbola
E3 (4.10c) for ak̃(t), am̃(t), and the hyperboloid E2 (4.10b) for am(t), am̃(t),
an(t).
Theorem 4.13 Let the 2-manifold E1 ∩E2 ∩E3 be coordinatized by (am, am̃).
On this 2-manifold, the resonant system (4.9) is Hamiltonian, and therefore
integrable. Its Hamiltonian vector field h is defined by
ιhω = Γ(λn − λk)
− Γ̃(λn − λk̃)
, (4.13)
where ιhω designates the inner product of the symplectic 2-form
dam ∧ dam̃
akanak̃
(4.14)
with the vector field h.
Proof of Theorem 4.13: Eliminating ak(t) via E1, an(t) via E2, ak̃(t) via
E3, the resonant system (4.9) reduces to:
ȧm = ±Γ(λn − λk)(E1 − (1− α)a2m)
2 (E2 − αa2m + (α̃− 1)a2m̃)
ȧm̃ = ±Γ̃(λn − λk̃)(E2 − αa
m + (α̃− 1)a2m̃)
2 (E3 − α̃a2m̃)
after changing the time variable into
(E1 − (1− α)a2m)
2 (E2 − αa2m + (α̃− 1)a2m̃)
2 (E3 − α̃a2m̃)
2 ds .
On each component of the manifold E1 ∩E2 ∩E3, the following functionals
are conserved:
H(am, am̃) = ± Γ̃(λn − λk̃)
(E1 − (1− α)a2m)1/2
± Γ(λn − λk)
(E3 − α̃a2m̃)1/2
Observe that the system of two coupled rigid bodies (4.9) does not seem to ad-
mit a simple Lie-Poisson bracket in the original variables (ak, am, an, am̃, ak̃).
Yet, when restricted to the 2-manifold E1 ∩E2 ∩E3 that is invariant under the
flow of (4.9), it is Hamiltonian and therefore integrable.
This raises the following interesting issue: according to the shadowing
Theorem 2.10, the Euler dynamics remains asymptotically close to that of
chains of coupled SO(3;R) and SO(3;C) rigid body systems. Perhaps some
new information could be obtained in this way. We are currently investigating
this question and will report on it in a forthcoming publication [G-M-N].
Already the simple 5-dimensional system (4.9) has interesting dynamical
properties, wich we could not find in the existing literature on systems related
to spinning tops.
Consider for instance the dynamics of the resonant system (4.9) with I.C.
topologically close to the hyperbola equilibria (±ak(0), 0, 0, 0,±ak̃(0)). Un-
der the conditions of (4.12) and with the help of the integrability Theorem
4.13, it is easy to construct equivariant families of homoclinic cycles at these
hyperbolic critical points:
Corollary 4.14 The hyperbolic critical points (±ak(0), 0, 0, 0,±ak̃(0)) pos-
sess 1-dimensional homoclinic cycles on the cones
a2n + (1− α̃)a2m̃ = −αa2m (4.15)
with α < 0, α̃ < 0.
Note that these are genuine homoclinic cycles, NOT sums of heteroclinic
connections. Initial conditions for the resonant system (4.9) are now chosen
in a small neighborhood of these hyperbolic critical points, the corresponding
orbits are topologically close to these cycles. With the ordering:
λm < λk < λn, (4.16a)
|λk| ≪ |λm|, |λk| ≪ λn, (4.16b)
λm̃ < λn < λk̃, (4.16c)
|λm̃| ≪ λk̃, (4.16d)
λk̃ ≫ λn, (4.16e)
which can be realized with |a
| ≫ 1 and | b
| ≪ 1 in the resonant triplets
(4.8), we can demonstrate bursting dynamics akin to Theorem 3.9 and 3.11 for
enstrophy and Hs norms, s ≥ 2. The interesting feature is the maximization
of |ak̃(t)| near the turning points of the homoclinic cycles on the cones (4.15).
This corresponds to transfer of energy to the smallest scale k̃, λk̃.
In a publication in preparation, we investigate infinite systems of the cou-
pled rigid bodies equations (4.9).
APPENDIX
We focus on a resonant wave number triplet (n, k,m) ∈ (Z∗)3 verifying
• the convolution relation
n = k +m, (A-1)
• the resonant 3-wave resonance relation
± n3√
1 + ϑ2n
2 + ϑ3n
± k3√
1 + ϑ2k
2 + ϑ3k
± m3√
1 + ϑ2m
2 + ϑ3m
(A-2)
• the condition of “non-catalyticity”
k3m3n3 6= 0, (A-3)
• and the degeneracy condition of [B-M-N4] (see p26)
Giri,j(k,m) = kinjml + klmjni = 0, (A-4)
where (i, j, l) is a permutation of (1, 2, 3).
Then, we know (see lemma 3.5 (2) of [B-M-N4]) that the system of equations
(A-3)-(A-4) for the unknown k and m, given the vector n, admits exactly 4
solutions in Z3 × Z3:
(k,m), (m, k), (k̃, m̃), (m̃, k̃).
Here k and m are the two vectors of the original resonant triplet, whereas
k̃ = ασi(k), m̃ = βσj(m)
where
mikl −mlki
mikl +mlki
/∈ {0,±1} and β = mlkj −mjkl
mlkj +mjkl
/∈ {0,±1}
and where the symmetries σi and σj are defined by
σi : u = (ul)l=1,2,3 →
(−1)δilul
l=1,2,3
One verifies that
σ2i = σ
j = Id, σiσj = σjσi = −σl.
That is, the group generated by σi and σj is the Klein group Z/2Z× Z/2Z.
Let us first write the irrational numbers α and β under the irreducible
representation
, β =
, with a, a′, b, b′ ∈ Z∗ and (a, a′) = (b, b′) = 1,
where ( , ) denotes the Greatest Common Denominator of the integer pair.
From k̃ ∈ Z3, it follows that a|a′k; but since (a, a′) = 1, the Euclid’s lemma
yields that a|k. Similarly, b|m. Now set
k ∈ Z3, m = 1
m ∈ Z3.
Hence the integer vector n admits the two decompositions
n = ak + bm = a′σi(k) + b
′σj(m).
Since the function
z 7−→ z3√
1 + ϑ2z
2 + ϑ3z
is homogeneous of degree 0, we see that within the resonance condition (A-2)
we can replace each vector k,m and n by any colinear vectors - either integer
or not. Suppose now that there exists some positive integer d 6= 1 such that
d|k; then d|n, so that by setting
n, k0 =
k, m0 =
we finally obtain
n0 = ak0 + bk0 = a
′σi(k0) + b
′σj(m0).
The triplets (n0, ak0, bm0) and (n0, a
′σi(k0), b
′σj(m0)) further verify from the
above remark, the convolution relation (A-1) and the resonance relation (A-2).
Hence, without loss of generality, we can assume that the only positive integer
d such that d|k and d|m is 1; which we denote by
(k,m) = 1.
Equivalently,
k1Z+ k2Z+ k3Z+m1Z+m2Z+m3Z+ = Z.
Finally, suppose there exists some positive integer d 6= 1 such that d|a and
d|b. Then d|n; set
n, a0 =
a, b0 =
Observe that
Giri,j(a0k, b0m) =
Giri,j(ak, bm) = 0.
It follows from lemma 3.5 (2) of [B-M-N4] that the vector n0 of the resonant
triplet (n0, a0k, b0m) can also be written as
n0 = k̂ + m̂ with (n0, k̂, m̂) verifying (A-2).
But then
n = dn0 = ak + bm = a
′σi(k) + b
′σj(m) = dk̂ + dm̂.
From lemma 3.5 (2) of [B-M-N4], (dk̂, dm̂) must coincide with either one of
the pairs
(a′σi(k), b
′σj(m)), (b
′σj(m), a
′σi(k)).
In particular, d|a′k and d|b′m. Since d|a and (a, a′) = 1, we have (d, a′);
similarly (d, b′) = 1. But then Euclid’s lemma yields that d|k and d|m, which
contradicts the fact that (k,m) = 1. Hence we have proven that (a, b) = 1.
In a similar way, one can show that (a′, b′) = 1.
Conclusion: It follows from the above study that n ∈ Z∗ admits the two
decompositions
n = ak + bm = a′σi(k) + b
′σj(m)
(a, a′) = (b, b′) = (a, b) = (a′, b′) = 1, (k,m) = 1.
The triplets (n, ak, bm) and (n, a′σi(k), b
′σj(m)) both verify the resonant
condition (A-2) (from the homogeneity of this condition) as well as the condi-
tion of non-catalyticity (A-3). Indeed, aba′b′ 6= 0 and the condition (A-3) on
the initial triplet (n, k,m) imply that the reduced triplet (n, k,m) also verifies
(A-3)). Finally, the degeneracy condition (A-4)
Giri,j(ak, bm) = 0
is verified.
Acknowledgments. We would like to thank A.I. Bobenko, C. Bardos
and G. Seregin for very useful discussions. The assistance of Dr. B. S. Kim
is gratefully acknowledged. A.M. and B.N. acknowledge the support of the
AFOSR contract FA9550-05-1-0047.
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Introduction
Vorticity waves and resonances of elementary swirling flows
Strictly resonant Euler systems: the SO(3) case
Strictly resonant Euler systems: the case of 3-waves resonances on small-scales
Infinite dimensional uncoupled SO(3) systems
Coupled SO(3) rigid body resonant systems
|
0704.0338 | Synergistic Effects of MoDTC and ZDTP on Frictional Behaviour of
Tribofilms at the Nanometer Scale | Microsoft Word - S_Bec_ZDTP_MoDTC_Tribology_Letters.doc
Synergistic effects of MoDTC and ZDTP on frictional behaviour
of tribofilms at the nanometer scale
S. Bec1*, A. Tonck1, J.M. Georges1 and G.W. Roper2
1Laboratoire de Tribologie et Dynamique des Systèmes, UMR CNRS 5513,
Ecole Centrale de Lyon, 36 av. Guy de Collongue, 69134 Ecully Cedex, France.
2Lubricants Technology Dept., Shell Global Solutions, Shell Research and Technology
Centre, Thornton, P. O. Box 1, Chester CH1 3SH, UK.
*To whom correspondence should be addressed
Abstract
The layered structure and the rheological properties of anti-wear films generated in a
rolling/sliding contact from lubricants containing zinc dialkyldithiophosphate (ZDTP) and/or
molybdenum dialkyldithiocarbamate (MoDTC) additives have been studied by dynamic
nanoindentation experiments coupled with a simple modelling of the stiffness measurements.
Local nano-friction experiments were conducted with the same device in order to determine
the evolution of the friction coefficient as a function of the applied pressure for the different
lubricant formulations. For the MoDTC film, the applied pressure in the friction test remains
low (<0.5 GPa) and the apparent friction coefficient is high (µ>0.4). For the tribofilms
containing MoDTC together with ZDTP, which permits the applied pressure to increase up to
a few GPa through some accommodation process, a very low friction domain appears
(0.01<µ<0.05), located a few nanometers below the surface of the tribofilm. This low friction
coefficient is attributed to the presence of MoS2 planes sliding over each other in a favourable
configuration obtained when the pressure is sufficiently high, which is made possible by the
presence of ZDTP.
Keywords : ZDTP, MoDTC, tribofilm structure, nanoindentation, mechanical properties,
nanofriction, low friction.
1. Introduction
In addition to zinc dialkyldithiophosphate (ZDTP) additives, extensively used for their
exceptional antioxidant and anti-wear properties under boundary conditions in automotive
engines, lubricating oils contain several additives, among which there are detergent and
dispersant additives whose main role is to keep oil insoluble contaminants and degradation
products in suspension, at elevated temperature for the detergent additives, and at low
temperatures for the dispersant ones. Organo molybdenum compounds such as molybdenum
dithiocarbamate (MoDTC) are also used as friction modifiers for energy saving. However,
when used together in formulated oils, additives interact in various ways resulting either in
synergies or in adverse effects affecting the oil performance regarding anti-wear and friction
behaviour, and modifying the characteristics of the protective surface films generated during
friction (tribofilms). A lot of investigations have been conducted to evaluate the performances
of additive mixtures and to determine the composition of associated tribofilms. Several
factors were identified as playing a role: additive structure [1, 2], additives concentration [3-
6], base oil nature [7, 8], …, or combinations of these parameters. A detailed review on
published information on that topic was written by Willermet [9]. Non-chemical parameters
such as characteristics of the solid antagonists (hardness, roughness) or test conditions (load,
temperature, sliding speed) [3, 10] also might influence the additive interactions.
Among this variety of additive interactions, we will focus on that between ZDTP and
MoDTC, extensively studied through chemical investigations. All published works agree
upon the fact that friction and anti-wear performances of oils are improved when ZDTP and
MoDTC are used together. The formation of molybdenum disulphide (MoS2) on the rubbing
surfaces has been evidenced by several authors [11, 12]. Using UHV friction tests, coupled
with high-resolution TEM observation of wear debris and spectroscopic studies, Grossiord et
al. has given evidence for the mechanism of single MoS2 sheet lubrication [13].
The aim of this paper is to enlarge the knowledge of the local mechanical and frictional
properties of anti-wear tribofilms to those of films obtained from lubricants containing
different additives (ZDTP, MoDTC, detergent/dispersant) or mixtures of additives, in order to
explore the ZDTP/MoDTC synergy on a mechanical point of view. The only published results
on that topic are the recent papers from Ye et al. who performed AFM observations and
nanoindentation measurements on ZDTP and ZDTP + MoDTC tribofilms [14, 15].
In the present study, nanoindentation tests with continuous stiffness measurements were
performed on unwashed and solvent-washed tribofilms to determine their mechanical
properties. The frictional behaviour of the tribofilms was investigated through local
nanofriction experiments, conducted with the same device. The evolution of the friction
coefficient as a function of the applied pressure for the different lubricant formulations
leading to different tribofilms has been determined.
2. Preliminary results obtained on ZDTP anti-wear tribofilms
The structure and the rheological properties of anti-wear films from a zinc
dialkyldithiophosphate (ZDTP) solution generated in a rolling/sliding contact, simulating
engine valve train conditions, have been studied in detail with analytical and surface force
tools and the results have been published by the authors in a previous paper [16]. As preamble
to the present paper, only the main points are summarised here.
The ZDTP solution was a commercial secondary alkyl ZDTP additive at 0.1% weight of
phosphorus in a highly refined base oil. The ZDTP anti-wear films have a complex structure
that has been determined by extensive use of surface analytical techniques. It has been shown
that the ZDTP films consisted of at least three non-homogeneous layers: on the steel surface,
there is a sulphide/oxide layer, which is almost completely covered by a protective phosphate
layer, with the addition of a viscous overlayer of ZDTP degradation precipitates (alkyl
phosphate precipitates). This latter layer was removed when the film was washed with an
alkane solvent. Therefore, the properties of the ZDTP films have been studied both before and
after solvent washing with n-heptane. First, sphere/plane squeeze experiments were
performed with a surface force apparatus (SFA) on unwashed films, showing that the
overlayer of alkyl phosphate precipitates was heterogeneous and discontinuous, with a
thickness of about 900 nm. Second, the mechanical properties were obtained from
nanoindentation experiments, performed after replacing the sphere by a diamond tip, and
coupled with in-situ topographic imaging procedures to measure the contact area. From the
indentation experiments, the properties of the films were determined from normal stiffness
measurements and through the application of a rheological film model. On the unwashed
specimens, the viscous layer of alkyl phosphate precipitates was detected by the indentation
tests. It is a very soft layer, mobile under the diamond tip, with a thickness of a few hundred
of nanometers, which was in good agreement with that of sphere/plane experiments. It was
also shown that indentation experiments removed this overlayer in the proximity of the tip,
probably through a shear flow mechanism. This procedure can be compared to a soft
"mechanical" sweep and the mechanical properties of ZDTP tribofilms after such a cleaning
were found to be similar to those of solvent washed specimens. The solvent washed
tribofilms, comprising sulphide and phosphate layers, exhibited an elastoplastic behaviour
and, during the loading stage of the indentation, the hardness and the Young's modulus of the
phosphate layer increased from their initial values of about 2 GPa for the hardness and
between 30 and 40 GPa for the Young's modulus. In particular, the initial hardness of the
polyphosphate layer at the beginning of the indentation tests was close to the mean applied
pressure during the films generation. This suggested that the layer accommodated the contact
pressure in the tribotest or during the loading stage of the indentation, and could thus be
regarded as a final and local pressure sensor. The characteristics of the full ZDTP films ensure
gradual changes in mechanical properties between the substrate, bonding layers and outer
layers with the viscous overlayer serving as the tribofilm's precursor. The properties of these
layered films can thus adapt to a wide range of imposed conditions and provide appropriate
level of resistance to contact between the metal surfaces. As the severity of loading increases,
so too do the resistive forces within the film. This ensures that the shear plane remains located
inside the ZDTP protective film, which explains the exceptional efficiency of ZDTP films as
anti-wear films.
3. Experimental
3.1. Tribofilms
The tribofilms were generated at Shell Research and Technology Centre, Thornton, U.K.,
with a reciprocating Amsler machine [17] designed to simulate the contact conditions of the
cam/follower system in an internal combustion engine valve train. A flat block specimen
(8 mm x 8 mm size, 4 mm thick) has a reciprocating motion in loaded contact with a rotating
disc. The block and the disc were made in through-hardened EN31 steel. Special care was
taken with the roughness of the blocks which were polished until the average roughness was
Ra = 0.01 µm. The movement of the block was driven by a crank linked to the motion of the
disc axis through a gearbox. The block motion was approximately sinusoidal and at the same
frequency as the disc rotation. Load was applied to the contact by a spring arrangement,
acting through a roller bearing. The surface in contact with the loading bearing (rear surface
of the reciprocating element) was curved to permit self-alignment between the block and the
disc. The films were generated at a normal load of 400 N (mean contact pressure of 0.36
GPa), speed of 600 rev/min., block temperature of approximately 100°C for 5 hours. The
lubricants consisted of a highly refined base oil with different commercial additives (details of
the oil formulation are not relevant to the present work):
- MoDTC solution,
- ZDTP + MoDTC solution,
- ZDTP + MoDTC + detergent/dispersant solution ("full formulation").
The rubbing area on the polished block was typically 5 mm long in the sliding direction.
Previous analyses have shown that the composition in the centre of the wear track was
reasonably uniform, while the composition within 1 mm of the ends of the wear track could
vary significantly. The mechanical measurements on the films with the Surface Force
Apparatus have been performed in the central area of the wear track. An additional unworn
and polished block was used to obtain reference values for the EN31 steel substrate.
To preserve the film structures, the blocks were stored in the base oil (containing
predominantly paraffinic hydrocarbons, with very low concentration of polar compounds)
immediately after production of the films in the reciprocating Amsler tests and they were
immersed again, when not in use.
3.2. Surface Force Apparatus
The Ecole Centrale de Lyon Surface Force Apparatus (SFA) used in these experiments has
been described in previous publications [18, 19]. The general principle is that a macroscopic
spherical body or a diamond tip can be moved toward and away from a planar one (the ZDTP
specimen) using the expansion and the vibration of a piezoelectric crystal, along the three
directions, Ox, Oy (parallel to the plane surface) and Oz (normal to the plane surface). The
plane specimen is supported by double cantilever sensors, measuring quasi-static normal and
tangential forces (respectively Fz and Fx). Each of these is equipped with a capacitive sensor.
The sensor's high resolution allows a very low compliance to be used for the force
measurement (up to 2 x 10-6 m/N). Three capacitive sensors were designed to measure relative
displacements in the three directions between the supports of the two solids, with a resolution
of 0.01 nm in each direction. Each sensor capacitance was determined by incorporating it in
an LC oscillator operating in the range 5 - 12 MHz [20].
3.3. Tests methodology
All the experiments were conducted at room temperature. Preliminary results obtained on
anti-wear films from a ZDTP solutions have shown that n-heptane washing damages the film
[16]. That is why the blocks were tested first as obtained from the Amsler friction test,
without any cleaning and second after washing with n-heptane. The unwashed specimens
were mounted on the SFA as taken from the storage base oil. Excess of base oil was simply
removed by placing the side of the specimen on absorbing paper, which allowed the surface to
be always preserved by an oil film (thickness > 10 µm).
Nanoindentation tests
The aim of these tests was to determine the elastoplastic properties of the tribofilms (hardness
and Young's modulus) and their “mechanical” structure (number of layers and estimation of
the thickness of each layer that constitutes the film). The method used to perform
nanoindentation experiment with the SFA has already been published in detail [21]. Specific
procedures have been developed for the characterisation of ZDTP tribofilms and have been
described in previous papers [16, 22]. In this study, the determination of the near surface
mechanical properties (first nanometers) was obtained through a specific tip shape calibration,
performed on a gold film deposited by magnetron sputtering onto a silicon substrate. This
film was very smooth (peak to valley roughness around 1 nm, measured on a scan length of
1 µm) and its hardness was constant versus depth from the surface and until the penetration
depth equals the gold film's thickness [21].
For the nanoindentation experiments, a trigonal diamond tip with an angle of 115.12° between
edges (Berkovitch type) was used. The indentation tests were performed in controlled
displacement mode. The standard set-up included the continuous quasi-static measurements
of the resulting normal force Fz versus the normal displacement Z, at a slow penetration
speed, generally 0.1 to 0.5 nm/s. It also included the simultaneous measurements of the
rheological behaviour (dissipative and conservative or elastic contributions) of the tested
surface, thanks to simultaneous small sinusoidal motions at a frequency of 37 Hz, with an
amplitude of about 0.2 nm RMS. Furthermore, using the Z feedback in the constant force
mode and the tangential displacement of the indenter, the surface topography was imaged
before and after the indentation test, with the same diamond tip. This was made practically
possible because of the partial elastic recovery during the unloading cycle and hence the
geometry of tip and indent were different which was necessary to permit resolution of the
indent. For this scanning procedure, a constant normal load of 0.5 µN was typically used.
Such in-situ imaging procedure enables the operator to choose precisely the location of the
indentation test on the surface and, after the test, to quantify the plastic pile-up around the
indent and thus to measure the actual contact area.
Rheological film model
The elastic properties of the films were very difficult to extract from the indentation tests
because of both the influence of the substrate and of the film structure itself. They were
obtained through the stiffness measurements, which are global (film+substrate)
measurements. To extract the properties of each layer of the film, a simple model has been
developed, and its main features are described as follows. The experimental stiffness versus
normal displacement curve was identified with the elastic response of a structure composed of
one or two homogeneous elastic layers on a substrate (semi-infinite elastic half space)
indented by a rigid cylindrical punch of radius a. For such a system, modelled by two springs
connected in series [23], the calculated global stiffness (Kz) depends on the reduced Young's
modulus of the substrate (Es* with Es*=Es/(1-νs2)), measured on an unworn steel block, over
the contact radius (a) and depends also on four unknown parameters which are the reduced
Young's modulus (Ef*, Ef*=Ef/(1-νf2)) and the thickness (t) of each layer. For each test, their
values were adjusted to obtain a good fit between the measured stiffness curve and the
calculated one. This procedure provided the structure (one or two layers), the thickness and
the reduced Young's modulus of each layer that constituted the tribofilms. Details are given in
a previous paper [16]. Following this model, the global stiffness of a single layer system is
given by:
22K t a
a E aEz f s
π π * *
(1)
This simple model describes perfectly the behaviour of model systems such as gold layers on
a silicon substrate [21]. In the case of tribofilms, deviations may be observed at a critical
pressure or at a critical depth from which the experimentally measured stiffness may be found
to exceed significantly the theoretical one. This is interpreted as a change in the surface
properties due to the applied pressure and appears to be related to a measured hardness
increase. Indeed, as the applied pressure can reach values much larger than the initial hardness
value of the surface, the resulting plastic flow may induce a small volume reduction and
molecular rearrangements which could be sufficient to induce a noticeable change in the
mechanical properties. From a threshold pressure value, H0, the stiffness curve was then
influenced both by the substrate's elasticity and by the change in mechanical properties. This
pressure dependence can be introduced in the model by writing that in the deformed volume
of material, when H>H0 (i.e. when the film accommodates the applied pressure through
hardness increase), the film modulus Ef* is proportional to the hardness (the ratio Ef*/H
remains constant). It gives the following equation:
EE = (2)
Ef0* is the reduced Young's modulus value, when the applied pressure is equal to or lower
than the threshold pressure H0. When necessary, by introducing this effect in our modelling
and by adjusting the value of the threshold pressure, we were able to fit correctly the whole
stiffness curve. An example of such fit is given on figure 1. The evolution of the film modulus
Ef* versus plastic depth can also be extracted from equation 1 using the experimentally
measured global (film+substrate) stiffness values Kz and the film's thickness, t, independently
of equation 2. This permits a check on whether it is proportional to the hardness as assumed in
equation 2. In the example shown figure 2, the calculated Young's modulus of the film (from
equation 1 with a film thickness t=25 nm) is found to be proportional to the measured
hardness with a mean ratio Ef*/H=16.5, in good agreement with the ratio
Ef0*/H0=17/1.05=16.2 obtained from the stiffness fit.
Full formulation (ZDTP+MoDTC+detergent/dispersant)
Solvent washed tribofilm
0 10 20 30 40 50
Penetration depth (nm)
Measured stiffness
Calculated stiffness, t=25 nm, Efo*=17 GPa,
without pressure accommodation
Calculated stiffness, t=25 nm, Efo*=17 GPa,
with pressure accommodation, Ho=1.05 GPa
Figure 1: Example of application of the rheological film model: measured and calculated
global stiffness for a tribofilm obtained from the full formulation (MoDTC + ZDTP +
detergent/dispersant). A good fit between the measured and the calculated values is obtained
with a single layer system (thickness t=25 nm and reduced Young's modulus Ef0*=17 GPa)
and a pressure accommodation effect from a threshold pressure H0=1.05 GPa.
t = 25 nmEf0* = 17 GPa H0 = 1.05 GPa
Full formulation (ZDTP+MoDTC+detergent/dispersant)
Solvent washed tribofilm
0 10 20 30 40 50 60
Plastic depth (nm)
Reduced Young's modulus of the tribofilm, Ef*
Hardness of the tribofilm, H
t = 25 nmEf0* = 17 GPa H0 = 1.05 GPa
Full formulation (ZDTP+MoDTC+detergent/dispersant)
Solvent washed tribofilm
0 10 20 30 40 50 60
Plastic depth (nm)
Reduced Young's modulus of the tribofilm, Ef*
Hardness of the tribofilm, H
Figure 2: Example of evolution of film's reduced Young's modulus and hardness versus
plastic depth, for a tribofilm obtained from the "full formulation" (MoDTC + ZDTP +
detergent/dispersant). The Young's modulus of the film is calculated using equation 1 with the
measured stiffness values and using only the film's thickness determined from the fit shown
figure 1 (t = 25 nm).
Nanofriction experiments
Nanofriction experiments were conducted on the blocks by moving the diamond tip along Ox
direction (parallel to the surface) at low speed (2 to 5 nm/s) along a distance of 0.5 µm. The
objective of these tests was to determine how the friction coefficient varies as a function of
the applied pressure. The tests were conducted at monitored increasing depth. During the
tests, the normal, Fz, and the tangential, Fx, forces were recorded, which allowed us to
calculate the apparent friction coefficient µ=Fx/Fz (see example figure 3).
Full formulation (ZDTP+MoDTC+detergent/dispersant)
Solvent washed tribofilm
0 20 40 60 80 100 120 140
Time (s)
µ=Fx/Fz
0 5 10 15 20 25
Penetration depth (nm)
Indentation test
Nanofriction test Smaller contact area
Figure 3: Procedure used for the nanofriction tests. The diamond tip is oriented edge first and
the nanofriction tests are conducted at monitored increasing depth. During the test, the
normal (Fz) and tangential (Fx) forces are recorded. The friction coefficient µ=Fx/Fz is
calculated.
Large pile-up was observed in the case of nanofriction with the diamond tip oriented face
first, which may induce large uncertainty in the calculation of the contact area. That is why
the nanofriction tests were conducted edge first. In these conditions, the estimation of the
applied pressure at a given depth was obtained using low load nanoindentation tests, made in
the near proximity of the nanofriction tests. Assuming that, at a given depth, the hardness of
the tribofilm should be the same for the friction test and for the near indentation test, the
contact area, and then the applied pressure, were obtained from the difference between the
normal force measured for the two tests at the same depth (see insert on figure 3). Using the
in-situ imaging procedure, figure 4 shows an example of an image of the surface of a
tribofilm after such a nanofriction test.
100 nm
Beginning
of the testBeginning of the wear
Direction of friction
100 nm
Beginning
of the testBeginning of the wear
Direction of friction
Figure 4: Typical image of the surface
of a tribofilm after a nanofriction
experiment. The image is obtained with
the in-situ imaging procedure.
4. Results
The first part presents the mechanical properties of the different tribofilms, determined from
the nanoindentation experiments. Their structure, one or two layers, and their thickness were
deduced from the use of our rheological film model.
The results concerning the frictional behaviour of the tribofilms are given in a second part.
4.1. Structure and mechanical properties of the tribofilms
MoDTC tribofilms
The tribofilm obtained from base oil + MoDTC has been tested without washing and after
washing with n-heptane. Even on the solvent washed block, it was not possible to make any
local topographic image nor line scanning preliminary to the indentations tests, revealing that
the film was very soft and was easily damaged by the diamond tip. Representative hardness
curves obtained on the MoDTC tribofilms are shown on figure 5.
MoDTC tribofilm
0 40 80 120 160
Plastic depth (nm)
Unwashed tribofilm
Solvent-washed tribofilm
Figure 5: Typical hardness curves obtained on the MoDTC tribofilms. Open symbols
correspond to hardness curves obtained on the unwashed film. Black symbols correspond to
hardness curves obtained on the solvent-washed film.
Very low mechanical properties were measured on the unwashed MoDTC tribofilm. The
surface hardness ranged from 0.02 to 0.1 GPa indicating the presence of a very soft overlayer
covering the tribofilm.
After washing with n-heptane, the indentation tests showed that this overlayer has been
removed by the washing procedure. The remaining tribofilm was a soft homogeneous layer,
whose hardness was typically in the range 0.4 - 0.5 GPa at the beginning of the tests.
Adhesion to the diamond tip was detected at the end of the unloading part of the tests. The
film thickness and the structure (number of layers) have been obtained from the stiffness
measurements performed during the experiments using the rheological film model.
The film appeared to be homogeneous in its thickness, and for most of the tests, its elastic
behaviour corresponded to that of a single layer, with constant properties versus depth. The
thickness of the film was found to be between 30 and 75 nm. The reduced Young's modulus
was typically equal to 7 – 8 GPa.
ZDTP + MoDTC tribofilms
From optical observation, the ZDTP + MoDTC unwashed film was very thin. This was
confirmed by the indentation tests. Prior to any contact, a very soft layer, 60 to 120 nm thick,
was detected at the surface of the unwashed film.
Indentation tests conducted after scanning or imaging the surface of the unwashed film
("mechanical sweep") showed that the film was spatially heterogeneous. Its thickness and its
mechanical properties varied depending on the test location:
- In some places, only a very thin layer (a few nanometers thick) with a reduced Young's
modulus of 50 GPa covered the work-hardened steel substrate (tests A and B on figure 6).
- A thicker layer (15 to 30 nm) with a reduced Young's modulus of 50 – 80 GPa was found
in other places (tests C and D on figure 6), sometimes with accommodation pressure
effect (threshold pressure H0 = 4.8 GPa). Such layer behaves like the sulphide-oxide layer
of the ZDTP tribofilm [16].
- Elsewhere, the structure of the tribofilm was more complex, with a soft layer covering a
stiffer one. For example, test E on figure 6 corresponds to a soft layer, 12 nm thick, with
properties comparable to those of the MoDTC tribofilm (hardness of 0.2 GPa and reduced
Young's modulus of 5 GPa) which covers a stiffer layer, 18 nm thick, with a reduced
Young's modulus of 50 GPa.
This heterogeneity was confirmed by the indentation tests conducted on the solvent-washed
ZDTP + MoDTC tribofilm, where, at least, three different types of film were identified:
- In some places, the film behaved like a one layer system, able to accommodate the
pressure (pressure threshold 2.8 GPa). Its thickness was between 35 nm and 150 nm. The
surface hardness was about 2 – 3 GPa and the reduced Young's modulus was about
55 - 65 GPa.
- In other places, the film behaved like a bilayered structure: a surface layer, about 25 nm
thick, with properties comparable to those of the MoDTC tribofilm (hardness of
0.3 - 0.4 GPa, reduced Young's modulus of 8 GPa), covers a stiffer layer, 150 nm thick,
with a reduced Young's modulus of about 80 GPa.
- Elsewhere, the surface film was between 3 and 15 nm thick, with properties comparable to
the lower properties measured on the ZDTP tribofilm (hardness about 1 – 1.5 GPa and
reduced Young's modulus about 10 GPa). For some tests, this surface film was able to
accommodate the pressure, with a pressure threshold of 1 – 1.5 GPa. It covers a stiffer
layer, 10 to 55 nm thick, with a reduced Young's modulus varying from 60 to 110 GPa.
ZDTP+MoDTC tribofilm
Unwashed block
0 20 40 60 80 100 120
Plastic depth (nm)
Prior to any contact
After imaging - test A
After imaging - test B
After imaging - test C
After imaging - test D
After imaging - test E
Figure 6: Representative hardness curves obtained on the unwashed ZDTP + MoDTC
tribofilm, prior to any contact and after the imaging procedure. The film is spatially
heterogeneous in thickness and in mechanical properties.
ZDTP + MoDTC + detergent/dispersant tribofilms ("full formulation" tribofilms)
Nanoindentation tests performed in fresh areas, prior to any contact showed that, at the
surface of the unwashed tribofilm, there was a very soft layer, mobile under the diamond tip,
with an apparent thickness of a few hundreds of nanometers.
Representative hardness curves obtained on the unwashed block near these initial contacts are
shown figure 7. Contrary to the ZDTP + MoDTC tribofilm, the film was found to be spatially
homogeneous. Only its thickness was found to vary, depending on the tested area. A very thin
softer layer was detected at the surface of the tribofilm, which did not resist to imaging nor
scanning, except if the normal load was very low (lower than 0.3 µN). This layer had a
hardness value (about 0.3 – 0.4 GPa) comparable to the hardness value of the MoDTC
tribofilm. The observed large hardness increase when the load increased also indicated that
the tribofilm had a great capability to accommodate the applied pressure. This result was
confirmed by the interpretation of the stiffness measurements using the rheological model,
which also showed that the tribofilm had a complex structure. At its surface, there was first a
layer with a thickness of only a few nanometers (2 nm to 7 nm) and a reduced Young's
modulus of 10 - 15 GPa. Then, there was a second layer (thickness between 20 nm and 140
nm) with a higher reduced Young modulus of 65 – 80 GPa.
A similar tribofilm was tested after n-heptane washing. It also had a great ability to
accommodate the applied pressure. From the stiffness measurements, on most places, the film
was found to behave like a film constituted by two layers. The surface layer was thin (5 to 25
nm) with a reduced Young's modulus value in the range 15 – 20 GPa. The thickness of the
underlayer was found to vary between 0 (no underlayer, example of figures 1 and 2) and 100
nanometers and its elastic modulus was in the range 110 - 120 GPa.
ZDTP + MoDTC + detergent/dispersant
Unwashed tribofilm
0 10 20 30 40 50 60
Plastic depth (nm)
First test, prior to any contact
Without preliminary scanning
After scanning or imaging
Figure 7: Representative hardness curves obtained on the unwashed ZDTP + MoDTC +
detergent/dispersant tribofilm ("full formulation"), prior to any contact and in the region near
the first contacts, either without preliminary surface scanning or after scanning/imaging
procedure.
Figure 8 compares representative hardness curves for all tested tribofilms. For the ZDTP +
MoDTC tribofilms, three curves are plotted because of the variety of obtained results
revealing the spatial heterogeneity of this tribofilm. A representative hardness curve for the
ZDTP tribofilm tested in the same conditions in a previous study [16] has been added for
comparison.
0 5 10 15 20 25 30 35 40
Total penetration depth (nm)
ZDTP, solvent washed MoDTC, solvent washed
ZDTP + MoDTC, unwashed (2 tests) ZDTP + MoDTC, solvent washed
Full formulation, unwashed Full formulation, solvent washed
Figure 8: Comparison of the hardness curves obtained on the different tribofilms. The
hardness curve obtained for a ZDTP anti-wear tribofilm obtained from a previous study is
plotted for comparison.
4.2. Nanofriction experiments
Nanofriction experiments were conducted on the three preceding tribofilms and also on a
ZDTP tribofilm and on a ZDTP + detergent/dispersant tribofilm. In order to simplify the
following graphs, only one representative curve was plotted for each tribofilm (or two when it
was necessary to illustrate the dispersion when it was significant).
Figure 9 shows the evolution of the friction force versus the normal force for the tested
tribofilms. For a given formulation, there was very little difference between the results
obtained on unwashed and on solvent washed tribofilms at low load, indicating that the
solvent washing does not seem to affect the frictional behaviour of the tribofilm. This agrees
with the idea that the soft viscous overlayer is supposed to serve as precursor for the tribofilm
rather than that it plays a mechanical role during friction.
0 3 6 9 12 15
Normal force, Fz (µN)
ZDTP, solvent washed ZDTP + Det/Disp, solvent washed
MoDTC, solvent washed ZDTP + MoDTC, solvent washed
ZDTP + MoDTC, solvent washed ZDTP + MoDTC, unwashed
Full formulation, unwashed Full formulation, solvent washed
Figure 9: Friction force (Fx) versus normal force (Fz) during nanofriction tests with
increasing penetration depth for different tribofilms.
It is also worth noting that the heterogeneity in mechanical properties found on the ZDTP +
MoDTC tribofilm also exists in the frictional properties. For this tribofilm, the friction force
at low normal loads may be comparable either to the friction force obtained for the ZDTP
tribofilm or to the friction force obtained for the "full formulation" tribofilm.
Under the present testing conditions, it can be observed that the lower friction forces were
obtained for films containing MoDTC together with ZDTP. The higher were obtained for the
tribofilm from MoDTC alone.
Figure 10 shows the evolution of the friction coefficient versus mean pressure. The existence
of low friction coefficient values (0.01<µ<0.05) appears to be related both to the presence of
MoDTC additive in the initial lubricant and to the ability for the tribofilm to reach sufficiently
high pressure values (1.5 – 3 GPa) during the friction test. Thus, the MoDTC tribofilm, which
is not able to resist to the contact pressure by increasing its mechanical properties seems to be
ineffective in reducing friction, contrary to the tribofilms containing ZDTP and MoDTC
together, which are able to accommodate the contact pressure by increasing their mechanical
properties. Nevertheless, both behaviours (high or low friction) were observed for the ZDTP
+ MoDTC tribofilms. This is certainly due to the spatial heterogeneity of these tribofilms,
which behave on some places like ZDTP tribofilms, or elsewhere like "full formulation"
tribofilms. It was also observed that tribofilms formed without MoDTC were ineffective in
reducing friction even if high contact pressures were reached during the friction tests.
0 1 2 3 4 5 6
Mean pressure P (GPa)
ZDTP, solvent washed ZDTP + Det/Disp, solvent washed
MoDTC, solvent washed ZDTP + MoDTC, solvent washed
ZDTP + MoDTC, solvent washed ZDTP + MoDTC, unwashed
Full formulation, unwashed Full formulation, solvent washed
Figure 10: Apparent friction coefficient versus mean pressure for the different tested
tribofilms.
When the evolution of the friction coefficient is plotted versus penetration depth (figure 11), it
appears that, when it existed, the low friction coefficient domain was detected a few
nanometers below the surface of the tribofilm. It also shows that, for the full formulation, the
low friction domain was deeper for the unwashed tribofilm than for the solvent washed one.
The unwashed tribofilm appears to be covered by a surface layer with rather bad frictional
properties, which can be removed by solvent washing or by "mechanical" sweep (low load
scanning procedures for example).
0 2 4 6 8 10 12 14 16 18 20
Penetration depth (nm)
ZDTP, solvent washed ZDTP + Det/Disp, solvent washed
MoDTC, solvent washed ZDTP + MoDTC, solvent washed
ZDTP + MoDTC, solvent washed ZDTP + MoDTC, unwashed
Full formulation, unwashed Full formulation, solvent washed
Figure 11: Apparent friction coefficient versus penetration depth for the different tested
tribofilms.
5. Discussion
Because of the inhomogeneous and patchy nature of anti-wear tribofilms and of their low
thickness, very few results are published concerning their mechanical properties [24-28].
Moreover, the differences in sample preparation and the diversity of used techniques and
experimental procedures render delicate the comparison of the obtained results. For example,
the Young’s modulus values given by Aktary et al. for a ZDTP tribofilm [28] are significantly
higher that those we measured but one explanation can be that they did not take into account
the substrate’s elasticity in their calculations, contrary to what is done in the current study. Or
if we attempt to compare our results with those recently published by Ye et al. on ZDTP and
ZDTP + MoDTC tribofilms [14, 15], this reveals significant differences. For example, Ye et
al. found that both tribofilms possess the same hardness and modulus depth distributions,
corresponding to continuously and functionally graded materials, when in the present work,
the hardness curves for similar tribofilms did not coincide and the use of our rheological film
model allowed us to describe the tribofilms as layered materials with properties adaptable to
contact conditions. The hardness and modulus values, respectively 10 GPa at a contact depth
of 30 nm and 215 GPa at a depth of 20 nm, that they reported are also significantly higher
than those we measured and also higher than those given by Aktary et al. This could be due to
differences in sample preparation and also certainly to the use of different methods and
assumptions for the treatment of the nanoindentation data.
Concerning the frictional behaviour of the tribofilms, the presented nanofriction tests were
conducted in unlubricated conditions, at very low speed (2 to 5 nm/s) and the measured
nanofriction coefficients corresponded to the friction between the diamond tip and the
tribofilm (over its steel substrate). That is why it also seems difficult to compare our values to
macroscopic friction coefficient values obtained on classical tribometers. The latter are
representative of steel on steel contact in the presence of a tribofilm and are averaged over the
whole contact surface. However, our local values are not far from the end of test Amsler
macroscopic friction coefficient values published by Pidduck and Smith [25] for ZDTP,
ZDTP + detergent/dispersant and ZDTP + friction modifier tribofilms. Moreover, these
macroscopic values were found to be proportional, with a factor 0.7, to micro-friction
coefficient values measured with Lateral Force Microscopy by the same authors, making
them suggest that there may be a link between macro and micro-frictional behaviour of
smooth regions of anti-wear tribofilms. Unfortunately, no tribofilm obtained from friction
modifier alone were tested in this study, with which we could compare our results.
Nevertheless, macroscopic friction coefficient values, in the range 0.10 – 0.14, measured on
an alternative ball on plane tribometer were reported by Muraki and Wada [6] for oil
containing MoDTC alone. They conclude that such lubricant was ineffective in reducing
friction, contrary to the oil containing MoDTC together with ZDTP. More recently, similar
high macroscopic friction coefficient values (in the range 0.095 – 0.2) were measured by
Unnikrishnan et al. for oil containing MoDTC alone [29]. On the other hand, Grossiord et al.
reported very low steady-state friction coefficient (0.04) measured for base oil + MoDTC
during SRV friction tests, and a lower steady-state value (0.02) for friction tests in a UHV
tribometer, carried out by sliding a macroscopic hemispherical steel pin again a flat covered
by a MoDTC tribofilm [13]. From tests carried out in a high frequency reciprocating rig,
Graham et al. [30] also reported that, in the absence of ZDTP, MoTDC additives were
effective in reducing friction at a combination of high additive concentration and high
temperature (up to 0.4% wt. and 200°C). Such diversity of results, certainly partly due to the
various tests conditions, makes unreasonable a comparison between the very high
nanofriction coefficient measured on the MoDTC tribofilm under the present testing
conditions and those published values. As, regarding the literature, the formation of MoS2 was
well established for MoDTC containing lubricants, the question is how can we explain such
high friction coefficient during the nanofriction tests ? Or what caused the very low friction
observed when ZDTP was used together with MoDTC ? From figure 10, the low friction
coefficient values (0.01<µ<0.05) were observed for the MoDTC containing lubricants when
the contact pressure was in the range 1.5 – 3 GPa (the question of the spatial heterogeneity of
the ZDTP + MoDTC tribofilm will be discussed latter). These high pressures were measured
for tribofilms able to increase their mechanical properties, thus accommodating the contact
conditions, which was demonstrated to be the case for ZDTP anti-wear tribofilms [16]. On the
other hand, high pressures were not reached for the soft MoDTC tribofilm. Thus, the easy
sliding of the MoS2 sheets could result from a favourable orientation induced by sufficiently
high contact pressure values. The ability of MoS2 sheets to orient in a favourable direction
was reported by Grossiord et al. [31] and Martin et al. [32], who recently investigated
tribochemical interactions between ZDTP, MoDTC and OCB (overbased detergent calcium
borate) additives. Using high resolution TEM observations of wear debris, coupled with wear
scar micro-spot XPS analysis, they observed perfectly oriented MoS2 sheets, with their basal
plane parallel to the flaky wear fragments. Such "mechanical" interpretation of the role of the
contact pressure agrees with previous work of Muraki et al. who studied the effect of roller
hardness on the rolling sliding characteristics of MoDTC in the presence of ZDTP and
concluded that the friction reduction effect increased with higher degree of roller hardness
[10]. Yamamoto also reported that a necessary condition for improving the friction and wear
characteristics of a lubricant was the formation of surface films composed of iron phosphates
with high hardness and Mo-S compounds [11]. Concerning the spatial heterogeneity of the
ZDTP + MoDTC tribofilms, it can be worth noting that using high resolution TEM
observations of wear debris collected after friction tests, coupled with AES and XPS studies
of rubbing surfaces, Grossiord et al. described the ZDTP + MoDTC tribofilm as being
composed of a mixture of glassy zinc phosphate zones containing molybdenum, and carbon-
rich zones containing zinc and highly-dispersed MoS2 single sheets [13, 33].
The observation that, during the nanofriction tests, the low friction domain was located a few
nanometers below the surface also corroborates this interpretation. As the nanofriction tests
were conducted at increasing depth, the sufficiently high pressures were obtained after a few
nanometers penetration depth inside the MoS2 containing layer (with properties similar to the
MoDTC tribofilm), thanks to the presence of the underneath resisting anti-wear layer, whose
characteristics are similar to those of the phosphate layer of the ZDTP tribofilm.
Finally, combining the results obtained from the nanoindentation and nanofriction
experiments, we can propose a possible schematic description of the anti-wear tribofilms
obtained from the "full formulation" oil. Some assumptions are also made on what happened
during nanofriction tests on such tribofilms (see figure 12 on which for convenient drawing,
as the Berkovitch diamond tip is not sharp, it was represented by a flat punch).
A soft layer containing non-oriented MoS2 sheets is present at the surface of the tribofilm
(layer (a) in figure 12). This layer, 0 to 25 nm thick, has mechanical properties comparable
with those of the MoDTC tribofilm (0.3 – 0.5 GPa for the hardness and 3 – 10 GPa for the
reduced Young's modulus). Its friction coefficient is rather high. This layer is easily damaged
or removed by the diamond tip during imaging or line-scanning procedures. When the contact
pressure is sufficiently high, friction induces a favourable orientation of the MoS2 sheets, over
a thickness of 1 or 2 nanometers (layer (b) in figure 12), resulting in very low friction
coefficient values which combine with the anti-wear efficiency of the tribofilm. Under this
layer, there is then an anti-wear layer (layer (c) in figure 12), with properties similar to those
of the polyphosphate layer of the ZDTP tribofilm. Then, just over the substrate (noted (e)in
figure 12), there is a bonding layer (layer (d) in figure 12) with high mechanical properties
(oxides, sulfides).
Figure 12: Possible schematic description of the anti-wear tribofilm obtained from the "full
formulation" and orientation of the MoS2 planes of the outer layer resulting from a
nanofriction tests (for convenient drawing, as the Berkovitch diamond tip is not sharp, it was
represented by a flat punch). The thickness of each layer is arbitrary drawn as it varies
significantly depending on the tested area (from zero when the layer is not present to a few
tens of nanometers).
(a) Soft layer containing non-oriented MoS2 sheets, with mechanical properties comparable
to those of the MoDTC tribofilm,
(b) Layer of favourably frictionally oriented MoS2 sheets with a typical thickness of 1 or
2 nm,
(c) Layer with properties similar to those of the polyphosphate layer of the ZDTP tribofilm,
(d) Bonding layer with high mechanical properties (oxides, sulfides),
(e) Steel substrate.
6. Conclusions
Thanks to the combined used of (i) nanoindentation experiments with continuous stiffness
measurements coupled with imaging procedures, (ii) a specifically developed rheological film
model and (iii) nanofriction tests, synergistic effects of ZDTP and MoDTC on frictional
behaviour of anti-wear tribofilms have been evidenced from mechanical considerations. One
original feature of this study lies in the characterisation of unwashed anti-wear tribofilms with
their full structure preserved.
The structure and nanomechanical properties (hardness and reduced Young's modulus) of
tribofilms formed with different mixtures of additives (ZDTP, MoDTC, detergent/dispersant)
were first determined.
Concerning the occurrence of very low friction (0.01<µ<0.05), the contact pressure was found
to be a critical parameter. The low friction coefficient values were attributed to a favourable
orientation of MoS2 sheets present in the outer layer of the tribofilms formed from MoDTC
containing lubricants. Such a favourable orientation occurred only if sufficiently high contact
pressure was reached. These high contact pressures were attained when ZDTP was used as oil
additive together with MoDTC because one of the main characteristics of ZDTP additives is
to form protective anti-wear tribofilms under boundary lubrication, with varying structure and
properties with depth, among which is an amazing ability to increase their mechanical
properties, thus accommodating the contact conditions.
A possible schematic description of the tribofilms containing both ZDTP and MoDTC was
deduced and a mechanism was proposed to account for the mechanical synergy that occurs
during nanofriction tests on such tribofilms.
Aknowledgement
The authors thank Shell Research Limited for financial support and permission to publish.
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|
0704.0339 | Lattice Boltzmann inverse kinetic approach for the incompressible
Navier-Stokes equations | Lattice Boltzmann inverse kinetic approach for the incompressible Navier-Stokes
equations
Enrico Fonda1,Massimo Tessarotto1,2 and Marco Ellero3
1Dipartimento di Matematica e Informatica,
Università di Trieste, Italy
2Consorzio di Magnetofluidodinamica, Trieste, Italy
3Institute of Aerodynamics,
Technical University of Munich, Munich, Germany
(Dated: August 18, 2021)
In spite of the large number of papers appeared in the past which are devoted to the lattice
Boltzmann (LB) methods, basic aspects of the theory still remain unchallenged. An unsolved theo-
retical issue is related to the construction of a discrete kinetic theory which yields exactly the fluid
equations, i.e., is non-asymptotic (here denoted as LB inverse kinetic theory). The purpose of this
paper is theoretical and aims at developing an inverse kinetic approach of this type. In principle
infinite solutions exist to this problem but the freedom can be exploited in order to meet important
requirements. In particular, the discrete kinetic theory can be defined so that it yields exactly the
fluid equation also for arbitrary non-equilibrium (but suitably smooth) kinetic distribution func-
tions and arbitrarily close to the boundary of the fluid domain. This includes the specification
of the kinetic initial and boundary conditions which are consistent with the initial and boundary
conditions prescribed for the fluid fields. Other basic features are the arbitrariness of the ”equi-
librium” distribution function and the condition of positivity imposed on the kinetic distribution
function. The latter can be achieved by imposing a suitable entropic principle, realized by means of
a constant H-theorem. Unlike previous entropic LB methods the theorem can be obtained without
functional constraints on the class of the initial distribution functions. As a basic consequence, the
choice of the the entropy functional remains essentially arbitrary so that it can be identified with
the Gibbs-Shannon entropy. Remarkably, this property is not affected by the particular choice of
the kinetic equilibrium (to be assumed in all cases strictly positive). Hence, it applies also in the
case of polynomial equilibria, usually adopted in customary LB approaches. We provide different
possible realizations of the theory and asymptotic approximations which permit to determine the
fluid equations with prescribed accuracy. As a result, asymptotic accuracy estimates of customary
LB approaches and comparisons with the Chorin artificial compressibility method are discussed.
PACS numbers: 47.27.Ak, 47.27.eb, 47.27.ed
1 - INTRODUCTION - INVERSE KINETIC
THEORIES
Basic issues concerning the foundations classical hy-
drodynamics still remain unanswered. A remarkable as-
pect is related the construction of inverse kinetic theo-
ries (IKT) for hydrodynamic equations in which the fluid
fields are identified with suitable moments of an appropri-
ate kinetic probability distribution. The topic has been
the subject of theoretical investigations both regarding
the incompressible Navier-Stokes (NS) equations (INSE)
[1, 2, 3, 4, 5, 6] and the quantum hydrodynamic equations
associated to the Schrödinger equation [7]. The impor-
tance of the IKT-approach for classical hydrodynamics
goes beyond the academic interest. In fact, INSE rep-
resent a mixture of hyperbolic and elliptic pde’s, which
are extremely hard to study both analytically and nu-
merically. As such, their investigation represents a chal-
lenge both for mathematical analysis and for computa-
tional fluid dynamics. The discovery of IKT [1] provides,
however, a new starting point for the theoretical and nu-
merical investigation of INSE. In fact, an inverse kinetic
theory yields, by definition, an exact solver for the fluid
equations : all the fluid fields, including the fluid pres-
sure p(r, t), are uniquely prescribed in terms of suitable
momenta of the kinetic distribution function, solution
of the kinetic equation. In the case of INSE this per-
mits, in principle, to determine the evolution of the fluid
fields without solving explicitly the Navier-Stokes equa-
tion, nor the Poisson equations for the fluid pressure [6].
Previous IKT approaches [2, 3, 4, 5, 7] have been based
on continuous phase-space models. However, the inter-
esting question arises whether similar concepts can be
adopted also to the development of discrete inverse ki-
netic theories based on the lattice Boltzmann (LB) the-
ory. The goal of this investigation is to propose a novel
LB theory for INSE, based on the development of an IKT
with discrete velocities, here denoted as lattice Boltzmann
inverse kinetic theory (LB-IKT). In this paper we intend
to analyze the theoretical foundations and basic proper-
ties of the new approach useful to display its relation-
ship with previous CFD and lattice Boltzmann methods
(LBM) for incompressible isothermal fluids. In particu-
lar, we wish to prove that it delivers an inverse kinetic
http://arxiv.org/abs/0704.0339v1
theory, i.e., that it realizes an exact Navier-Stokes and
Poisson solver.
1a - Motivations: difficulties with LBM’s
Despite the significant number of theoretical and nu-
merical papers appeared in the literature in the last few
years, the lattice Boltzmann method [8, 9, 10, 11, 12,
13, 14] - among many others available in CFD - is prob-
ably the one for which a complete understanding is not
yet available. Although originated as an extension of
the lattice gas automaton [15, 16] or a special discrete
form of the Boltzmann equation [17], several aspects re-
garding the very foundation of LB theory still remain to
be clarified. Consequently, also the comparisons and ex-
act relationship between the various lattice Boltzmann
methods (LBM) and other CFD methods are made dif-
ficult or, at least, not yet well understood. Needless to
say, these comparisons are essential to assess the relative
value (based on the characteristic computational com-
plexity, accuracy and stability) of LBM and other CFD
methods. In particular the relative performance of the
numerical methods depend strongly on the characteris-
tic spatial and time discretization scales, i.e., the minimal
spatial and time scale lengths required by each numerical
method to achieve a prescribed accuracy. On the other
hand, most of the existing knowledge of the LBM’s prop-
erties originates from numerical benchmarks (see for ex-
ample [18, 19, 20]). Although these studies have demon-
strated the LBM’s accuracy in simulating fluid flows, few
comparisons are available on the relative computational
efficiency of the LBM and other CFD methods [17, 21].
The main reason [of these difficulties] is probably because
current LBM’s, rather than being exact Navier-Stokes
solvers, are at most asymptotic ones (asymptotic LBM’s),
i.e., they depend on one or more infinitesimal parame-
ters and recover INSE only in an approximate asymptotic
sense.
The motivations of this work are related to some of
the basic features of customary LB theory representing,
at the same time, assets and weaknesses. One of the
main reasons of the popularity of the LB approach lays
in its simplicity and in the fact that it provides an ap-
proximate Poisson solver, i.e., it permits to advance in
time the fluid fields without explicitly solving numeri-
cally the Poisson equation for the fluid pressure. How-
ever customary LB approaches can yield, at most, only
asymptotic approximations for the fluid fields. This is
because of two different reasons. The first one is the dif-
ficulty in the precise definition of the kinetic boundary
conditions in customary LBM’s, since sufficiently close to
the boundary the form of the distribution function pre-
scribed by the boundary conditions is not generally con-
sistent with hydrodynamic equations. The second reason
is that the kinetic description adopted implies either the
introduction of weak compressibility [8, 9, 11, 12, 13, 14]
or temperature [22] effects of the fluid or some sort of
state equation for the fluid pressure [23]. These assump-
tions, although physically plausible, appear unacceptable
from the mathematical viewpoint since they represent a
breaking of the exact fluid equations.
Moreover, in the case of very small fluid viscosity
customary LBM’s may become inefficient as a conse-
quence of the low-order approximations usually adopted
and the possible presence of the numerical instabilities
mentioned above. These accuracy limitations at low vis-
cosities can usually be overcome only by imposing severe
grid refinements and strong reductions of the size of the
time step. This has the inevitable consequence of rais-
ing significantly the level of computational complexity
in customary LBM’s (potentially much higher than that
of so-called direct solution methods), which makes them
inefficient or even potentially unsuitable for large-scale
simulations in fluids.
A fundamental issue is, therefore, related to the con-
struction of more accurate, or higher-order, LBM’s, ap-
plicable for arbitrary values of the relevant physical
(and asymptotic) parameters. However, the route which
should permit to determine them is still uncertain, since
the very existence of an underlying exact (and non-
asymptotic) discrete kinetic theory, analogous to the con-
tinuous inverse kinetic theory [2, 3], is not yet known.
According to some authors [24, 25, 26] this should be
linked to the discretization of the Boltzmann equation, or
to the possible introduction of weakly compressible and
thermal flow models. However, the first approach is not
only extremely hard to implement [27], since it is based
on the adoption of higher-order Gauss-Hermite quadra-
tures (linked to the discretization of the Boltzmann equa-
tion), but its truncations yield at most asymptotic the-
ories. Other approaches, which are based on ’ad hoc’
modifications of the fluid equations (for example, intro-
ducing compressibility and/or temperature effects [28]),
by definition cannot provide exact Navier-Stokes solvers.
Another critical issue is related to the numerical sta-
bility of LBM’s [29], usually attributed to the violation of
the condition of strict positivity (realizability condition)
for the kinetic distribution function [29, 30]. Therefore,
according to this viewpoint, a stability criterion should
be achieved by imposing the existence of an H-theorem
(for a review see [31]). In an effort to improve the ef-
ficiency of LBM numerical implementations and to cure
these instabilities, there has been recently a renewed in-
terest in the LB theory. Several approaches have been
proposed. The first one involves the adoption of entropic
LBM’s (ELBM [30, 32, 33, 34] in which the equilibrium
distribution satisfies also a maximum principle, defined
with respect to a suitably defined entropy functional.
However, usually these methods lead to non-polynomial
equilibrium distribution functions which potentially re-
sult in higher computational complexity [35] and less nu-
merical accuracy[36]. Other approaches rely on the adop-
tion of multiple relaxation times [37, 38]. However the
efficiency, of these methods is still in doubt. Therefore,
the search for new [LB] models, overcoming these limita-
tions, remains an important unsolved task.
1b - Goals of the investigation
The aim of this work is the development of an inverse
kinetic theory for the incompressible Navier-Stokes equa-
tions (INSE) which, besides realizing an exact Navier-
Stokes (and Poisson) solver, overcomes some of the lim-
itations of previous LBM’s. Unlike Refs. [2, 3], where a
continuous IKT was considered, here we construct a dis-
crete theory based on the LB velocity-space discretiza-
tion. In such a type of approach, the kinetic description
is realized by a finite number of discrete distribution func-
tions fi(r, t), for i = 0, k, each associated to a prescribed
discrete constant velocity ai and defined everywhere in
the existence domain of the fluid fields (the open set Ω×I
). The configuration space Ω is a bounded subset of the
Euclidean space R3and the time interval I is a subset of
R. The kinetic theory is obtained as in [2, 3] by introduc-
ing an inverse kinetic equation (LB-IKE) which advances
in time the distribution function and by properly defin-
ing a correspondence principle, relating a set of velocity
momenta with the relevant fluid fields.
To achieve an IKT for INSE, however, also a proper
treatment of the initial and boundary conditions, to be
satisfied by the kinetic distribution function, must be in-
cluded. In both cases, it is proven that they can be de-
fined to be exactly consistent - at the same time - both
with the hydrodynamic equations (which must hold also
arbitrarily close to the boundary of the fluid domain) and
with the prescription of the initial and Dirichlet bound-
ary conditions set for the fluid fields. Remarkably, both
the choice of the initial and equilibrium kinetic distri-
bution functions and their functional class remain essen-
tially arbitrary. In other words, provided suitable min-
imal smoothness conditions are met by the kinetic dis-
tributions function, for arbitrary initial and boundary
kinetic distribution functions, the relevant moment equa-
tions of the kinetic equation coincide identically with the
relevant fluid equations. This includes the possibility
of defining a LB-IKT in which the kinetic distribution
function is not necessarily a Galilean invariant.
This arbitrariness is reflected also in the choice of pos-
sible ”equilibrium” distribution functions, which remain
essentially free in our theory, and can be made for exam-
ple in order to achieve minimal algorithmic complexity.
A possible solution corresponds to assume polynomial-
type kinetic equilibria, as in the traditional asymptotic
LBM’s. These kinetic equilibria are well-known to be
non-Galilean invariant with respect to arbitrary finite
velocity translations. Nevertheless, as discussed in detail
in Sec.4, Subsection 4A, although the adoption of Galilei
invariant kinetic distributions is in possible, this choice
does not represent an obstacle for the formulation of a
LB-IKT. Actually Galilean invariance need to be fulfilled
only by the fluid equations. The same invariance prop-
erty must be fulfilled only by the moment equations of
the LB-IKT and not necessarily by the whole LB inverse
kinetic equation (LB-IKE).
Another significant development of the theory is the
formal introduction of an entropic principle, realized by
a constant H-theorem, in order to assure the strict pos-
itivity of the kinetic distribution function in the whole
existence domain Ω× I. The present entropic principle
departs significantly from the literature. Unlike previ-
ous entropic LBM’s it is obtained without imposing any
functional constraints on the class of the initial kinetic
distribution functions. Namely without demanding the
validity of a principle of entropy maximization (PEM,
[39]) in a true functional sense on the form of the distri-
bution function. Rather, it follows imposing a constraint
only on a suitable set of extended fluid fields, in particu-
lar the kinetic pressure p1(r, t).The latter is uniquely re-
lated to the actual fluid pressure p(r, t) via the equation
p1(r, t) = p(r, t) + Po(t), with Po(t) > 0 to be denoted
as pseudo-pressure. The constant H-theorem is therefore
obtained by suitably prescribing the function Po(t) and
implies the strict positivity. The same prescription as-
sures that the entropy results maximal with respect in the
class of the admissible kinetic pressures, i.e., it satisfies a
principle of entropy maximization. Remarkably, since
this property is not affected by the particular choice of
the kinetic equilibrium, the H-theorem applies also in the
case of polynomial equilibria. We stress that the choice
of the entropy functional remains essentially arbitrary,
since no actual physical interpretation can be attached to
it. For example, without loss of generality it can always
be identified with the Gibbs-Shannon entropy. Even pre-
scribing these additional properties, in principle infinite
solutions exist to the problem. Hence, the freedom can
be exploited to satisfy further requirements (for example,
mathematical simplicity, minimal algorithmic complex-
ity, etc.). Different possible realizations of the theory and
comparisons with other CFD approaches are considered.
The formulation of the inverse kinetic theory is also use-
ful in order to determine the precise relationship between
the LBM’s and previous CFD schemes and in particular
to obtain possible improved asymptotic LBM’s with pre-
scribed accuracy. As an application, we intend to con-
struct asymptotic models which satisfy with prescribed
accuracy the required fluid equations [INSE] and possi-
bly extend also the range of validity of traditional LBM’s.
In particular, this permits to obtain asymptotic accuracy
estimates of customary LB approaches. The scheme of
presentation is as follows. In Sec.2 the INSE problem
is recalled and the definition of the extended fluid fields
{V, p1} is presented. In Sec. 3 the basic assumptions
of previous asymptotic LBM’s are recalled. In.Sec.4 and
5 the foundations of the new inverse kinetic theory are
laid down and the integral LB inverse kinetic theory is
presented, while in Sec. 6 the entropic theorem is proven
to hold for the kinetic distribution function for properly
defined kinetic pressure. Finally, in Sec.7 various asymp-
totic approximations are obtained for the inverse kinetic
theory and comparisons are introduce with previous LB
and CFD methods and in Sec. 8 the main conclusions
are drawn.
2 - THE INSE PROBLEM
A prerequisite for the formulation of an inverse kinetic
theory [2, 3] providing a phase-space description of a clas-
sical (or quantum) fluid is the proper identification of the
complete set of fluid equations and of the related fluid
fields. For a Newtonian incompressible fluid, referred to
an arbitrary inertial reference frame, these are provided
by the incompressible Navier-Stokes equations (INSE) for
the fluid fields {ρ,V,p}
∇ ·V = 0, (1)
NV = 0, (2)
ρ(r,t) = ρo. (3)
There are supplemented by the inequalities
p(r,t) ≥ 0, (4)
ρo > 0. (5)
Equations (1)-(3) are defined in a open connected set
Ω ⊆ R3 (defined as the subset of R3 where ρ(r,t) > 0)
with boundary δΩ, while Eqs. (4) and (5) apply on its
closure Ω. Here the notation is standard. Thus, N is the
NS operator
NV ≡ρo
V +∇p+ f − µ∇2V, (6)
with D
+V · ∇ the convective derivative, f denotes
a suitably smooth volume force density acting on the fluid
element and µ ≡ νρo > 0 is the constant fluid viscosity.
In particular we shall assume that f can be represented
in the form
f = −∇Φ(r) + f1(r,t)
where we have separated the conservative ∇Φ(r) and the
non-conservative f1 parts of the force. Equations (1)-(3)
are assumed to admit a strong solution in Ω × I, with
I ⊂ R a possibly bounded time interval. By assumption
{ρ,V,p} are continuous in the closure Ω. Hence if in Ω×I,
f is at least C(1,0)(Ω×I), it follows necessarily that {V,p}
must be at least C(2,1)(Ω × I). In the sequel we shall
impose on {V,p} the initial conditions
V(r,to) = Vo(r), (7)
p(r, to) = po(r).
Furthermore, for greater mathematical simplicity, here
we shall impose Dirichlet boundary conditions on δΩ
V(·,t)|
δΩ = VW (·,t)|δΩ
p(·,t)|
δΩ = pW (·,t)|δΩ .
Eqs.(3) and (7)-(8) define the initial-boundary value
problem associated to the reduced INSE (reduced INSE
problem). It is important to stress that the previous
problem can also formulated in an equivalent way by re-
placing the fluid pressure p(r, t) with a function p1(r, t)
(denoted kinetic pressure) of the form
p1(r, t) = Po + p(r, t), (9)
where Po = Po(t) is prescribed (but arbitrary) real func-
tion of time and is at least Po(t) ∈ C
(1)(I). {V,p1} will
be denoted hereon as extended fluid fields and Po(t) will
be denoted as pseudo-pressure.
3 - ASYMPTOTIC LBM’S
3A - Basic assumptions
As is well known, all LB methods are based on a dis-
crete kinetic theory, using a so-called lattice Boltzmann
velocity discretization of phase-space (LB discretization).
This involves the definition of a kinetic distribution func-
tion f, which can only take the values belonging to a
finite discrete set {fi(r, t), i = 0, k} (discrete kinetic dis-
tribution functions). In particular, it is assumed that the
functions fi, for i = 0, k, are associated to a discrete set
of k+1 different ”velocities” {ai, i = 0, k} . Each ai is an
’a priori’ prescribed constant vector spanning the vector
space Rn (with n = 2 or 3 respectively for the treatment
of two- and three-dimensional fluid dynamics),and each
fi(r, t) is represented by a suitably smooth real function
which is defined and continuous in Ω×I and in particular
is at least C(k,j)(Ω× I) with k ≥ 3.
The crucial aspect which characterizes customary LB
approaches [8, 9, 10, 11, 12, 13, 14, 17, 40, 41] involves the
construction of kinetic models which allow a finite sound
speed in the fluid and hence are based on the assumption
of a (weak) compressibility of the same fluid. This is
realized by assuming that the evolution equation (kinetic
equation) for the discrete distributions fi(r, t) (i = 1, k),
depends at least one (or more) infinitesimal (asymptotic)
parameters (see below). Such approaches are therefore
denoted as asymptotic LBM’s. They are characterized
by a suitable set of assumptions, which typically include:
1. LB assumption #1: discrete kinetic equation and
correspondence principle: the first assumption con-
cerns the definition of an appropriate evolution
equation for each fi(r, t) which must hold (together
with all its moment equations) in the whole open
set Ω× I. In customary LB approaches it takes the
form of the so-called LB-BGK equation [13, 41, 42]
L(i)fi = Ωi(fi), (10)
where i = 0, k. Here L(i) is a suitable streaming
operator,
Ωi(fi) = −νc(fi − f
i ) (11)
(with νc ≥ 0 a constant collision frequency) is
known as BKG collision operator (after Bhatba-
gar, Gross and Krook [43]) and f
i is an ”equi-
librium” distribution to be suitably defined. In
customary LBM’s it is implicitly assumed that
the solution of Eq.(10), subject to suitable initial
and boundary conditions exists and is unique in
the functional class indicated above. In partic-
ular, usually L(i) is either identified with the fi-
nite difference streaming operator (see for example
[8, 11, 13, 42]), i.e., L(i)fi(r, t) = LFD(i)fi(r, t) ≡
[fi(r+ ai∆t, t+∆t)− fi(r, t)] or with the dif-
ferential streaming operator (see for instance [17,
40, 41])
L(i) = LD(i) ≡
+ ai ·
. (12)
Here the notation is standard. In particular, in the
case of the operator LFD(i), ∆t and c∆t ≡ Lo are
appropriate parameters which define respectively
the characteristic time- and length- scales associ-
ated to the LBM time and spatial discretizations.
A common element to all LBM’s is the assump-
tion that all relevant fluid fields can be identified,
at least in some approximate sense, with appro-
priate momenta of the discrete kinetic distribu-
tion function (correspondence principle). In par-
ticular, for neutral and isothermal incompressible
fluids, for which the fluid fields are provided re-
spectively by the velocity and pressure fluid fields
{Yj(r, t), j = 1, 4} ≡ {V(r, t), p(r, t)} , it is as-
sumed that they are identified with a suitable set
of discrete velocity momenta (for j = 1, 4)
Yj(r, t) =
i=0,k
Xji(r, t)fi(r, t), (13)
where Xji(r, t) (with i = 0, k and j = 1, k) are ap-
propriate, smooth real weight functions. In the
literature several examples of correspondence prin-
ciples are provided, a particular case being provided
by the so-called D2Q9 (V, p)-scheme [44, 45]
p(r, t) = c2
i=0,k
fi = c
i=0,k
i , (14)
V(r,t) =
i=1,k
aifi =
i=1,k
i , (15)
where k = 8 and c = min {|ai| > 0, i = 0, k} is a
characteristic parameter of the kinetic model to be
interpreted as test particle velocity. In customary
LBM’s the parameter cs =
(with D the dimen-
sion of the set Ω) is interpreted as sound speed of
the fluid. In order that the momenta (14) and (15)
recover (in some suitable approximate sense) INSE
, however, appropriate subsidiary conditions must
be met.
2. LB assumption #2: Constraints and asymptotic
conditions: these are based on the introduc-
tion of a dimensionless parameter ε, to be consid-
ered infinitesimal, in terms of which all relevant
parameters can be ordered. In particular, it is
required that the following asymptotic orderings
[17, 40, 41] apply respectively to the fluid fields
ρo,V(r, t), p(r, t), the kinematic viscosity ν = µ/ρo
and Reynolds number Re = LV/ν:
ρo,V(r, t), p(r, t) ∼ o(ε
0), (16)
[1 + o(ε)] ∼ o(εαR), (17)
Re ∼ 1/o(ε
αR), (18)
where αR ≥ 0. Here we stress that the position
for ν holds in the case of D2Q9 only, while the
generalization to 3D and other LB discretizations.
is straightforward. Furthermore, the velocity c and
collision frequency νc are ordered so that
c ∼ 1/o(εαc), (19)
νc ∼ 1/o(ε
αν ), (20)
∼ o(εα), (21)
with α ≡ αν−αc > 0; the characteristic length and
time scales, Lo ≡ c∆t and ∆t for the spatial and
time discretization are assumed to scale as
∼ o(εαL), (22)
∼ o(εαt), (23)
with αt, αL > 0. Here L and T are the (smallest)
characteristic length and time scales, respectively
for spatial and time variations of V(r, t) and p(r.t).
Imposing also that 1
results infinitesimal at least
of order
∼ o(εα)
it follows that it must be also αt − αL > 0. These
assumptions imply necessarily that the dimension-
less parameter M eff ≡ V
(Mach number) must be
ordered as
M eff ∼ O(εαc) (24)
(small Mach-number expansion).
3. LB assumption #3: Chapman-Enskog expansion -
Kinetic initial conditions, relaxation conditions: it
is assumed that the kinetic distribution function
fi(r, t) admits a convergent Chapman-Enskog ex-
pansion of the form
fi = f
i + δf
i + δ
i + .., (25)
where δ ≡ εα and the functions f
i (j ∈ N)
are assumed smooth functions of the form (multi-
scale expansion) f
i (ro, r1, r2, ..to, t1, t2, ..), where
rn = δ
nr, tn = δ
nt and n ∈ N. In typical LBM’s
the parameter δ is usually identified with ε (which
requires letting α = 1), while the Chapman-Enskog
expansion is usually required to hold at least up to
order o(δ2). In addition the initial conditions
fi(r, to) = f
i (r, to), (26)
(for i = 0, k) are imposed in the closure of the fluid
domain Ω. It is well known [46] that this position
generally (i.e., for non-stationary fluid fields), im-
plies the violation of the Chapman-Enskog expan-
sion close to t = to, since the approximate fluid
equations are recovered only letting δf
0, i.e., assuming that the kinetic distribution func-
tion has relaxed to the Chapman-Enskog form (25).
This implies a numerical error (in the evaluation of
the correct fluid fields) which can be overcome only
discarding the first few time steps in the numerical
simulation.
4. LB assumption #5: Equilibrium kinetic distribu-
tion: a possible realization for the equilibrium dis-
tributions f
i (i = 0, k) is given by a polynomial of
second degree in the fluid velocity [44]
i (r, t) = wi
[p− Φ(r)] + (27)
+wiρo
ai ·V
ai ·V
Here, without loss of generality, the case of the
D2Q9 LB discretization will be considered, with wi
and ai (for i = 0, 8) denoting prescribed dimension-
less constant weights and discrete velocities. Notice
that, by definition, f
i is not a Galilei scalar. Nev-
ertheless, it can be considered approximately in-
variant, at least with respect to low-velocity trans-
lations which do not violate the low-Mach number
assumption (24).
5. LB assumption #6: Kinetic boundary conditions:
They are specified by suitably prescribing the form
of the incoming distribution function at the bound-
ary δΩ. [47, 48, 49, 50, 51, 51, 52, 53, 54, 54, 55, 56,
57, 58, 59]. However, this position is not generally
consistent with the Chapman-Enskog solution (25)
(see related discussion in Appendix A). As a con-
sequence violations of the hydrodynamic equations
may be expected sufficiently close to the boundary,
a fact which may be only alleviated (but not com-
pletely eliminated) by adopting suitable grid refine-
ments near the boundary. An additional potential
difficulty is related to the condition of strict posi-
tivity of the kinetic distribution function [57] which
is not easily incorporated into the no-slip boundary
conditions [50, 51, 52].
3B - Computational complexity of asymptotic
LBM’s
The requirements posed by the validity of these hy-
potheses may strongly influence the computational com-
plexity of asymptotic LBM’s which is usually associated
to the total number of ”logical” operations which must
be performed during a prescribed time interval. There-
fore, a critical parameter of numerical simulation meth-
ods is their discretization time scale ∆t. This is - in turn
- related to the Courant number NC =
, where V
and Lo.denote respectively the sup of the magnitude of
the fluid velocity and the amplitudes of the spatial dis-
cretization. As is well known ”optimal” CFD simulation
methods typically allow Lo ∼ L and a definition of the
time step ∆t = ∆tOpt such that NC ∼
V ∆tOpt
∼ 1. In-
stead, for usual LBM’s satisfying the low-M eff assump-
tion (24), the Courant number is very small since it re-
sults NC = M
eff Lo
∼ O(εα)Lo
. This means that their
discretization time scale of ∆t is much smaller than ∆tOpt
and reads
∆t ∼M eff
∆tOpt. (28)
In addition, depending on the accuracy of the numeri-
cal algorithms adopted for the construction of the dis-
crete kinetic distribution function, also the ratio Lo
sults infinitesimal in the sense Lo
∼ o(εαL), with suitable
αL > 0. Finally, we stress that LB approaches based
on the adoption of the finite-difference streaming opera-
tor LFD(i) are usually only accurate to order o(∆t
2). For
them, therefore, the requirement placed by Eq.(28) might
be even stronger. This implies that traditional LBM’s
may involve a vastly larger computation time than that
afforded by more efficient numerical methods.
4 - NEW LB INVERSE KINETIC THEORY
(LB-IKT)
A basic issue in LB approaches [8, 11, 13, 42] con-
cerns the choice of the functional class of the discrete
kinetic distribution functions fi (i = 0, k) as well as the
related definition of the equilibrium discrete distribution
function f
i [which appears in the BGK collision opera-
tor; see Eq.(11)]. This refers in particular to their trans-
formation properties with respect to arbitrary Galilean
transformations, and specifically to their Galilei invari-
ance with respect to velocity translations with constant
velocity.
In statistical mechanics it is well known that the ki-
netic distribution function is usually assumed to be a
Galilean scalar. The same assumption can, in principle,
be adopted also for LB models. However, the kinetic
distribution functions fi and f
i do not necessarily re-
quire a physical interpretation of this type. In the se-
quel we show that for a discrete inverse kinetic theory it
is sufficient that fi and f
i be so defined that the mo-
ment equations coincide with the fluid equations (which
by definition are Galilei covariant). It is sufficient to de-
mand that both fi and f
i are identified with a ordinary
scalars with respect to the group of rotation in R2, while
they need not be necessarily invariant with respect to
arbitrary velocity translations. This means that fi is in-
variant only for a particular subset of inertial reference
frames. For example for a fluid which at the initial time
moves locally with constant velocity an element of this
set can be identified with the inertial frame which in the
same position is locally co-moving with the fluid.
The adoption of non-translationally invariant discrete
distributions fi is actually already well known in LBM
and results convenient for its simplicity. This means,
manifestly, that in general no obvious physical interpre-
tation can be attached to the other momenta of the dis-
crete kinetic distribution function. As a consequence,
the very definition of the concept of statistical entropy to
be associated to the f ′is is essentially arbitrary, as well as
the related principle of entropy maximization, typically
used for the determination of the equilibrium distribution
function f
i . Several authors, nevertheless, have investi-
gated the adoption of possible alternative formulations,
which are based on suitable definitions of the entropy
functional and/or the requirement of approximate or ex-
act Galilei invariance (see for example [29, 32, 62]).
4A - Foundations of LB-IKT
As previously indicated, there are several important
motivations for seeking an exact solver based on LBM.
The lack of a theory of this type represents in fact a
weak point of LB theory. Besides being a still unsolved
theoretical issue, the problem is relevant in order to de-
termine the exact relationship between the LBM’s and
traditional CFD schemes based on the direct discretiza-
tion of the Navier–Stokes equations. Following ideas re-
cently developed [2, 3, 4, 5, 7], we show that such a theory
can be formulated by means of an inverse kinetic theory
(IKT) with discrete velocities. By definition such an IKT
should yield exactly the complete set of fluid equations
and which, contrary to customary kinetic approaches in
CFD (in particular LB methods), should not depend on
asymptotic parameters. This implies that the inverse ki-
netic theory must also satisfy an exact closure condition.
As a further condition, we require that the fluid equa-
tions are fulfilled independently of the initial conditions
for the kinetic distribution function (to be properly set)
and should hold for arbitrary fluid fields. The latter re-
quirement is necessary since we must expect that the
validity of the inverse kinetic theory should not be lim-
ited to a subset of possible fluid motions nor depend on
special assumptions, like a prescribed range of Reynolds
numbers. In principle a phase-space theory, yielding an
inverse kinetic theory, may be conveniently set in terms of
a quasi-probability, denoted as kinetic distribution func-
tion, f(x, t). A particular case of interest (investigated in
Refs.[2, 3]) refers to the case in which f(x, t) can actu-
ally be identified with a phase-space probability density.
In the sequel we address both cases, showing that, to a
certain extent, in both cases the formulation of a generic
IKT can actually be treated in a similar fashion. This
requires the introduction of an appropriate set of consti-
tutive assumptions (or axioms). These concern in par-
ticular the definitions of the kinetic equation - denoted
as inverse kinetic equation (IKE) - which advances in
time f(x, t) and of the velocity momenta to be identified
with the relevant fluid fields (correspondence principle).
However, further assumptions, such as those involving
the regularity conditions for f(x, t) and the prescription
of its initial and boundary conditions must clearly be
added. The concept [of IKT] can be easily extended to
the case in which the kinetic distribution function takes
on only discrete values in velocity space. In the sequel
we consider for definiteness the case of the so-called LB
discretization, whereby - for each (r, t) ∈ Ω × I - the
kinetic distribution function is discrete, and in particu-
lar admits a finite set of discrete values fi(r, t) ∈ R, for
i = 0, k, each one corresponding to a prescribed constant
discrete velocity ai ∈ R
3 for i = 0, k.
4B - Constitutive assumptions
Let us now introduce the constitutive assumptions (ax-
ioms) set for the construction of a LB-IKT for INSE,
whose form is suggested by the analogous continuous
inverse kinetic theory [2, 3]. The axioms, define the
”generic” form of the discrete kinetic equation, its func-
tional setting, the momenta of the kinetic distribution
function and their initial and boundary conditions, are
the following ones:
Axiom I - LB–IKE and functional setting.
Let us require that the extended fluid fields {V,p1}
are strong solutions of INSE, with initial and boundary
conditions (7)-(8) and that the pseudo pressure po(t) is
an arbitrary, suitably smooth, real function. In particu-
lar we impose that the fluid fields and the volume force
belong to the minimal functional setting:
p1,ΦǫC
(2,1)(Ω× I),
VǫC(3,1)(Ω× I), (29)
(1,0)(Ω× I).
We assume that in the set Ω×I the following equation
LD(i)fi = Ωi(fi) + Si (30)
[LB inverse kinetic equation (LB-IKE)] is satisfied iden-
tically by the discrete kinetic distributions fi(r, t) for
i = 0, k. Here Ωi(fi) and LD(i) are respectively the BGK
and the differential streaming and operators [Eqs.(11)
and (12)], while Si is a source term to be defined. We
require that KB-IKE is defined in the set Ω× I, so that
Ωi(fi) and Si are at least that C
(1)(Ω × I) and contin-
uous in Ω × I. Moreover Ωi(fi), defined by Eq.(11), is
considered for generality and will be useful for compar-
isons with customary LB approaches. We remark that
the choice of the equilibrium kinetic distribution f
the BGK operator remains completely arbitrary. We
assume furthermore that in terms of fi the fluid fields
{V, p1} are determined by means of functionals of the
form MXj [fi] =
i=0,8
Xjfi (denoted as discrete velocity
momenta). For X = X1, X2 (with X1 = c
2, X2 =
these are related to the fluid fields by means of the equa-
tions (correspondence principle)
p1(r, t)− Φ(r) = c
i=0,8
fi = c
i=0,8
i , (31)
V(r,t)=
i=1,8
aifi =
i=1,8
i , (32)
where c = min {|ai| , i = 1, 8} is the test particle veloc-
ity and f
i is defined by Eq.(27) but with the kinetic
pressure p1 that replaces the fluid pressure p adopted
previously [44]. These equations are assumed to hold
identically in the set Ω × I and by assumption, fi and
i belong to the same functional class of real functions
defined so that the extended fluid fields belong to the
minimal functional setting (29). Moreover, without loss
of generality, we consider the D2Q9 LB discretization.
Axiom II - Kinetic initial and boundary conditions.
The discrete kinetic distribution function satisfies, for
i = 0, k and for all r belonging to the closure Ω, the
initial conditions
fi(r, to) = foi(r,to) (33)
where foi(r,to) (for i = 0, k) is a initial distribution func-
tion defined in such a way to satisfy in the same set the
initial conditions for the fluid fields
p1o(r) ≡ Po(to) + po(r)− Φ(r) = (34)
i=0,8
foi(r),
Vo(r) =
i=1,8
aifoi(r) . (35)
To define the analogous kinetic boundary conditions on
δΩ, let us assume that δΩ is a smooth, possibly moving,
surface. Let us introduce the velocity of the point of the
boundary determined by the position vector rw ∈ δΩ, de-
fined by Vw(rw(t), t) =
rw(t) and denote by n(rw, t)
the outward normal unit vector, orthogonal to the bound-
ary δΩ at the point rw. Let us denote by f
i (rw, t)
and f
i (rw , t) the kinetic distributions which carry the
discrete velocities ai for which there results respectively
(ai −Vw) ·n(rw , t) > 0 (outgoing-velocity distributions)
and (ai −Vw) · n(rw, t) ≤ 0 (incoming-velocity distribu-
tions) and which are identically zero otherwise. We as-
sume for definiteness that both sets, for which |ai| > 0,
are non empty (which requires that the parameter c
be suitably defined so that c > |Vw|). The bound-
ary conditions are obtained by prescribing the incom-
ing kinetic distribution f
i (rw , t), i.e., imposing (for all
(rw, t) ∈ δΩ× I)
i (rw, t) = f
oi (rw , t). (36)
Here f
oi (rw, t) are suitable functions, to be assumed
non-vanishing and defined only for incoming discrete ve-
locities for which (ai −Vw)·n(rw , t) ≤ 0. Manifestly, the
functions f
oi (rw, t) (i = 0, k) must be defined so that
the Dirichlet boundary conditions for the fluid fields are
identically fulfilled, namely there results
p1w(rw, t) = Po(t) + pw(rw, t)− Φ(r) = (37)
i=0,k
oi (rw, t) + f
i (rw, t)
Vw(rw, t) = (38)
i=1,k
oi (rw, t) + f
i (rw, t)
Here, again, the functions foi(r) and f
oi (rw, t) (for i =
0, k) must be assumed suitably smooth. A particular case
is obtained imposing identically for i = 0, k
foi(r,to) = f
i (r, to), (39)
oi (rw, t) = f
i (rw , t), (40)
where the identification with f
oi (rw, t) and f
oi (rw, t)
is intended respectively in the subsets ai ·n(rw, t) > 0 and
ai ·n(rw , t) ≤ 0. Finally, we notice that in case Neumann
boundary conditions are imposed on the fluid pressure,
Eq.(37) still holds provided pw(rw, t) is intended as a
calculated value.
Axiom III - Moment equations.
If fi(r, t), for i = 0, k, are arbitrary solutions of LB-
IKE [Eq.(30)] which satisfy Axioms I and II validity of
Axioms I and II, we assume that the moment equations of
the same LB-IKE, evaluated in terms of the moment op-
erators MXj [·] =
i=0,8
Xj ·, with j = 1, 2, coincide iden-
tically with INSE, namely that there results identically
[for all (r, t) ∈ Ω× I]
MX1 [Lifi − Ωi(fi)− Si] = ∇ ·V = 0, (41)
MX2 [Lifi − Ωi(fi)− Si] = NV = 0. (42)
Axiom IV - Source term.
The source term is required to depend on a finite num-
ber of momenta of the distribution function. It is as-
sumed that these include, at most, the extended fluid
fields {V,p1} and the kinetic tensor pressure
Π = 3
fiaiai − ρoVV. (43)
• Furthermore, we also normally require (except
for the LB-IKT described in Appendix B) that
Si(r, t) results independent of f
i (r,t), foi(r) and
fwi(rw , t) (for i = 0, k).
Although, the implications will made clear in the fol-
lowing sections, it is manifest that these axioms do not
specify uniquely the form (and functional class) of the
equilibrium kinetic distribution function f
i (r,t), nor
of the initial and boundary kinetic distribution func-
tions (33),(36). Thus, both f
i (r,t), foi(r,to) and the
related distribution they still remain in principle com-
pletely arbitrary. Nevertheless, by construction, the
initial and (Dirichlet) boundary conditions for the fluid
fields are satisfied identically. In the sequel we show that
these axioms define a (non-empty) family of parameter-
dependent LB-IKT’s, depending on two constant free pa-
rameters νc, c > 0 and one arbitrary real function Po(t).
The examples considered are reported respectively in the
following Sec. 5,6 and in the Appendix B.
5 - A POSSIBLE REALIZATION: THE
INTEGRAL LB-IKT
We now show that, for arbitrary choices of the distri-
butions fi(r,t) and f
i (r,t) which fulfill axioms I-IV, an
explicit (and non-unique) realization of the LB-IKT can
actually be obtained. We prove, in particular, that a pos-
sible realization of the discrete inverse kinetic theory, to
be denoted as integral LB-IKT, is provided by the source
Si = (44)
− ai ·
f1−µ∇
V −∇ ·Π+∇p
≡ S̃i,
where wi
is denoted as first pressure term. Holds, in
fact, the following theorem.
Theorem 1 - Integral LB-IKT
In validity of axioms I-IV the following statements
hold. For an arbitrary particular solution fi and for ar-
bitrary extended fluid fields :
A) if fi is a solution of LB-IKE [Eq.(30)] the moment
equations coincide identically with INSE in the set Ω×I;
B) the initial conditions and the (Dirichlet) boundary
conditions for the fluid fields are satisfied identically;
C) in validity of axiom IV the source term S̃i is non-
uniquely defined by Eq.(44).
Proof
A) We notice that by definition there results identically
S̃i =
aiS̃i = (46)
f−µ∇2V−∇ ·Π+∇p
On the other hand, by construction (Axiom I) fi (i =
1, k) is defined so that there results identically
i=0 Ωi =
0 and
i=0 aiΩi = 0. Hence the momenta MX1 ,MX2 of
LB-IKE deliver respectively
i=1,8
aifi = 0 (47)
i=1,8
aifi + ρoV · ∇V +∇p1 + f−µ∇
V = 0 (48)
where the fluid fields V,p1 are defined by Eqs.(31),(32).
Hence Eqs.(47) and (48) coincide respectively with the
isochoricity and Navier-Stokes equations [(1) and (2)].
As a consequence, fi is a particular solution of LB-IKE
iff the fluid fields {V,p1} are strong solutions of INSE.
B) Initial and boundary conditions for the fluid fields
are satisfied identically by construction thanks to Axiom
C) However, even prescribing νc, c > 0 and the real
function Po(t), the functional form of the equation can-
not be unique The non uniqueness of the functional form
of the source term S̃i(r, t) is assumed to be indepen-
dent of f
i (r,t) [and hence of Eq.(30)] is obvious. In
fact, let us assume that S̃i is a particular solution for
the source term which satisfies the previous axioms I-
IV. Then, it is always possible to add to Si arbitrary
terms of the form S̃i + δSi, with δSi 6= 0 which depends
only on the momenta indicated above, and gives van-
ishing contributions to the first two moment equations,
namely MXj [δSi] =
i=0,8
XjδSi = 0, with j = 1, 2. To
prove the non-uniqueness of the source term Si, it is suf-
ficient to notice that, for example, any term of the form
δSi =
F (r, t), with F (r, t) an arbitrary real
function (to be assumed, thanks to Axiom IV, a linear
function of the fluid velocity), gives vanishing contribu-
tions to the momentaMX1 ,MX2 . Hence S̃i is non-unique.
The implications of the theorem are straightforward.
First, manifestly, it holds also in the case in which the
BGK operator vanishes identically. This occurs letting
νc = 0 in the whole domain Ω × I. Hence the inverse
kinetic equation holds independently of the specific defi-
nition of f
i (r,t).
An interesting feature of the present approach lies in
the choice of the boundary condition adopted for fi(r,t),
which is different from that usually adopted in LBM’s
[see for example [14] for a review on the subject]. In par-
ticular, the choice adopted is the simplest permitting to
fulfill the Dirichlet boundary conditions [imposed on the
fluid fields]. This is obtained prescribing the functional
form of fi(r,t) on the boundary of the fluid domain (δΩ),
which is identified with a function foi(r, t).
Second, the functional class of fi(r,t), f
i (r,t) and of
foi(r, t) remains essentially arbitrary. Thus, in particu-
lar, the initial and boundary conditions, specified by the
same function foi(r, t), can be defined imposing the po-
sitions (39),(40). As further basic consequence, f
i (r,t)
and fi(r,t) need not necessarily be Galilei-invariant (in
particular they may not be invariant with respect to ve-
locity translations), although the fluid equations must
be necessarily fully Galilei-covariant. As a consequence
it is always possible to select f
i (r,t) and foi(r, t) based
on convenience and mathematical simplicity. Thus, be-
sides distributions which are Galilei invariant and sat-
isfy a principle of maximum entropy (see for example
[22, 30, 32, 34, 60, 61]), it is always possible to iden-
tify them [i.e., f
i (r,t), foi(r, t)] with a non-Galilean in-
variant polynomial distribution of the type (27) [mani-
festly, to be exactly Galilei-invariant each f
i (r,t) should
depend on velocity only via the relative velocity ui =
ai −V].
We mention that the non-uniqueness of the source term
S̃i can be exploited also by imposing that f
i (r,t) re-
sults a particular solution of the inverse kinetic equation
Eq.(30) and there results also foi(r, t) = f
i (r,t). In Ap-
pendix B we report the extension of THM.1 which is ob-
tained by identifying again f
i (r,t) with the polynomial
distribution (27).
6 - THE ENTROPIC PRINCIPLE - CONDITION
OF POSITIVITY OF THE KINETIC
DISTRIBUTION FUNCTION
A fundamental limitation of the standard LB ap-
proaches is their difficulty to attain low viscosities, due to
the appearance of numerical instabilities [14]. In numeri-
cal simulations based on customary LB approaches large
Reynolds numbers is usually achieved by increasing nu-
merical accuracy, in particular strongly reducing the time
step and the grid size of the spatial discretization (both
of which can be realized by means of numerical schemes
with adaptive time-step and using grid refinements).
Hence, the control [and possible inhibition] of numerical
instabilities is achieved at the expense of computational
efficiency. This obstacle is only partially alleviated by
approaches based on ELBM [22, 30, 32, 34, 60, 61]. Such
methods are based on the hypothesis of fulfilling an H-
theorem, i.e., of satisfying in the whole domain Ω × I
the condition of strict positivity for the discrete kinetic
distribution functions. This requirement is considered,
by several authors (see for example [26, 29, 62]), an es-
sential prerequisite to achieve numerical stability in LB
simulations. However, the numerical implementation of
ELBM typically induce a substantial complication of the
original algorithm, or require a cumbersome fine-tuning
of adjustable parameters [22, 37].
6A - The constant entropy principle and PEM
A basic aspect of the IKT’s here developed is the possi-
bility of fulfilling identically the strict positivity require-
ment by means of a suitable H-theorem which provides
also a maximum entropy principle. In particular, in this
Section, extending the results of THM.1 and 2, we intend
to prove that a constant H-theorem can be established
both for the integral and differential LB-IKT’s defined
above. The H-theorem can be reached by imposing for
the Gibbs-Shannon entropy functional the requirement
that for all t ∈ I there results
S(f) = −
i=0,8
fi ln(fi/wi) = 0, (49)
which implies that S(f) is necessarily maximal in a suit-
able functional set {f} . The result can be stated as fol-
lows:
Theorem 2 - Constant H-theorem
In validity of THM.1, let us assume that:
1) the configuration domain Ω is bounded;
2) at time to the discrete kinetic distribution functions
fi, for i = 0, 8, are all strictly positive in the set Ω.
Then the following statements hold:
A) by suitable definition of the pseudo pressure
Po(t), the Gibbs-Shannon entropy functional S(f) =
i=0,8
fi ln(fi/wi) can be set to be constant in the
whole time interval I. This holds provided the pseudo-
pressure Po(t) satisfies the differential equation
(1 + log fi) = (50)
ai · ∇fi − Ŝi
(1 + log fi) ,
where Ŝi = Si +
B) if the entropy functional S(f) =
i=0,8
fi ln(fi/wi) is constant in the whole
time interval I the discrete kinetic distribution functions
fi are all strictly positive in the whole set Ω× I;
C) an arbitrary solution of LB-IKE [Eq.(30)] which
satisfies the requirement A) is extremal in a suitable func-
tional class and maximizes the Gibbs-Shannon entropy .
Proof:
A) Invoking Eq.(30), there results
∂S(t)
[1 + log fi] = (51)
(ai · ∇fi − Si) (1 + log fi) ,
where Si is the source term, provided by Eq.(44). By
direct substitution it follows the thesis.
B) If Eq.(50) holds identically in there results ∀t ∈
I, S (t) = S (t0) , which implies the strict positivity of fi,
for all i = 0, 8.
C) Let us introduce the functional class
{f + αδf} = {fi = fi(t) + αδfi(t), i = 0, 8} , (52)
where α is a finite real parameter and the syn-
chronous variation δfi(t) is defined δfi(t) = dfi(t) ≡
∂fi(t)
dt. Introducing the synchronous variation of the en-
tropy, defined by δS (t) = ∂
, with ψ(α) =
S (f + αδf) , it follows
δS (t) = dt
∂S(t)
. (53)
Since in validity of Eq.(50) there results
∂S(t)
which in view of Eq.(53) implies also δS (t) = 0. It is im-
mediately follows that there results necessarily δ2S (t) ≤
0, i.e., S (t) is maximal. Therefore, the kinetic distribu-
tion function which satisfies IKE (Eq.(30)] is extremal in
the functional class of variations (52) and maximizes the
Gibbs-Shannon entropy functional.
6B - Implications
In view of statement B, THM.2 warrants the strict pos-
itivity of the discrete distribution functions fi (i = 0, 8)
only in the open set Ω × I, while nothing can be said
regarding their behavior on the boundary δΩ (on which
fi might locally vanish). However, since the inverse ki-
netic equation actually holds only in the open set Ω× I,
this does not affect the validity of the result. While the
precise cause of the numerical instability of LBM’s is still
unknown,the strict positivity of the distribution function
is usually considered important for the stability of the nu-
merical solution [29, 30]. It must be stressed that the nu-
merical implementation of the condition of constant en-
tropy Eq.(50) should be straightforward, without involv-
ing a significant computational overhead for LB simula-
tions. Therefore it might represent a convenient scheme
to be adopted also for customary LB methods.
7 - ASYMPTOTIC APPROXIMATIONS AND
COMPARISONS WITH PREVIOUS CFD
METHODS
A basic issue is the relationship with previous CFD nu-
merical methods, particularly asymptotic LBM’s. Here
we consider, for definiteness, only the case of the inte-
gral LB-IKT introduced in Sec.5. Another motivation is
the possibility of constructing new improved asymptotic
models, which satisfy with prescribed accuracy the re-
quired fluid equations [INSE], of extending the range of
validity of traditional LBM’s and fulfilling also the en-
tropic principle (see Sec.6). The analysis is useful in
particular to establish on rigorous grounds the consis-
tency of previous LBM’s. The connection [with previ-
ous LBM’s] can be reached by introducing appropriate
asymptotic approximations for the IKT’s, obtained by
assuming that suitable parameters which characterize the
IKT’s are infinitesimal (or infinite) (asymptotic parame-
ters). A further interesting feature is the possibility of
constructing in principle a class of new asymptotic LBM’s
with prescribed accuracy , i.e., in which the distribution
function (and the corresponding momenta) can be de-
termined with predetermined accuracy in terms of per-
turbative expansions in the relevant asymptotic parame-
ters. Besides recovering the traditional low-Mach number
LBM’s [17, 21, 40], which satisfy the isochoricity condi-
tion only in an asymptotic sense and are closely related to
the Chorin artificial compressibility method, it is possible
to obtain an improved asymptotic LBM’s which satisfy
exactly the same equation.
We first notice that the present IKT is characterized
by the arbitrary positive parameters νc, c and the initial
value Po(to), which enter respectively in the definition
of the BGK operator [see (11)], the velocity momenta
and equilibrium distribution function f
i . Both c and
Po(to) must be assumed strictly positive, while, to assure
the validity of THM.2, Po(to) must be defined so that
(for all i = 0, 8) f
i (r,to) > 0 in the closure Ω. Thanks
to THM.1.and 2 the new theory is manifestly valid for
arbitrary finite value of these parameters. This means
that they hold also assuming
o(εαν )
, (54)
o(εαc)
, (55)
Po(to) ∼ o(ε
0), (56)
where ε denotes a strictly positive real infinitesimal,
αν , αc > 0 are real parameters to be defined, while the
extended fluid fields {ρ,V, p1} and the volume force f
are all assumed independent of ε. Hence, with respect to
ε they scale
ρo,V,p1, f ∼ o(ε
0). (57)
As a result, for suitably smooth fluid fields (i.e., in va-
lidity of Axiom 1) and appropriate initial conditions for
fi(r, t), it is expected that the first requirement actually
implies in the whole set Ω× I the condition of closeness
fi(r, t) ∼= f
i (r, t) [1 + o(ε)] , consistent with the LB As-
sumption #4. To display meaningful comparisons with
previous LBM’s let us introduce the further assumption
that the fluid viscosity is small in the sense
µ ∼ o(εαµ), (58)
with αµ ≥ 1 another real parameter to be defined.
Asymptotic approximations for the corresponding LB-
IKE [Eq.(30)] can be directly recovered by introducing
appropriate asymptotic orderings for the contributions
appearing in the source term Si = S̃i. Direct inspec-
tion shows that these are provided by the (dimensional)
parameters
M effp,a ≡
, (59)
∣∣∇ ·Π−∇p
∣∣ , (60)
∣∣µ∇2V
∣∣ . (61)
The first two M effp,a and M
are here denoted respec-
tively as (first and second) pressure effective Mach num-
bers, driven respectively by the pressure time-derivative
and by the divergence of the pressure anisotropy Π−p1.
Furthermore, M
is denoted as velocity effective Mach
number. Physically relevant examples [of asymptotic
LBM’s] can be achieved by introducing suitable orderings
in terms of the single infinitesimal ε for the parameters
M effp,a ,M
.We stress that these orderings, in
principle, can be introduced without actually introducing
restrictions on the fluid fields, i.e., retaining the assump-
tion that the extended fluid fields are independent of ε.
Interesting cases are provided by the asymptotic order-
ings indicated below.
7A - Small effective Mach numbers (Meffp,a ,M
An important aspect of LB theory is the possibility
of constructing asymptotic LBM’s with prescribed accu-
racy with respect to the infinitesimal parameter ε, in the
sense that the fluid equations are satisfied at least cor-
rect up to terms of order o(εn) included, with n = 1
or 2, namely ignoring error terms of order o(εn+1) or
higher. Let us, first, consider the case in which all pa-
rametersM effp,a ,M
and M
are all infinitesimal w.r.
to ε (low-effective-Mach numbers). Since the parameters
c and νc are free, they can be defined so that that there
results c ∼ νc ∼ 1/o(ε) [which implies αc = αν = 1].
This requires
M effp,a ∼M
∼ o(ε2). (62)
If, we consider a low-viscosity fluid for which the kine-
matic viscosity ν = µ/ρo can be assumed of order ε [and
hence αµ = 1] it follows that
∼ o(ε2). (63)
Thanks to the assumptions (54)-(58) there follows ∇ ·
Π − ∇p ∼ o(ε) and µ∇2V ∼ o(ε),which implies that
the source term S̃i, ignoring corrections of order o(ε
becomes
S̃i ∼= S̃Ai [1 + o(ε)] , (64)
S̃Ai ≡ −
ai · f . (65)
It is immediate to determine the corresponding moment
equations, which read:
+∇ ·V = 0, (66)
NV = 0+ o(ε2), (67)
Formally the first equation can be interpreted as an evo-
lution equation for the kinetic pressure p1. Nevertheless,
in view of the ordering (62) it actually implies the iso-
choricity condition
∇ ·V = 0 + o(ε2). (68)
Instead, the second one [Eq.(67)]. due to the asymp-
totic approximation (63), reduces to the Euler equation.
Therefore in this case the asymptotic approximation (64)
is not adequate. To recover the correct Navier-Stokes
equation a more accurate approximation is needed, real-
ized requiring that the hydrodynamic equations are sat-
isfied correct to order o(ε3). A fist possibility is to con-
sider a more accurate approximation for the source term.
Restoring the pressure and viscous source terms in (64)
there results the asymptotic source term
S̃Bi ≡
− ai ·
f1−µ∇
, (69)
where in validity of the previous orderings
S̃i ∼= S̃Bi [1 + o(ε)] . (70)
The corresponding moment equations become therefore
∇ ·V = 0, (71)
NV = 0+ o(ε3). (72)
It is remarkable that in this case the isochoricity condi-
tion is exactly fulfilled, even if the source term is not the
exact one. For the sake of reference, it is interesting to
mention another possible small-Mach-number ordering.
This is obtained imposing for the parameters c and νc
, (73)
o(ε2)
, (74)
while requiring for ν = µ/ρo the same constraint adopted
by asymptotic LBM’s, namely Eq.(17). In this case one
can show that the moment equation (72) is actually satis-
fied correct to order o(ε3), while the isochoricity condition
is only satisfied to order o(ε2). The following theorem
can, in fact, be proven:
Theorem 3 - Low effective-Mach-numbers asymptotic
approximation
In validity of THM.1, let us invoke the following as-
sumptions:
1) LB assumptions #3 and #4 for the discrete kinetic
distributions fi ( i = 0, 8);
2) the free parameters c and νc are assumed to satisfy
the asymptotic orderings (73),(74);
3) the fluid viscosity µ is assumed of order µ ∼ o(ε)
4) the fluid viscosity µ is prescribed so that the kine-
matic viscosity ν = µ/ρo is defined in accordance to
Eq.(17);
5) the kinetic pressure p1 is assumed slowly varying in
the sense
∂ ln p1
∼ o(ε). (75)
It follows that the source term is approximated by
Eq.(64) and moment equations are provided by the
asymptotic equations:
+∇ ·V = 0 + o(ε3), (76)
NV = 0+ o(ε3), (77)
i.e., the isochoricity and NS equation are recovered re-
spectively correct to order o(ε2) and o(ε3).
Proof
First we notice that the ordering assumptions 2)-5)
require
M effp,a ∼ o(ε
3) (78)
∼ o(ε2), (79)
∼ o(ε4), (80)
which imply at least the validity of Eqs.(64)-(67). The
proof of Eqs.(76) and (77) is immediate. In both cases
it sufficient to notice that in validity of hypotheses 1)-3)
and in terms of a Chapman-Enskog perturbative solution
of Eq.(30) there results actually
−µ∇2V −∇ ·Π+∇p = O + o(ε3), (81)
and hence S̃i reduces to Eq.(64).
The predictions of THM.3 are relevant for comparisons
and to provide asymptotic accuracy estimates for previ-
ous asymptotic LBM’s [see Refs. [17, 21, 40]]. In fact,
the asymptotic moment equations (76) and (77) formally
coincide with the analogous moment equations predicted
by such theories, when the kinetic pressure p is replaced
by the fluid pressure p1 (i.e., if the function Po(t) is set
identically equal to zero). [17, 21, 40]. Nevertheless, the
accuracy of customary LBM’s depends on the properties
of the solutions of INSE. In fact, if one assumes
∂ ln p1
∼ o(ε0) (82)
the customary (V, p) asymptotic LBM [17, 21, 40] result
actually accurate only to order o(ε2). Therefore, in such
case to reach an accuracy of order o(ε3) the approxima-
tion (69) must be invoked for the source term.
The other interesting feature of Eqs.(76) and (77) is
that they provide a connection with the artificial com-
pressibility method (ACM) postulated by Chorin [63],
previously motivated merely on the grounds of an asymp-
totic LBM [21]. In fact, these coincides with the Chorin’s
pressure relaxation equation where c can be interpreted
as sound speed of the fluid. However - in a sense - this
analogy is purely formal and is only due to the neglect of
the first pressure source term in Si. It disappears alto-
gether in Eq.(71) if we adopt the more accurate asymp-
totic source term (69). A further difference is provided
by the adoption of the kinetic pressure p1 which replaces
the fluid pressure p (used in Chorin approach). We stress
that the choice of p1 here adopted, with Po(t) determined
by the entropic principle, represents an important differ-
ence, since it permits to satisfy everywhere in Ω× I the
condition of strict positivity for the discrete kinetic dis-
tribution functions.
7B - Finite pressure-Mach number Meffp,a
Another possible asymptotic ordering, usually not per-
mitted by customary asymptotic LBM’s, is the one in
which the test particle velocity is finite, namely c ∼ o(ε0),
the viscosity remains arbitrary and is taken of order
µ ∼ o(ε0) while again νc is assumed νc ∼ 1/o(ε
2) [i.e.,
αc = νc = 0, αν = 2]. In this case the pressure Mach
M effp,a number results finite, while velocity and the sec-
ond pressure Mach numbers are considered infinitesimal,
respectively of first and second order in ε, namely
M effp,a ∼ o(ε
∼ o(ε), (83)
∼ o(ε2).
To obtain the fluid equation with the prescribed accu-
racy, say of order o(ε2), it is sufficient to approximate
the source term S̃i in terms of S̃i ∼= S
1 + o(ε2)
. The
set of asymptotic moment equations coincide therefore
with Eqs.(71),(72). Again, the isochoricity condition is
exactly fulfilled, while in this case the NS equation is
accurate only to order o(ε2).
7C - Small effective pressure-Mach numbers
(Meffp,a ,M
) and finite velocity-Mach number (M
Finally, another interesting case is the one in which the
fluid viscosity µ remains finite (strongly viscous fluid),
i.e., in the sense µ ∼ o(ε0) [i.e., αµ = 0] while both
parameters c and νc are suitably large, and respectively
scale as c ∼ 1/o(ε), νc ∼ 1/o(ε
2) [i.e., αc = 1, αν = 2].
Due to assumptions (54)-(58) one obtains ∇ ·Π−∇p ∼
o(ε2) and µ∇2V ∼ o(ε0). It follows that the effective
Mach numbers scale respectively as
∼ o(ε3) (84)
M effp,a ∼ M
∼ o(ε2),
If we impose on µ also the same constraint set by Eq.(17),
the customary asymptotic LBM’s can be invoked also in
this case. However, since the first pressure and veloc-
ity Mach numbers are only second order accurate, the
NS equation is recovered to order o(ε2) only. Never-
theless, it is possible to recover with prescribed accuracy
the fluid equations (71),(72). This is obtained adopt-
ing the source term S̃i ∼= S̃Bi [see Eq.(69)]. As a basic
consequence, the isochoricity equation is satisfied exactly
(hence no meaningful analogy with Chorin’s approach
arises), while the NS equation results correct to order
o(ε3). These results provide a meaningful extension of
the customary asymptotic LBM’s. We stress that the
entropic approach here developed holds independently of
the asymptotic orderings here considered [for the param-
etersM effp,a ,M
]. Thus it can be used in all cases
to assure the strict positivity of the discrete distribution
function.
8 - CONCLUSIONS
In this paper we have presented the theoretical foun-
dations of a new phase-space model for incompressible
isothermal fluids, based on a generalization of customary
lattice Boltzmann approaches.We have shown that many
of the limitations of traditional (asymptotic) LBM’s can
be overcome. As a main result, we have proven that
the LB-IKT can be developed in such a way that it
furnishes exact Navier-Stokes and Poisson solvers, i.e.,
it is - in a proper sense - an inverse kinetic theory for
INSE. The theory exhibits several features, in particular
we have proven that the integral LB-IKT (see Sec.5):
1. determines uniquely the fluid pressure p(r, t) via
the discrete kinetic distribution function without
solving explicitly (i.e., numerically) the Poisson
equation for the fluid pressure. Although analo-
gous to traditional LBM’s, this is interesting since
it is achieved without introducing compressibility
and/or thermal effects. In particular the present
theory does not rely on a state equation for the
fluid pressure.
2. is complete, namely all fluid fields are expressed as
momenta of the distribution function and all hy-
drodynamic equations are identified with suitable
moment equations of the LB inverse kinetic equa-
tion.
3. allows arbitrary initial and boundary conditions for
the fluid fields.
4. is self-consistent : the kinetic theory holds for ar-
bitrary, suitably smooth initial conditions for the
kinetic distribution function. In other words, the
initial kinetic distribution function must remain ar-
bitrary even if a suitable set of its momenta are
prescribed at the initial time.
5. the associated the kinetic and equilibrium distri-
bution functions can always be chosen to belong to
the class of non-Galilei-invariant distributions. In
particular the equilibrium kinetic distribution can
always be identified with a polynomial of second
degree in the velocity.
6. is non-asymptotic, i.e., unlike traditional LBM’s it
does not depend on any small parameter, in partic-
ular it holds for finite Mach numbers.
7. fulfills an entropic principle, based on a constant-H
theorem. This theorem assures, at the same time,
the strict positivity of the discrete kinetic distri-
bution function and the maximization of the as-
sociated Gibbs-Shannon entropy in a properly de-
fined functional class. Remarkably the constant H-
theorem is fulfilled for arbitrary (strictly positive)
kinetic equilibria. This includes also the case of
polynomial kinetic equilibria.
A further remarkable aspect of the theory concerns the
choice of the kinetic boundary conditions to be satisfied
by the distribution function (Axiom II) and obtained by
prescribing the form of the incoming-velocity distribution
[see Eq.(36)]. Thanks to Eqs.(34),(35), this requirement
[of the LB-IKT] the boundary conditions for the fluid
fields are satisfied exactly while the fluid equations are
by construction identically fulfilled also arbitrarily close
to the boundary. This result, in a proper sense, applies
only to Dirichlet boundary conditions for the fluid fields
[see Eqs.(8)]. Nevertheless the same approach can be
in principle extended to the case of mixed or Neumann
boundary conditions for the fluid fields.
Moreover, we have shown that a useful implication of
the theory is provided by the possibility of constructing
asymptotic approximations to the inverse kinetic equa-
tion. This permits to develop a new class of asymptotic
LBM’s which satisfy INSE with prescribed accuracy, to
obtain useful comparisons with previous CFD methods
(Chorin’s ACM) and to achieve accuracy estimates for
customary asymptotic LBM’s. The main results of the
paper are represented by THM’s 1-3, which refer respec-
tively to the construction of the integral LB-IKT, to the
entropic principle and to construction of the low effective-
Mach-numbers asymptotic approximations. For the sake
of reference, also another type of LB-IKT, which admits
as exact particular solution the polynomial kinetic equi-
librium, has been pointed out (THM.1bis).
The construction of a discrete inverse kinetic theory of
this type for the incompressible Navier-Stokes equations
represents an exciting development for the phase-space
description of fluid dynamics, providing a new starting
point for theoretical and numerical investigations based
on LB theory. In our view, the route to more accurate,
higher-order LBM’s, here pointed out, will be important
in order to achieve substantial improvements in the effi-
ciency of LBM’s in the near future.
APPENDIX A
The basic argument regarding the accuracy of the
boundary conditions adopted by customary asymptotic
LBM’s is provided by Ref.[46]. In fact. let us assume that
on the boundary δΩ the incoming distribution function
i (rw , t) is prescribed according to Eqs.(33),(37) and
(38), being f
oi (rw, t) prescribed suitably smooth func-
tions which are non vanishing only only for incoming dis-
crete velocities ai for which (ai −Vw) ·n(rw, t) ≤ 0. For
definiteness, let us assume that f
oi (rw , t) ≡ f
i (rw, t)
where f
i (rw, t) denotes a suitable equilibrium distribu-
tion. It follows that suitably close to the boundary the
kinetic distribution differs from the Chapman-Enskog so-
lution (25). The numerical error can be overcome only
discarding the first few spatial grid (close to the bound-
ary) in the numerical simulation [46].
APPENDIX B
Unlike standard kinetic theory, the distinctive feature
of LB-IKT’s is the possibility of adopting a non-Galilei
invariant kinetic distribution function (i.e., non-invariant
with respect to velocity translations). Here we re-
port another example of discrete inverse kinetic theory
of this type. Let us modify Axiom IV so that to permit
that a particular solution of LB-IKE [Eq.(30)] is pro-
vided by fi = f
i . Here we identify f
i with the (non-
Galilei invariant) polynomial kinetic distribution defined
by Eq.(27) but with the kinetic pressure p1 that replace
the fluid pressure p. In this case one can prove that the
source term Si reads
Si = S
i ≡ S̃i +∆Si, (85)
where
∆Si =
(ai −V) · ∇V−
ai∇·V
· ai+
V − ai
3ai ·V
+ (86)
wiρoai · ∇
ai ·V
Here N1 ≡ N−ρo
, where N is the Navier-Stokes opera-
tor (6), namely N1 is the nonlinear operator which acting
onV yields N1V = ρoV·∇V+∇ [p1 − Φ (r)]+f1−µ∇
Hence, invoking INSE, ∆Si can also be written in the
equivalent form
∆Si =
(ai −V) · ∇V−
ai∇·V
· ai+
V − ai
3ai ·V
+ (87)
wiρoai · ∇
ai ·V
The following result holds:
Theorem 1bis - Differential LB-IKT
In validity of axioms I-IV and the assumption that
fi = f
i is a particular solution of Eq.(30), the following
statements hold:
i is a particular solution of LB-IKE [Eq.(30)]
if and only if the extended fluid fields {V,p1} are strong
solutions of INSE of class (29), with initial and boundary
conditions (7)-(8), and arbitrary pseudo pressure po(t) of
class C(1)(I).
Moreover, for an arbitrary particular solution fi and
for arbitrary extended fluid fields :
For an arbitrary particular solution fi :
B) fi is a solution of LB-IKE [Eq.(30)] if and only
if the extended fluid fields {V,p1} are arbitrary strong
solutions of INSE of class (29), with initial and boundary
conditions (7)-(8), and arbitrary pseudo pressure po(t) of
class C(1)(I);
C) the moment equations of L-B IKE coincide identi-
cally with INSE in the set Ω× I;
D) the initial conditions and the (Dirichlet) boundary
conditions for the fluid fields are satisfied identically;
E) the source term Si is uniquely defined by
Eqs.(85),(86);
Proof:
The proof of propositions A,B, C and D is analogous
to that provided in THM.1. Assuming Si = S
i , the
proof of B follows from straightforward algebra. In fact,
letting fi(r, t) = f
i (r, t) for all (r, t) ∈ Ω × I in the
LB-IKE [Eq.(30)], one finds that Eq.(30) is fulfilled iff
the fluid fields satisfy the Navier-Stokes, isochoricity and
incompressibility equations (1),(2) and (3). The proof
of proposition E can be reached in a similar way. The
uniqueness of the source term Si is an immediate conse-
quence of the uniqueness of the solutions for INSE.
ACKNOWLEDGEMENTS Useful comments and
stimulating discussions with K.R. Sreenivasan, Direc-
tor, ICTP (International Center of Theoretical Physics,
Trieste, Italy) are warmly acknowledged. Research de-
veloped in the framework of PRIN Project Fundamen-
tals of kinetic theory and applications to fluid dynamics,
magnetofluid dynamics and quantum mechanics (MIUR,
Ministry for University and Research, Italy), with the
support of the Consortium for Magnetofluid Dynamics,
Trieste, Italy.
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|
0704.0340 | Phonon-mediated decay of an atom in a surface-induced potential | Phonon-mediated decay of an atom in a surface-induced potential
Fam Le Kien,1,∗ S. Dutta Gupta,1,2 and K. Hakuta1
Department of Applied Physics and Chemistry, University of Electro-Communications, Chofu, Tokyo 182-8585, Japan
School of Physics, University of Hyderabad, Hyderabad, India
(Dated: August 4, 2021)
We study phonon-mediated transitions between translational levels of an atom in a surface-induced
potential. We present a general master equation governing the dynamics of the translational states
of the atom. In the framework of the Debye model, we derive compact expressions for the rates
for both upward and downward transitions. Numerical calculations for the transition rates are
performed for a deep silica-induced potential allowing for a large number of bound levels as well
as free states of a cesium atom. The total absorption rate is shown to be determined mainly by
the bound-to-bound transitions for deep bound levels and by bound-to-free transitions for shallow
bound levels. Moreover, the phonon emission and absorption processes can be orders of magnitude
larger for deep bound levels as compared to the shallow bound ones. We also study various types of
transitions from free states. We show that, for thermal atomic cesium with temperature in the range
from 100 µK to 400 µK in the vicinity of a silica surface with temperature of 300 K, the adsorption
(free-to-bound decay) rate is about two times larger than the heating (free-to-free upward decay)
rate, while the cooling (free-to-free downward decay) rate is negligible.
PACS numbers: 34.50.Dy,33.70.Ca
I. INTRODUCTION
Over the past few years, tight confinement of cold
atoms has drawn considerable attention. The interest in
this area is motivated not only by the fundamental na-
ture of the problem, but also by its potential applications
in atom optics and quantum information. A method
for microscopic trapping and guiding of individual atoms
along a nanofiber has been proposed [1]. Surface–atom
quantum electrodynamic effects have constituted another
interesting area, where a great deal of work has been
carried out. Modification of spontaneous emission of
an atom [2] and radiative exchange between two distant
atoms [3] mediated by a nanofiber have been investigated.
Surface-induced deep potentials have played a major role
and have received due attention in recent years. Oria
et al. have studied various theoretical schemes to load
atoms into such potentials [4, 5]. A rigorous theory of
spontaneous decay of an atom in a surface-induced po-
tential invoking the density-matrix formalism has been
developed [6]. The role of interference between the emit-
ted and reflected fields and also the role of transmission
into the evanescent modes were identified. Further cal-
culations on the excitation spectrum have been carried
out [7]. Bound-to-bound transitions were shown to lead
to significant effects like a large red tail of the excita-
tion spectrum as compared to the weak consequences of
free-to-bound transitions. A crucial step in this direction
was the experimental observation of the excitation spec-
trum and the channeling of the fluorescent photons along
the nanofiber [8], opening up avenues for novel quantum
information devices.
In most of the problems involving surface–atom inter-
action, the macroscopic surface is usually kept at room
temperature. Thus the pertinent question that can be
asked is what would be the effect of heating on the cold
atoms. It is understood that transfer of heat to the
trapped atoms will lead to a change in the occupation
probability of the vibrational levels as well as their co-
herence. Phonon-induced changes in the populations of
the vibrational levels have been studied by several groups
[5, 9, 10]. In a nice and compact treatment based on the
dyadic Green function and the Fermi golden rule, Henkel
et al. showed that the effects can be very different de-
pending on the nature of the atomic/molecular species
[9]. The time scales for various species were estimated.
It should be stressed that the trap considered by Henkel
et al. was not necessarily a surface trap and misses out
on many of the aspects of the surface–atom interaction
[9]. Based on the assumption that the surface–atom in-
teraction can be represented by a Morse potential, the
phonon-mediated decay was estimated by Oria et al. [5].
Their estimate was based on the formalism developed by
Gortel et al. [10]. However, all the previous theories
focus on only the transition rates and thus are not gen-
eral enough. In this paper, we present a general density-
matrix formalism to calculate the phonon-mediated de-
cay of populations as well as the changes in coherence.
We derive the relevant master equation for the density
matrix of the atom. We emphasize that our density-
matrix equation describes the full dynamics of the cou-
pling between trapped atoms and phonons and does not
assume any particular form of the trapping potential.
Under the Debye approximation, we derive compact ex-
pressions for the phonon-mediated decay rates. Numer-
ical calculations are carried out assuming the potential
model considered in [4]. In contrast to the previous work,
we include a large number of vibrational levels due to the
deep surface–atom potential. We show that there can be
significant differences in the decay rates when the initial
level is chosen as one of the shallow or deep bound levels.
We also calculate and analyze the decay rates for various
http://arxiv.org/abs/0704.0340v1
types of transitions from free states.
The paper is organized as follows. In Sec. II we de-
scribe the model. In Sec. III we derive the basic dynam-
ical equations for the phonon-mediated decay processes.
In Sec. IV we present the results of numerical calcula-
tions. Our conclusions are given in Sec. V.
II. DESCRIPTION OF THE MODEL SYSTEM
We assume the whole space to be divided into two re-
gions, namely, the half-space x < 0, occupied by a nondis-
persive nonabsorbing dielectric medium (medium 1), and
the half-space x > 0, occupied by vacuum (medium 2).
We examine a single atom moving in the empty half-
space x > 0. We assume that the atom is in a fixed
internal state |i〉 with energy h̄ωi. Without loss of gen-
erality, we assume that the energy of the internal state
|i〉 is zero, i.e. ωi = 0. We describe the interaction be-
tween the atom and the surface. We first consider the
surface-induced interaction potential and then add the
atom-phonon interaction.
A. Surface-induced interaction potential
In this subsection, we describe the interaction between
the atom and the surface in the case where thermal vi-
brations of the surface are absent. The potential en-
ergy of the surface–atom interaction is a combination of
a long-range van der Waals attraction and a short-range
repulsion [11]. Despite a large volume of research on the
surface–atom interaction, due to the complexity of sur-
face physics and the lack of data, the actual form of the
potential is yet to be ascertained [11]. For the purpose
of numerical demonstration of our formalism, we choose
the following model for the potential [4, 11]:
U(x) = Ae−αx − C3
. (1)
Here, C3 is the van der Waals coefficient, while A and
α determine the height and range, respectively, of the
surface repulsion. The potential parameters C3, A, and
α depend on the nature of the dielectric and the atom.
In numerical calculations, we use the parameters of fused
silica, for the dielectric, and the parameters of ground-
state atomic cesium, for the atom. The parameters for
the interaction between silica and ground-state atomic
cesium are theoretically estimated to be C3 = 1.56 kHz
µm3, A = 1.6× 1018 Hz, and α = 53 nm−1 [6].
We introduce the notation ϕν(x) for the eigenfunc-
tions of the center-of-mass motion of the atom in the
potential U(x). They are determined by the stationary
Schrödinger equation
+ U(x)
ϕν(x) = Eνϕν(x). (2)
Here m is the mass of the atom. In the numerical ex-
ample with atomic cesium, we have m = 132.9 a.u.
= 2.21 × 10−25 kg. The eigenvalues Eν are the center-
of-mass energies of the translational levels of the atom.
These eigenvalues are the shifts of the energies of the
translational levels from the energy of the internal state
|i〉. Without loss of generality, we assume that the
center-of-mass eigenfunctions ϕν(x) are real functions,
i.e. ϕ∗ν(x) = ϕν(x).
In Fig. 1, we show the potential U(x) and the wave
functions ϕν(x) of a number of bound levels with en-
ergies in the range from −1 GHz to −5 MHz. We also
plot the wave function of a free state with energy of about
4.25 MHz. In order to have some estimate about the spa-
tial extent of a wave function ϕν(x), we define a crossing
point xcross, which corresponds to the rightmost solution
of the equation U(x) = Eν . Note that, for shallow lev-
els, the wave function generally peaks close to the point
xcross. We plot the eigenvalue modulus |Eν | and the cross-
ing point xcross in Figs. 2(a) and 2(b), respectively. It is
clear from the figure that, for ν in the range from 0 to
300, the eigenvalue varies dramatically from about 158
THz to about 322 kHz, while the wave function extends
only up to 170 nm.
FIG. 1: Energies and wave functions of the center-of-mass
motion of an atom in a surface-induced potential. The pa-
rameters of the potential are C3 = 1.56 kHz µm
3, A =
1.6 × 1018 Hz, and α = 53 nm−1. The mass of the atom
is m = 2.21 × 10−25 kg. We plot bound levels with energies
in the range from −1 GHz to −5 MHz and also a free state
with energy of about 4.25 MHz.
FIG. 2: Eigenvalue modulus |Eν | (a) and crossing point xcross
(b) as functions of the vibrational quantum number ν. The
parameters used are as in Fig. 1.
We introduce the notation |ν〉 = |ϕν〉 and ων = Eν/h̄
for the state vectors and frequencies of translational lev-
els. Then, the Hamiltonian of the atom in the surface-
induced potential can be represented in the diagonal form
h̄ωνσνν . (3)
Here, σνν = |ν〉〈ν| is the population operator for the
translational level ν. We emphasize that the summation
over ν includes both the discrete (Eν < 0) and continuous
(Eν > 0) spectra. The levels ν with Eν < 0 are called
the bound (or vibrational) levels. In such a state, the
atom is bound to the surface. It is vibrating, or more
exactly, moving back and forth between the walls formed
by the van der Waals part and the repulsive part of the
potential. The levels ν with Eν > 0 are called the free (or
continuum) levels. The center-of-mass wave functions of
the bound states are normalized to unity. The center-of-
mass wave functions of the free states are normalized to
the delta function of energy.
B. Atom–phonon interaction
In this subsection, we incorporate the thermal vibra-
tions of the solid into the model. Due to the thermal
effects, the surface of the dielectric vibrates. The surface-
induced potential for the atom is then U(x− xs), where
xs is the displacement of the surface from the mean po-
sition 〈xs〉 = 0. We approximate the vibrating potential
U(x− xs) by expanding it to the first order in xs,
U(x− xs) = U(x) − U ′(x)xs. (4)
The first term, U(x), when combined with the kinetic
energy p2/2m, yields the Hamiltonian HA [see Eq. (3)],
which leads to the formation of translational levels of
the atom. The second term, −U ′(x)xs, accounts for the
thermal effects in the interaction of the atom with the
solid. Note that the quantity F = −U ′(x) is the force
of the surface upon the atom. Hence, the force of the
atom upon the surface is −F = U ′(x) and, consequently,
U ′(x)xs is the work required to displace the surface for
a small distance xs.
It is well known that, for a smooth surface, the gas
atom interacts only with the phonons polarized along
the x direction [10]. In the harmonic approximation, we
2MNωq
iqR + b†qe
−iqR). (5)
Here, M is the mass of a particle of the solid, N is
the particle number density, ωq and q are the frequency
and wave vector of the x-polarized acoustic phonons, re-
spectively, R = (0, y, z) is the lateral component of the
position vector (x, y, z) of the atom, and bq and b
q are
the annihilation and creation phonon operators, respec-
tively. Without loss of generality, we choose R = 0.
Meanwhile, the operator U ′ can be decomposed as U ′ =
νν′ σνν′ 〈ν|U ′|ν′〉, where σνν′ = |ν〉〈ν′| is the operator
for the translational transition ν ↔ ν′. Hence, the en-
ergy term −U ′(x)xs leads to the atom–phonon interac-
tion Hamiltonian [10]
HI = h̄
S(bq + b
q), (6)
gνν′σνν′ . (7)
Here we have introduced the atom–phonon coupling co-
efficients
gνν′ =
Fνν′√
2MNh̄
, (8)
Fνν′ = −
ϕν(x)U
′(x)ϕν′ (x)dx (9)
being the matrix elements for the force of the surface
upon the atom. We note that Fνν′ = −mω2νν′xνν′ , where
xνν′ = 〈ν|x|ν′〉 and ωνν′ = ων −ων′ are the surface–atom
dipole matrix element and the translational transition
frequency, respectively. Hence, the coupling coefficient
gνν′ depends on the dipole matrix element xνν′ and the
transition frequency ωνν′ . Since ωνν = 0, we have gνν =
We note that the Hamiltonian of the x-polarized acous-
tic phonons is given by
h̄ωqb
qbq. (10)
The total Hamiltonian of the atom–phonon system is
H = HA +HI +HB. (11)
We use the above Hamiltonian to study the phonon-
mediated decay of the atom.
III. DYNAMICS OF THE ATOM
In this section, we present the basic equations for the
phonon-mediated decay processes. We derive a general
master equation for the reduced density operator of the
atom in subsection IIIA, obtain analytical expressions
for the relaxation rates and frequency shifts in subsec-
tion III B, and calculate the rates and the shifts in the
framework of the Debye model in subsection III C.
A. Master equation
In the Heisenberg picture, the equation for the phonon
operator bq(t) is
ḃq(t) = −iωqbq(t)−
S(t), (12)
which has a solution of the form
bq(t) = bq(t0)e
−iωq(t−t0) − iWq(t). (13)
Here, t0 is the initial time and Wq is given by
Wq(t) =
e−iωq(t−τ)S(τ) dτ. (14)
Consider an arbitrary atomic operator O which acts only
on the atomic states but not on the phonon states. The
time evolution of this operator is governed by the Heisen-
berg equation
∂O(t)
[HA(t) +HI(t),O(t)], (15)
which, with account of Eqs. (6) and (13), yields
∂O(t)
[HA(t),O(t)]
[S(t),O(t)][bq(t0)e−iωq(t−t0) − iWq(t)]
[b†q(t0)e
iωq(t−t0) + iW †q(t)][O(t), S(t)].
We assume the initial density of the atom–phonon sys-
tem to be the direct product state
ρΣ(t0) = ρ(t0)ρB(t0), (17)
with the atom in an arbitrary state ρ(t0) and the phonons
in a thermal state
ρB(t0) = Z
−1 exp[−HB(t0)/kBT ]. (18)
Here, Z is the normalization constant and T is the tem-
perature of the phonon bath. For the initial condition
(17), the Bogolubov’s lemma [12], applied to an arbitrary
operator Θ(t), asserts the following:
〈Θ(t)bq(t0)〉 = n̄q〈[bq(t0),Θ(t)]〉, (19)
where the mean number of phonons in the mode q is
given by
n̄q =
exp(h̄ωq/kBT )− 1
. (20)
Let Θ be an atomic operator. We then have the commu-
tation relation [bq(t),Θ(t)] = 0, which yields
[bq(t0),Θ(t)] = ie
iωq(t−t0)[Wq(t),Θ(t)]. (21)
Combining Eq. (19) with Eq. (21) leads to
〈Θ(t)bq(t0)〉 = ieiωq(t−t0)n̄q〈[Wq(t),Θ(t)]〉. (22)
We perform the quantum mechanical averaging for ex-
pression (16) and use Eq. (22) to eliminate the phonon
operators bq(t0) and b
q(t0). The resulting equation can
be written as
∂〈O(t)〉
〈[HA(t),O(t)]〉
n̄q + 1√
〈[S(t),O(t)]Wq(t) +W †q(t)[O(t), S(t)]〉
〈Wq(t)[O(t), S(t)] + [S(t),O(t)]W †q(t)〉.
We note that Eq. (23) is exact. It does not contain
phonon operators explicitly. The dependence on the
phonon operators is hidden in the time shift of the oper-
ator S(τ) in expression (14) for the operator Wq(t).
We now show how the dependence of the operator
Wq(t) on the phonon operators can be approximately
eliminated. We assume that the atom–phonon coupling
coefficients gνν′ are small. The use of the zeroth-order
approximation σνν′(τ) = σνν′ (t)e
iωνν′(τ−t) in the expres-
sion for S(τ) [see Eq. (7)] yields
S(τ) =
gνν′σνν′ (t)e
iωνν′(τ−t), (24)
which is accurate to first order in the coupling coeffi-
cients. Inserting Eq. (24) into Eq. (14) gives
Wq(t) =
gνν′σνν′(t)δ−(ων′ν − ωq), (25)
where
δ−(ω) = lim
e−i(ω+iǫ)τ dτ
δ(ω). (26)
Here, in order to take into account the effect of adiabatic
turn-on of interaction, we have added a small positive
parameter ǫ to the integral and have used the limit t0 →
−∞. Introducing the notation
gνν′σνν′δ−(ων′ν − ωq), (27)
we can rewrite Eq. (23) in the form
∂〈O(t)〉
〈[HA(t),O(t)]〉
(n̄q + 1)〈[S(t),O(t)]Kq(t) +K†q(t)[O(t), S(t)]〉
n̄q〈Kq(t)[O(t), S(t)] + [S(t),O(t)]K†q(t)〉. (28)
In order to examine the time evolution of the reduced
density operator ρ(t) of the atom in the Schrödinger
picture, we use the relation 〈O(t)〉 = Tr[O(t)ρ(0)] =
Tr[O(0)ρ(t)], transform to arrange the operator O(0) at
the first position in each operator product, and eliminate
O(0). Then, we obtain the Liouville master equation
∂ρ(t)
= − i
[HA, ρ(t)]
(n̄q + 1){[Kqρ(t), S] + [S, ρ(t)K†q]}
n̄q{[S, ρ(t)Kq] + [K†qρ(t), S]}. (29)
Equations (28) and (29) are valid to second order in
the coupling coefficients. These equations allow us to
study the time evolution and dynamical characteristics of
the atom interacting with the thermal phonon bath. We
note that Eq. (29) is a particular form of the Zwanzig’s
generalized master equation, which can be obtained by
the projection operator method [13].
B. Relaxation rates and frequency shifts
We use Eq. (29) to derive an equation for the matrix
elements ρjj′ ≡ 〈j|ρ|j′〉 of the reduced density operator
of the atom. The result is
∂ρjj′
= −iωjj′ρjj′ +
(γejj′νν′ + γ
jj′νν′)ρνν′
[(γejν + γ
jν)ρνj′ + (γ
j′ν + γ
j′ν)ρjν ], (30)
where the coefficients
γejj′νν′ = 2π
n̄q + 1
gjνgj′ν′ [δ−(ωνj − ωq)
+ δ+(ων′j′ − ωq)],
γejν = 2π
n̄q + 1
gjµgνµδ−(ωνµ − ωq) (31)
γajj′νν′ = 2π
gjνgj′ν′ [δ−(ωj′ν′ − ωq)
+ δ+(ωjν − ωq)],
γajν = 2π
gjµgνµδ+(ωµν − ωq) (32)
are the decay parameters associated with the phonon
emission and absorption, respectively. Here, the nota-
tion δ+(ω) = δ
−(ω) has been used.
Equation (30) describes phonon-induced variations in
the populations and coherences of the translational levels
of the atom. We analyze the characteristics of the relax-
ation processes. For simplicity of mathematical treat-
ment, we first consider only transitions from discrete lev-
els. The equation for the diagonal matrix element ρjj for
a discrete level j can be written in the form
(γejjνν + γ
jjνν )ρνν
− (γejj + γajj + c.c.)ρjj
+ off-diagonal terms. (33)
When the off-diagonal terms are neglected, Eq. (33) re-
duces to a simple rate equation. It is clear from Eq. (33)
that the rate for the downward transition from an upper
level l to a lower level k (k < l) is
Rekl = γ
kkll = 2π
n̄q + 1
g2lkδ(ωlk − ωq), (34)
while the rate for the upward transition from a lower level
k to an upper level l (l > k) is
Ralk = γ
llkk = 2π
g2lkδ(ωlk − ωq). (35)
Equations (34) and (35) are in agreement with the re-
sults of Gortel et al. [10], obtained by using the Fermi
golden rule. We note that Rekl and R
lk with l ≤ k are
mathematically equal to zero because they have no phys-
ical meaning. For convenience, we introduce the notation
Rlk = R
lk, R
lk, or 0 for l < k, l > k, or l = k, respec-
tively. It is clear that the off-diagonal coefficients Rlk
with l 6= k are the rates of transitions. However, the di-
agonal coefficients Rkk have no physical meaning and are
mathematically equal to zero.
As seen from Eq. (33), the phonon-mediated depletion
rate of a level k is Γkk = 2Re(γ
kk + γ
kk). The explicit
expression for this rate is
Γkk = 2π
n̄q + 1
g2kµδ(ωkµ − ωq)
g2µkδ(ωµk − ωq). (36)
We note that Γkk =
µk +R
µk) =
µ Rµk. We can
write Γkk = Γ
kk + Γ
kk, where
Γekk =
Reµk (37)
Γakk =
Raµk (38)
are the contributions due to downward transitions
(phonon emission) and upward transitions (phonon ab-
sorption), respectively. In the above equations, the sum-
mation over µ can be extended to cover not only the
discrete levels but also the continuum levels.
Meanwhile, the equation for the off-diagonal matrix
element ρlk for a pair of discrete levels l and k can be
written in the form ∂ρlk/∂t = −(iωlk + γell + γall + γe∗kk +
γa∗kk)ρlk + . . . , or, equivalently,
= −i(ωlk +∆lk − iΓlk)ρlk + . . . . (39)
Here the frequency shift ∆lk is given by
∆lk =
n̄q + 1
ωlµ − ωq
ωµk + ωq
ωlµ + ωq
ωµk − ωq
, (40)
while the coherence decay rate Γlk is expressed as
Γlk = π
n̄q + 1
g2lµδ(ωlµ − ωq) + g2kµδ(ωkµ − ωq)
g2µlδ(ωµl − ωq) + g2µkδ(ωµk − ωq)
When we set l = k in Eq. (40), we find ∆kk = 0.
When we set l = k in Eq. (41), we recover Eq. (36).
We note that Γlk =
µl + R
µk + R
µl + R
µk)/2 =
µ(Rµl + Rµk)/2. Comparison between Eqs. (41) and
(36) yields the relation Γlk = (Γll +Γkk)/2. We can also
write Γlk = Γ
lk + Γ
lk, where Γ
µl + R
µk)/2
and Γalk =
µl + R
µk)/2 are the contributions due
to downward transitions (phonon emission) and upward
transitions (phonon absorption), respectively. In the
above equations, the summation over µ can be extended
to cover not only the discrete levels but also the contin-
uum levels.
We now discuss phonon-mediated transitions from con-
tinuum (free) levels. We start by considering free-to-
bound transitions. For a continuum level f with energy
Ef > 0, the center-of-mass wave function ϕf (x) is nor-
malized per unit energy. In this case, the quantity Rνf
becomes the density of the transition rate. A free level f
can be approximated by a level of a quasicontinuum [14].
A discretization of the continuum can be realized by using
a large box of length L with reflecting boundary condi-
tions [15]. We label En the energies of the eigenstates
in the box and φn(x) the corresponding wave functions.
Note that such states are standing-wave states [14, 15].
The relation between a quasicontinuum-state wave func-
tion φnf (x), normalized to unity in the box, and the cor-
responding continuum-state wave function ϕf (x), nor-
malized per unit energy, with equal energies Enf = Ef ,
is [15]
ϕf (x) ∼=
]−1/2
φnf (x)
)1/2 (
φnf (x). (42)
Consequently, for a single atom initially prepared in the
quasicontinuum standing-wave state |nf 〉 = |φnf 〉, the
rate for the transition to an arbitrary bound state |ν〉 is
approximately given by
Gνf =
vfRνf , (43)
where vf = (2Ef/m)1/2 is the velocity of the atom in the
initial standing-wave state |f〉. The phonon-mediated
free-to-bound decay rate (adsorption rate) is then given
Gνf , (44)
where the summation includes only bound levels. It is
clear from Eq. (43) that, in the continuum limit L → ∞,
the rate Gνf tends to zero. This is because a free atom
can be anywhere in free space and therefore the effect of
phonons on a single free atom is negligible.
In order to get deeper insight into the free-to-bound
transition rate density Rνf , we consider a macroscopic
atomic ensemble in the thermodynamic limit [14]. Sup-
pose that there are N0 atoms in a volume with a large
length L and a transverse cross section area S0. Assume
that all the atoms are in the same quasicontinuum state
|nf 〉 and interact with the dielectric independently. The
rate for the transitions of the atoms from the quasicon-
tinuum state |nf〉 to an arbitrary bound state |ν〉, defined
as the time derivative of the number of atoms in the state
|ν〉, is Dνf = N0Gνf . In order to get the rate for the con-
tinuum state |f〉, we need to take the thermodynamical
limit, where L → ∞ and N0 → ∞ but N0/L remains
constant. Then, the rate for the transitions of the atoms
from the continuum state |f〉 to an arbitrary bound
state |ν〉 is given by Dνf = πh̄ρ0S0vfRνf = 2πh̄NfRνf .
Here, ρ0 = N0/LS0 is the atomic number density and
Nf = ρ0S0vf/2 is the number of atoms incident into the
dielectric surface per unit time. It is clear that the tran-
sition rate Dνf is proportional to the incidence rate Nf
as well as the transition rate density Rνf . We emphasize
that Dνf is a characteristics for a macroscopic atomic en-
semble in the thermodynamic limit while Gνf is a mea-
sure for a single atom. When the length of the box, L,
and the number of atoms, N0, are finite, the dynamics of
the atoms cannot be described by the free-to-bound rate
Dνf directly. Instead, we must use the transition rate
per atom Gνf = Dνf/N0, which depends on the length
L of the box that contains the free atoms [see Eq. (43)].
In a thermal gas, the atoms have different velocities
and, therefore, different energies. For a thermal Maxwell-
Boltzmann gas with temperature T0, the distribution of
the kinetic energy Ef of the atomic center-of-mass motion
along the x direction is
P (Ef ) =
πkBT0
e−Ef/kBT0
. (45)
The transition rate to an arbitrary bound state |ν〉 is
then given by GνT0 =
GνfP (Ef ) dEf , i.e.
GνT0 =
e−Ef/kBT0RνfdEf , (46)
where λD = (2πh̄
2/mkBT0)
1/2 is the thermal de Broglie
wavelength. The phonon-mediated free-to-bound decay
rate (adsorption rate) is given by
GT0 =
GνT0 =
GfP (Ef ) dEf . (47)
In the above equation, the summation over ν includes
only bound levels. Note that Eq. (46) is in qualitative
agreement with the results of Refs. [5, 14].
It is easy to extend the above results to the case of
free-to-free transitions. Indeed, it can be shown that the
density of the rate for the transition from a quasicontin-
uum state |nf 〉, which corresponds to a free state |f〉, to
a different free state |f ′〉 is given by
Qf ′f =
vfRf ′f . (48)
For convenience, we introduce the notation Qef ′f = Qf ′f
or 0 for Ef ′ < Ef or Ef ′ ≥ Ef , respectively, and Qaf ′f =
Qf ′f or 0 for Ef ′ > Ef or Ef ′ ≤ Ef , respectively. Then, we
have Qf ′f = Q
f ′f , 0, or Q
f ′f for Ef ′ < Ef , Ef ′ = Ef , or
Ef ′ > Ef , respectively. The downward (phonon-emission)
and upward (phonon-absorption) free-to-free decay rates
for the free state |f〉 are given by
Qef =
Qef ′fdEf ′ (49)
Qaf =
Qaf ′fdEf ′ , (50)
respectively. The total free-to-free decay rate for the free
state |f〉 is Qf = Qef +Qaf =
Qf ′fdEf ′ .
For a thermal gas, we need to replace the transition
rate density Qf ′f and the decay rate Qf by Qf ′T0 =
Qf ′fP (Ef ) dEf and QT0 =
QfP (Ef ) dEf , respec-
tively, which are the averages of Qf ′f and Qf , respec-
tively, with respect to the energy distribution P (Ef )
of the initial state. Like in the other cases, we have
Qf ′T0 = Q
f ′T0
+Qaf ′T0 and QT0 = Q
+QaT0 , where
Qef ′T0 =
Qef ′fP (Ef ) dEf ,
Qaf ′T0 =
∫ Ef′
Qaf ′fP (Ef ) dEf (51)
are the downward and upward transition rate densities
QeT0 =
QefP (Ef ) dEf ,
QaT0 =
QafP (Ef ) dEf (52)
are the downward and upward decay rates. The thermal
decay ratesQeT0 andQ
describe the cooling and heating
processes, respectively. It can be easily shown thatQeT0 <
QaT0 , Q
> QaT0 , and Q
= QaT0 when T0 < T , T0 >
T , and T0 = T , respectively. The relation Q
< QaT0
(QeT0 > Q
), obtained for T0 < T (T0 > T ), indicates
the dominance of heating (cooling) of free atoms by the
surface.
C. Relaxation rates and frequency shifts in the
framework of the Debye model
In order to get insight into the relaxation rates and
frequency shifts, we approximate them using the Debye
model for phonons. In this model, the phonon frequency
ωq is related to the phonon wave number q as ωq = vq,
where v is the sound velocity. Furthermore, the summa-
tion over the first Brillouin zone is replaced by an integral
over a sphere of radius qD = (6π
2N/V )1/3, where V is the
volume of the solid. The Debye frequency and the Debye
temperature are given by ωD = vqD and TD = h̄ωD/kB,
respectively. For fused silica, we have v = 5.96 km/s,
NM/V = 2.2 g/cm3, and M = 9.98× 10−26 kg [16]. Us-
ing these parameters, we find qD = 109.29 × 106 cm−1,
ωD = 10.4 THz, and TD = 498 K. In order to perform
the summation over phonon states in the framework of
the Debye model, we invoke the thermodynamic limit,
i.e., replace
· · · = V
|q|≤qD
. . . dq =
. . . ω2qdωq. (53)
Then, for transitions between an upper level l and a lower
level k, where 0 < ωlk < ωD, Eqs. (34) and (35) yield
Rekl =
Mh̄ω3D
(n̄lk + 1)ωlkF
lk (54)
Ralk =
Mh̄ω3D
n̄lkωlkF
lk. (55)
Here, n̄lk is given by Eq. (20) with ωq replaced by ωlk.
We emphasize that, according to Eqs. (54) and (55),
the phonon-emission rate Rekl and the phonon-absorption
rate Ralk depend not only on the matrix element Flk of
the force but also on the translational transition fre-
quency ωlk. The frequency dependences of the transi-
tion rates are comprised of the frequency dependences
of the mean phonon number n̄lk, the phonon mode den-
sity 3Nω2lk/ω
D, and the matrix element Flk = −U ′lk =
−mω2lkxlk of the force. An additional factor comes from
the presence of the phonon frequency in Eq. (5) for
the surface displacement and, consequently, in the atom–
phonon interaction Hamiltonian (6). It is clear that an
increase in the phonon frequency leads to a decrease in
the mean phonon number and an increase in the phonon
mode density. The matrix element of the force usu-
ally first increases and then decreases with increasing
phonon frequency. Due to the existence of several com-
peting factors, the frequency dependences of the tran-
sition rates are rather complicated. They usually first
increase and then decrease with increasing phonon fre-
quency. We note that, for transitions with ωlk > ωD, we
have Rekl = R
lk = 0.
We conclude this section by noting that the use of Eq.
(53) in Eq. (40) yields the frequency shift
∆lk = ∆
lk +∆
lk , (56)
where
2Mh̄ω3D
F 2lµ
ωlµ − ω
F 2µk
ωµk + ω
ωdω (57)
Mh̄ω3D
ω2lµ − ω2
ω2µk − ω2
n̄ωωdω
are the zero- and finite-temperature contributions, re-
spectively. In Eq. (58), n̄ω is given by Eq. (20) with ωq
replaced by ω.
IV. NUMERICAL RESULTS AND
DISCUSSIONS
In this section, we present the numerical results based
on the analytical expressions derived in the previous
section for the phonon-mediated relaxation rates of the
translational levels of the atom. In particular, we use
Eqs. (54) and (55), obtained in the framework of the De-
bye model, for our numerical calculations. We consider
transitions from bound states as well as free states. The
transitions from bound states to other translational lev-
els occur in the case where the atom is initially already
adsorbed or trapped near the surface. The transitions
from free states to other translational levels occur in the
processes of adsorbing, heating, and cooling of free atoms
by the surface. Due to the difference in physics of the ini-
tial situations, we study the transitions from bound and
free states separately.
A. Transitions from bound states
FIG. 3: Phonon-emission rates Reν′ν from the vibrational lev-
els (a) ν = 280 and (b) ν = 120 to other levels ν′ as functions
of the lower-level energy Eν′ . The arrows mark the initial
states. The parameters of the solid are M = 9.98 × 10−26 kg
and ωD = 10.4 THz. The temperature of the phonon bath is
T = 300 K. Other parameters are as in Fig. 1.
FIG. 4: Phonon-absorption rates Raν′ν from the vibrational
levels (a) ν = 280 and (b) ν = 120 to other levels ν′ as func-
tions of the upper-level energy Eν′ . The left (right) panel
in each row corresponds to bound-to-bound (bound-to-free)
transitions. The arrows mark the initial states. The param-
eters used are as in Fig. 3. The temperature of the phonon
bath is T = 300 K.
We start from a given bound level and calculate the
rates of phonon-mediated atomic transitions, both down-
ward and upward. The profiles of the phonon-emission
(downward-transition) rate Reν′ν [see Eq. (54)] and the
phonon-absorption (upward-transition) rateRaν′ν [see Eq.
(55)] are shown in Figs. 3 and 4, respectively. The upper
(lower) part of each of these figures corresponds to the
case of the initial level ν = 280 (ν = 120), with energy
Eν = −156 MHz (Eν = −8.4 THz). The left (right) panel
of Fig. 4 corresponds to bound-to-bound (bound-to-free)
upward transitions. The temperature of the surface is
assumed to be T = 300 K. As seen from Figs. 3 and 4,
the transition rates have pronounced localized profiles.
Due to the competing effects of the mean phonon num-
ber, the phonon mode density, and the matrix element
of the force, the transition rates usually first increase
and then decrease with increasing phonon frequency. It
is clear from a comparison of Figs. 3(a) and 3(b) and
also a comparison of Figs. 4(a) and 4(b) that transitions
from shallow levels have probabilities orders of magni-
tude lower than those from deeper levels. The main rea-
son is that the wave functions of the shallow states are
spread further away from the surface than those for the
deep states. Due to this difference, the effects of the sur-
face vibrations are weaker for the shallow levels than for
the deep levels. Another pertinent feature that should
be noted from the figure is the following: Since transi-
tion frequencies involved are large, they may overshoot
the Debye frequency ωD = 10.4 THz, leading to a cutoff
on the lower (higher) side of the frequency axis for the
emission (absorption) curve.
In order to see the overall effect of the individual tran-
sition rates shown above, we add them up. First we ex-
amine the phonon-absorption rates of bound levels. The
total phonon-absorption rate Γaνν of a bound level ν is the
sum of the individual absorption rates Raµν over all the
upper levels µ, both bound and free [see Eq. (38)]. We
plot in Fig. 5 the contributions to Γaνν from two types
of transitions, bound-to-bound and bound-to-free (des-
orption) transitions. The solid curve of the figure shows
that the bound-to-bound phonon-absorption rate is large
(above 1010 s−1) for deep and intermediate levels. How-
ever, it reduces dramatically with increasing ν in the
region of large ν and becomes very small (below 10−5
s−1) for shallow levels. Meanwhile, the dashed curve of
Fig. 5 shows that the bound-to-free phonon-absorption
rate (i.e., the desorption rate) is zero for deep levels, since
the energy required for the transition is greater than the
Debye energy [5]. However, the desorption rate is sub-
stantial (above 105 s−1) for intermediate and shallow lev-
els. Thus, the total phonon-absorption rate Γaνν is mainly
determined by the bound-to-bound transitions in the case
of deep levels and by the bound-to-free transitions in the
case of shallow levels. One of the reasons for the dramatic
reduction of the bound-to-bound phonon-absorption rate
in the region of shallow levels is that the number of up-
per bound levels µ becomes small. The second reason is
that the frequency of each individual transition becomes
small, leading to a decrease of the phonon mode den-
sity. The third reason is that the center-of-mass wave
functions of shallow levels are spread far away from the
surface, leading to a reduction of the effect of phonons
on the atom.
Unlike the bound-to-bound phonon-absorption rate,
the bound-to-free phonon-absorption rate is substantial
in the region of shallow levels. This is because the free-
state spectrum is continuous and the range of the bound-
to-free transition frequency can be large (up to the De-
bye frequency ωD = 10.4 THz). The gradual reduction of
the bound-to-free phonon-absorption rate in the region of
shallow levels is mainly due to the reduction of the time
that the atom spends in the proximity of the surface.
FIG. 5: Contributions of bound-to-bound (solid curve) and
bound-to-free (dashed curve) transitions to the total phonon-
absorption rate Γaνν versus the vibrational quantum number
ν of the initial level. The parameters used are as in Fig. 3.
The temperature of the phonon bath is T = 300 K.
The total phonon-emission rate Γeνν [see Eq. (37)] and
the total phonon-absorption rate Γaνν [see Eq. (38)] are
shown in Fig. 6 by the solid and dashed curves, respec-
tively. It is clear from the figure that emission is com-
parable to but slightly stronger than absorption. Such a
dominance is due to the fact that phonon emission moves
the atom to a center-of-mass state closer to the surface
while phonon absorption changes the atomic state in the
opposite direction (see Figs. 1 and 2). Our results for the
rates are in good qualitative agreement with the results
of Oria et al., albeit with the Morse potential [5]. We
stress that we include a large number of vibrational lev-
els as a consequence of the deep silica–cesium potential.
Note that the earlier work on this theme involved much
fewer levels [5].
FIG. 6: Phonon-emission decay rate Γeνν (solid lines) and
phonon-absorption decay rate Γaνν (dashed lines) of a bound
level as functions of the vibrational quantum number ν. The
inset shows the rates in the linear scale to highlight the dif-
ferences in the dissociation limit. The parameters used are as
in Fig. 3. The temperature of the phonon bath is T = 300 K.
FIG. 7: Same as in Fig. 6 except that T = 30 K.
We next study the effect of temperature on the decay
rates. The results for the phonon-mediated decay rates
for T = 30 K are shown in Fig. 7. In contrast to Fig. 6,
the absorption rate is now much smaller than the corre-
sponding emission rate for both shallow and deep levels.
Thus, while it is difficult to distinguish the two log-scale
curves for deep and shallow levels at room temperature
(see Fig. 6), they are well resolved at low temperature.
B. Transitions from free states
We now calculate the rates for transitions from free
states to other levels. We first examine free-to-bound
transitions, which correspond to the adsorption process.
According to Eq. (43), the free-to-bound (more exactly,
quasicontinuum-to-bound) transition rate Gνf depends
not only on the continuum-to-bound transition rate den-
sity Rνf but also on the length L of the free-atom quan-
tization box. To be specific, we use in our numerical
calculations the value L = 1 mm, which is a typical size
of atomic clouds in magneto-optical traps [17].
FIG. 8: Free-to-bound transition rates Gνf for transitions
from the free plane-wave states with energies (a) Ef = 2 MHz
and (b) Ef = 3.1 THz to bound levels ν as functions of the
bound-level energy Eν . The arrows mark the energies of the
initial free states. The insets show Gνf on the log scale versus
Eν in the range from −200 MHz to −0.2 MHz to highlight
the rates to shallow bound levels. The length of the free-
atom quantization box is L = 1 mm. The temperature of the
phonon bath is T = 300 K. Other parameters are as in Fig. 3.
We plot in Fig. 8 the free-to-bound transition rate Gνf
[see Eq. (43)] as a function of the vibrational quantum
number ν. The upper (lower) part of the figure corre-
sponds to the case of the initial-state energy Ef = 2
MHz (Ef = 3.1 THz), which is close to the average ki-
netic energy per atom in an ideal gas with temperature
T0 = 200 µK (T0 = 300 K). We observe that the free-to-
bound transition rate first increases and then decreases
with increasing transition frequency ωfν = (Ef − Eν)/h̄.
Such behavior results from the competing effects of the
mean phonon number, the phonon mode density, and the
matrix element of the force, like in the case of bound-
to-bound transitions (see Fig. 3). We also see a cut-
off of the transition frequency, which is associated with
the Debye frequency. Comparison of Figs. 8(a) and 8(b)
shows that the transitions from low-energy free states
have probabilities orders of magnitude smaller than those
from high-energy free states. One of the reasons is that
the transition rate Gνf is proportional to the velocity
vf = (2Ef/m)1/2 [see Eq. (43)]. The dependence of
the transition rate density Rνf on the transition fre-
quency ωfν also plays an important role. Because of this,
the rates for the transitions from low-energy free states
to shallow bound levels are very small [see the inset of
Fig. 8(a)].
FIG. 9: Free-to-bound decay rate Gf as a function of the
free-state energy Ef . The inset highlights the magnitude and
profile of the decay rate for Ef in the range from 0 to 20
MHz. The temperature of the phonon bath is T = 300 K.
Other parameters are as in Fig. 8.
We show in Fig. 9 the free-to-bound decay rate Gf
[see Eq. (44)], which is a characteristic of the adsorp-
tion process, as a function of the free-state energy Ef .
We see that Gf first increases and then decreases with
increasing Ef . The increase of Gf with increasing Ef in
the region of small Ef (see the inset) is mainly due to
the increase in the atomic incidence velocity vf . In this
region, we have Gf ∝ vf ∝
Ef [see Eqs. (43) and
(44)]. For Ef in the range from 0 to 20 MHz, which is
typical for atoms in magneto-optical traps, the maximum
value of Gf is on the order of 10
4 s−1 (see the inset of
Fig. 9). Such free-to-bound (adsorption) rates are sev-
eral orders of magnitude smaller than the bound-to-free
(desorption) rates (see the dashed curve in Fig. 5). The
decrease of Gf with increasing Ef in the region of large
Ef is mainly due to the reduction of the atom–phonon
coupling coefficients.
FIG. 10: Free-to-bound transition rates GνT0 for transitions
from the thermal states with temperatures (a) T0 = 200 µK
and (b) T0 = 300 K to bound levels ν as functions of the
bound-level energy Eν . The insets show GνT0 on the log scale
versus Eν in the range from −200 MHz to −0.2 MHz to high-
light the rates to shallow bound levels. The temperature of
the phonon bath is T = 300 K. Other parameters are as in
Fig. 8.
FIG. 11: Free-to-bound decay rate GT0 as a function of the
atomic temperature T0 in the ranges (a) from 100 µK to 400
µK and (b) from 50 K to 350 K. The temperature of the
phonon bath is T = 300 K. Other parameters are as in Fig. 8.
In a thermal gas, the adsorption process is charac-
terized by the transition rate GνT0 [see Eq. (46)] and
the decay rate GT0 [see Eq. (47)], which are the av-
erages of the free-to-bound transition rate Gνf and the
free-to-bound decay rate Gf , respectively, over the free-
state energy distribution (45). We plot the free-to-bound
transition rate GνT0 and the free-to-bound decay rate
GT0 in Figs. 10 and 11, respectively. Comparison be-
tween Figs. 10(a) and 9(a) shows that the transition rates
from low-temperature thermal states and low-energy free
states look quite similar to each other. The reason is that
the spread of the energy distribution is not substantial
in the case of low temperatures. The spread of the en-
ergy distribution is however substantial in the case of
high temperatures, leading to the softening of the cut-
off frequency effect [compare Fig. 10(b) with Fig. 9(b)].
Figure 11 shows that the free-to-bound decay rate GT0
first increases and then reduces with increasing atomic
temperature T0. For T0 in the range from 100 µK to 400
µK, which is typical for atoms in magneto-optical traps,
the maximum value of GT0 is on the order of 10
4 s−1 [see
Fig. 11(a)]. Such free-to-bound (adsorption) rates are
several orders of magnitude smaller than the bound-to-
free (desorption) rates (see the dashed curve in Fig. 5).
Figure 11(a) shows that, in the region of low atomic tem-
perature T0, one has GT0 ∝
T0, in agreement with the
asymptotic behavior of Eqs. (46) and (47).
FIG. 12: Free-to-free transition rate densities Qf ′f for the
upward (solid lines) and downward (dashed lines) transitions
from the free states |f〉 with energies (a) Ef = 2 MHz and (b)
Ef = 3.1 THz to other free states |f
′〉 as functions of the final-
level energy Ef ′ . The arrows mark the energies of the initial
free states. The inset in part (a) shows Qf ′f versus Ef ′ in
the range from 0 to 4 MHz to highlight the small magnitude
of the rate density for downward transitions (dashed line).
The temperature of the phonon bath is T = 300 K. Other
parameters are as in Fig. 8.
We now examine free-to-free transitions, both upward
and downward, which corresponding to the heating and
cooling processes of free atoms by the surface. We plot in
Fig. 12 the free-to-free transition rate density Qf ′f [see
Eq. (48)] as a function of the final-level energy Ef ′ . The
upper (lower) part of the figure corresponds to the case
of the initial-state energy Ef = 2 MHz (Ef = 3.1 THz),
which is close to the average kinetic energy per atom in
an ideal gas with temperature T0 = 200 µK (T0 = 300
K). The rate densities are shown for the upward (phonon-
absorption) and downward (phonon-emission) transitions
by the solid and dashed lines, respectively. The fig-
ure shows that the free-to-free transition rate density in-
creases or decreases with increasing transition frequency
if the latter is not too large or is large enough, respec-
tively. We also observe a signature of the Debye cutoff
of the phonon frequency. Comparison of Figs. 12(a) and
12(b) shows that transitions from low-energy free states
have probabilities orders of magnitude smaller than those
from high-energy free states. Figure 12(a) and its inset
show that, when the energy of the free state is low, the
free-to-free downward (cooling) transition rate is very
small as compared to the free-to-free upward (heating)
transition rate.
FIG. 13: Free-to-free upward and downward decay rates Qaf
(solid lines) and Qef (dashed lines) as functions of the energy
Ef of the initial free state. The insets highlight the magni-
tudes and profiles of the decay rates for Ef in the range from
0 to 20 MHz. The temperature of the phonon bath is T = 300
K. Other parameters are as in Fig. 8.
We show in Fig. 13 the free-to-free upward (phonon-
absorption) and downward (phonon-emission) decay
rates Qaf [see Eq. (50)] and Q
f [see Eq. (49)] as functions
of the free-state energy Ef . We observe that Qaf and Qef
increase with increasing Ef in the range from 0 to 8 THz.
The increase of Qaf with increasing Ef in the region of
small Ef (see the left inset) is mainly due to the increase
in the atomic incidence velocity vf . In this region, we
have Qaf ∝ vf ∝
Ef [see Eqs. (48) and (50)]. The
increase of Qef with increasing Ef in the region of small
Ef (see the right inset) is due to not only the increase in
the atomic incidence velocity vf [see Eq. (48)] but also
the increase of the transition rate density Qef ′f and the
increase of the integration interval (0, Ef) [see Eq. (49)].
In this region, the dependence of Qef on the energy Ef
is of higher order than E3/2f . The left inset of Fig. 13
shows that, for Ef in the range from 0 to 20 MHz, the
maximum value of Qaf is on the order of 10
4 s−1. Such
free-to-free upward (heating) decay rates are comparable
to but about two times smaller than the corresponding
free-to-bound (adsorption) decay rates (see the inset of
Fig. 9). Meanwhile, the right inset of Fig. 13 shows that,
in the region of small Ef , the free-to-free downward (cool-
ing) decay rate Qef is very small.
FIG. 14: Free-to-free transition rate densities QafT0 for up-
ward transitions (solid lines) and QefT0 for downward tran-
sitions (dashed lines) from the thermal states with tempera-
tures (a) T0 = 200 µK and (b) T0 = 300 K to free levels f
as functions of the free-level energy Ef . The inset in part (a)
shows the rate densities versus Ef in the range from 0 to 8
MHz to highlight the small magnitude of QefT0 (dashed line).
The temperature of the phonon bath is T = 300 K. Other
parameters are as in Fig. 8.
FIG. 15: Free-to-free decay rates QaT0 (solid lines) and
(dashed lines) for upward and downward transitions, re-
spectively, as functions of the atomic temperature T0 in the
ranges (a) from 100 µK to 400 µK and (b) from 50 K to 350
K. For comparison, the free-to-bound decay rate GT0 is re-
plotted from Fig. 11 by the dotted lines. The temperature of
the phonon bath is T = 300 K. Other parameters are as in
Fig. 8.
In the case of a thermal gas, the phonon-mediated heat
transfer between the gas and the surface is characterized
by the free-to-free transition rate densitiesQafT0 andQ
[see Eqs. (51)] and the free-to-free decay rates QaT0and
QeT0 [see Eqs. (52)]. We plot the free-to-free transition
rate densities QafT0 and Q
in Fig. 14. Comparison
between Figs. 14(a) and 12(a) shows that the transition
rate densities from low-temperature thermal states and
low-energy free states are quite similar to each other.
The spread of the initial-state energy distribution is not
substantial in this case. However, the energy spread of
the initial state is substantial in the case of high tem-
peratures, concealing the cutoff frequency effect [com-
pare Fig. 14(b) with Fig. 12(b)]. We display the free-
to-free decay rates QaT0 and Q
in Fig. 15. The solid
and dashed lines correspond to the upward (heating) and
downward (cooling) transitions, respectively. For com-
parison, the free-to-bound decay rate (adsorption rate)
GT0 is re-plotted from Fig. 11 by the dotted lines. We
observe that, for T0 in the range from 100 µK to 400 µK
[see Fig. 15(a)], the adsorption rate GT0 (dotted line) is
about two times larger than the heating rate QaT0 (solid
line), while the cooling rate QeT0 (dashed line) is negligi-
ble. Figure 15(a) shows that, in the region of low atomic
temperatures, one has QT0
∼= QaT0 ∝
T0, in agreement
with the asymptotic behavior of expressions (52). The
figure also shows that QeT0 quickly increases with increas-
ing atomic temperature T0. The relation Q
< QaT0 ,
obtained for T0 < T , indicates the dominance of heating
of cold free atoms by the surface. The substantial mag-
nitude of the free-to-bound transition rate GT0 (dotted
line) indicates that a significant number of atoms can be
adsorbed by the surface. According to Fig. 15(b), the
free-to-free downward transition rate QeT0 (dashed line)
crosses the upward transition rate QaT0 (solid line) when
T0 = T = 300 K, and then becomes the dominant decay
rate. The relation QeT0 > Q
, obtained for T0 > T , indi-
cates the dominance of cooling of hot free atoms by the
surface.
V. CONCLUSIONS
In conclusion, we have studied the phonon-mediated
transitions of an atom in a surface-induced potential.
We developed a general formalism, which is applicable
for any surface–atom potential. A systematic derivation
of the corresponding density-matrix equation enables us
to investigate the dynamics of both diagonal and off-
diagonal elements. We included a large number of vi-
brational levels originating from the deep silica–cesium
potential. We calculated the transition and decay rates
from both bound and free levels. We found that the
rates of phonon-mediated transitions between transla-
tional levels depend on the mean phonon number, the
phonon mode density, and the matrix element of the force
from the surface upon the atom. Due to the effects of
these competing factors, the transition rates usually first
increase and then reduce with increasing transition fre-
quency. We focused on the transitions from bound states.
Two specific examples, namely, when the initial level is
a shallow level also when it can be one of the deep levels
have been worked out. We have shown that there can be
marked differences in the absorption and emission behav-
ior in the two cases. For example, both the absorption
and emission rates from the deep bound levels can be sev-
eral orders (in our case, six orders) of magnitude larger
than the corresponding rates from the shallow bound lev-
els. We also analyzed various types of transitions from
free states. We have shown that, for thermal atomic ce-
sium with temperature in the range from 100 µK to 400
µK in the vicinity of a silica surface with temperature of
300 K, the adsorption (free-to-bound decay) rate is about
two times larger than the heating (free-to-free upward de-
cay) rate, while the cooling (free-to-free downward decay)
rate is negligible.
Acknowledgments
We thank M. Chevrollier for fruitful discussions. This
work was carried out under the 21st Century COE pro-
gram on “Coherent Optical Science.”
[∗] Also at Institute of Physics and Electronics, Vietnamese
Academy of Science and Technology, Hanoi, Vietnam.
[1] V. I. Balykin, K. Hakuta, Fam Le Kien, J. Q. Liang, and
M. Morinaga, Phys. Rev. A 70, 011401(R) (2004); Fam
Le Kien, V. I. Balykin, and K. Hakuta, Phys. Rev. A 70,
063403 (2004).
[2] Fam Le Kien, S. Dutta Gupta, V. I. Balykin, and K.
Hakuta, Phys. Rev. A 72, 032509 (2005).
[3] Fam Le Kien, S. Dutta Gupta, K. P. Nayak, and K.
Hakuta, Phys. Rev. A 72, 063815 (2005).
[4] E. G. Lima, M. Chevrollier, O. Di Lorenzo, P. C. Se-
gundo, and M. Oriá, Phys. Rev. A 62, 013410 (2000).
[5] T. Passerat de Silans, B. Farias, M. Oriá, and M.
Chevrollier, Appl. Phys. B 82, 367 (2006).
[6] Fam Le Kien and K. Hakuta, Phys. Rev. A 75, 013423
(2007).
[7] Fam Le Kien, S. Dutta Gupta, and K. Hakuta, e-print
quant-ph/0610067.
[8] K. P. Nayak, P. N. Melentiev, M. Morinaga, Fam Le Kien,
V. I. Balykin, and K. Hakuta, e-print quant-ph/0610136.
[9] C. Henkel and M. Wilkens, Europhys. Lett. 47, 414
(1999).
[10] Z. W. Gortel, H. J. Kreuzer, and R. Teshima, Phys. Rev.
B 22, 5655 (1980).
[11] H. Hoinkes, Rev. Mod. Phys. 52, 933 (1980).
[12] N. N. Bogolubov, Commun. of JINR, E17-11822, Dubna
(1978); N. N. Bogolubov and N. N. Bogolubov Jr., Ele-
mentary Particles and Nuclei (USSR) 11, 245 (1980).
[13] R. Zwanzig, Lectures in Theoretical Physics, eds. W. E.
Brittin, B. W. Downs, and J. Downs (Interscience, New
York, 1961) Vol. 3, p. 106; G. S. Agarwal, Progress in
Optics, ed. E. Wolf (North-Holland, Amsterdam, 1973)
Vol. 11, p. 3; L. Mandel and E. Wolf, Optical Coherence
and Quantum Optics (Cambridge, New York, 1995) p.
[14] J. Javanainen and M. Mackie, Phys. Rev. A 58, R789
(1998); M. Mackie and J. Javanainen, ibid. 60, 3174
(1999).
[15] E. Luc-Koenig, M. Vatasescu, and F. Masnou-Seeuws,
Eur. Phys. J. D 31, 239 (2004).
[16] See, for example, G. P. Agrawal, Nonlinear Fiber Optics
(Academic, New York, 2001).
[17] H. J. Metcalf and P. van der Straten, Laser Cooling and
Trapping (Springer, New York, 1999).
http://arxiv.org/abs/quant-ph/0610067
http://arxiv.org/abs/quant-ph/0610136
|
0704.0341 | Infrared Evolution Equations: Method and Applications | Infrared Evolution Equations: Method and Applications
B.I. Ermolaev
Ioffe Physico-Technical Institute, 194021 St.Petersburg, Russia
M. Greco
Department of Physics and INFN, University Rome III, Rome, Italy
S.I. Troyan
St.Petersburg Institute of Nuclear Physics, 188300 Gatchina, Russia
It is a brief review on composing and solving Infrared Evolution Equations. They can be used in
order to calculate amplitudes of high-energy reactions in different kinematic regions in the double-
logarithmic approximation.
PACS numbers: 12.38.Cy
I. INTRODUCTION
Double-logarithmic (DL) contributions are of a special interest among radiative corrections. They are interesting
in two aspects: first, in every fixed order of the perturbation theories they are the largest terms among the radiative
corrections depending on the total energy and second, they are easiest kind of the corrections to sum up. DL corrections
were discovered by V.V. Sudakov in Ref. [1] in the QED context. He showed that DL terms appear from integrations
over soft, infrared (IR) -divergent momenta of virtual photons. All-order resummation of such contributions led to
their exponentiations.
Next important step was done in Refs. [2] where calculation and summation of DL contributions was considered
in a systematic way. They found a complementary source of DL terms: soft virtual fermions. This situation appears
in the Regge kinematics. The all-order resummations of DL contributions in the Regge kinematic are quite involved
and yield more complicated expressions than the Sudakov exponentials. Nonetheless important was the proof of the
factorization of bremsstrahlung photons with small k⊥ in the high-energy hadronic reactions found in Ref. [3] and
often addressed as the Gribov’s bremsstrahlung theorem. This statement, suggested originally in the framework of
the phenomenological QED of hadrons was extended to QCD in Refs. [4].
Calculation in the double-logarithmic approximation (DLA) amplitudes of the fermion-antifermion annihilation in
the Regge forward and backward kinematics involves accounting for DL contributions from soft quarks and soft gluons.
These reactions in QED and QCD have many common features. The e+e− -annihilation was studied in Refs. [2].
The quark-aniquark annihilation DLA was investigated in Ref. [5]. The method of calculation here was based on
factorization of virtual quarks and gluons with minimal k⊥. Generally speaking, the results obtained in Ref. [5] could
be obtained with the method of Ref. [2], however the technique of calculations suggested in Ref. [5] was much more
elegant and efficient. Although Ref. [5] is about quark scattering only, it contains almost all technical ingredients
necessary to compose Infrared Evolution Equations for any of elastic scattering amplitudes. Nevertheless it could not
directly be applied to inelastic processes involving emission of soft particles. Such a generalization was obtained in
Refs. [4, 6]. The basic idea of the above-mentioned method was suggested by L.N. Lipatov: to investigate evolution
with respect to the infrared cut-off. The present, sounding naturally term ”Infrared Evolution Equations” (IREE) for
this method was suggested by M. Krawczyk in Ref. [7] where amplitudes for the backward Compton scattering were
calculated in DLA.
The aim of the present brief review is to show how to compose and solve IREE for scattering amplitudes in different
field theories and kinematic regions. The paper is organized as follows: in Sect. II we consider composing IREE in the
technically simplest hard kinematics. In Sect. III we consider composing IREE in the forward kinematics and apply
it to studying the structure function g1 of the polarized Deep-Inelastic scattering (DIS) at small x. The point is that
the commonly used theoretical instrument to study g1 is DGLAP [11]. It collects logarithms of Q
2 to all orders in
αs but does not include the total resummation of logarithms of 1/x, though it is important at small x. Accounting
for such a resummaton leads to the steep rise of g1 at the small-x region. As is shown in Sect. IV, DGLAP lacks the
resummaion but mimics it inexplicitly, through the special choice of fits for the initial parton densities. Invoking such
peculiar fits together with DGLAP to describe g1 at x ≪ 1 led to various misconceptions in the literature. They are
enlisted and corrected in Sect. V. The total resummaion of the leading logarithms is essential in the region of small x.
In the opposite region of large x, DGLAP is quite efficient. It is attractive to combine the resummation with DGLAP.
http://arxiv.org/abs/0704.0341v1
The manual for doing it is given in Sect. VI. Finally, Sect. VII is for concluding remarks.
II. IREE FOR SCATTERING AMPLITUDES IN THE HARD KINEMATICS
From the technical point of view, the hard kinematics, where all invariants are of the same order, is the easiest for
analysis. For the simplest, 2 → 2 -processes, the hard kinematics means that the Mandelstamm variables s, t, u obey
s ∼ −t ∼ −u . (1)
In other words, the cmf scattering angles θ ∼ 1 in the hard kinematics. This kinematics is the easiest because the
ladder Feynman graphs do not yield DL contributions here and usually the total resummation of DL contributions
leads to multiplying the Born amplitude by exponentials decreasing with the total energy. Let us begin with composing
and solving an IREE for the well-known object: electromagnetic vertex Γµ of an elementary fermion (lepton or quark).
As is known,
Γµ = ū(p2)
γµf(q
2)− σµνqν
g(q2)
u(p1) (2)
where p1,2 are the initial and final momenta of the fermion, m stands for the fermion mass and the transfer momentum
q = p2−p1. Scalar functions f and g in Eq. (2) are called form factors. Historically, DL contributions were discovered
by V. Sudakov when he studied the QED radiative corrections to the form factor f at |q2| ≫ |p21,2|. Following him,
let us consider vertex Vµ at
|q2| ≫ p21 = p22 = m2 (3)
i.e. we assume the fermion to be on–shell and account for DL electromagnetic contributions. We will drop m for the
sake of simplicity.
A. IREE for the form factor f(q2) in QED
Step 1 is to introduce the infrared cut-off µ in the transverse (with respect to the plane formed by momenta p1,2)
momentum space for all virtual momenta ki:
ki ⊥ > µ (4)
where i = 1, 2, ...
Step 2 is to look for the softest virtual particle among soft external and virtual particles. The only option we have
is the softest virtual photon. Let denote its transverse momenta ≡ k⊥. By definition,
k⊥ = min ki ⊥ . (5)
Step 3: According to the Gribov theorem, the propagator of the softest photon can be factorized (i.e. it is attached
to the external lines in all possible ways) whereas k⊥ acts as a new cut-off for other integrations. Adding the Born
contribution fBorn = 1 we arrive at the IREE for f in the diagrammatic form. It is depicted in Fig. 1. IREE in the
analytic form are written in the gauge-invariant way, but their diagrammatical writing depends on the gauge. In the
present paper we use the Feynman gauge.
Applying to it the standard Feynman rules, we write it in the analytic form:
f(q2, µ2) = fBorn − e
dαdβdk2
− µ2) f(q2, k2
(sαβ − k2
+ ıǫ)(−sα+ sαβ − k2
+ ıǫ)(sβ + sαβ − k2
+ ıǫ)
where we have used the Sudakov parametrization k = αp2 + βp1 + k⊥ and denoted s = −q2 ≈ 2p1p2. As f(q2, k2⊥)
does not depend on α and β, the DL integration over them can be done with the standard way, so we are left with a
simple integral equation to solve:
f(q2, µ2) = fBorn − e
ln(s/k2
)f(q2, k2
) . (7)
FIG. 1: The IREE for the Sudakov form factor. The letters in the blobs stand for IR cut-off.
Differentiation of Eq. (7) over µ2 (more exactly, applying −µ2∂/∂µ2) reduces it to a differential equation
∂f/∂(ln(s/µ2)) = −(e2/8π2) ln(s/µ2)f (8)
with the obvious solution
f = fBorn exp[−(α/4π) ln2(q2/m2)] (9)
where we have replaced µ by m and used α = e2/4π. Eq. (9) is the famous Sudakov exponential obtained in Ref. [1].
B. IREE for the form factor g(q2) in QED
Repeating the same steps (see Ref. [8] for detail) leads to a similar IREE for the form factor g:
g(q2,m2, µ2) = gBorn(s,m2)− e
ln(s/k2
)g(q2,m2, k2
) (10)
where gBorn(s,m2) = −(m2/s)(α/π) ln(s/m2). Solving this equation and putting µ = m in the answer leads to the
following relation between form factors f and g:
g(s) = −2
, (11)
with ρ = s/m2. Combining Eqs. (9,11) allows to write a simple expression for the DL asymptotics of the vertex Γµ:
Γµ = ū(p2)
σµνqν
u(p1) exp[−(α/4π) ln2 ρ] . (12)
C. e+e− -annihilation into a quark-antiquark pair
Let us consider the e+e− -annihilation into a quark q(p1) and q̄(p2) at high energy when 2p1p2 ≫ p21,2. We consider
the channel where the e+e− -pair annihilates into one heavy photon which decays into the q(p1) q̄(p2) -pair:
e+e− → γ∗ → q(p1) q̄(p2) . (13)
We call this process elastic. In this case the most sizable radiative corrections arise from the graphs where the quark
and antiquark exchange with gluons and these graphs look absolutely similar to the graphs for the electromagnetic
vertex Γµ considered in the previous subsection. As a result, accounting for the QCD radiative corrections in DLA
to the elastic form factors fq, gq of quarks can be obtained directly from Eqs. (9,11) by replacement
α → αsCF , (14)
with CF = (N
2 − 1)/2N = 4/3.
D. e+e− -annihilation into a quark-antiquark pair and gluons
In addition to the elastic annihilation (13), the final state can include gluons:
e+e− → γ∗ → q(p1) q̄(p2) + g(k1), ..g(kn) . (15)
We call this process the inelastic annihilation. The QED radiative corrections to the inelastic annihilation (15) in
DLA are absolutely the same as the corrections to the elastic annihilation. On the contrary, the QCD corrections
account for gluon exchanges between all final particles. This makes composing the IREE for the inelastic annihilation
be more involved (see Ref. [4]). The difference to the considered elastic case appears at Step 2: look for the softest
virtual particle among soft external and virtual particles. Indeed, now the softest particle can be both a virtual gluon
and an emitted gluon. For the sake of simplicity let us discuss the 3-particle final state, i.e. the process
e+e− → γ∗ → q(p1) q̄(p2) + g(k1) . (16)
The main ingredient of the scattering amplitude of this process is the new electromagnetic vertex Γ
µ of the quark.
In DLA, it is parameterized by new form factors F (1) and G(1)
Γµ = B1(k1)ū(p2)
(1)(q, k1)−
σµνqν
G(1)(q, k1)
u(p1) (17)
where (1) corresponds to the number of emitted gluons, q = p1 + p2 and l is the polarization vector of the emitted
gluon. The bremsstrahlung factor B1 in Eq. (17) at high energies is expressed through k1 ⊥:
( p2l
− p1l
. (18)
We call F (n), G(n) inelastic form factors. Let us start composing the IREE for F (1). Step 1 is the same like in
the previous case. Step 2 opens more options. Let us first choose the softest gluon among virtual gluons and denote
its transverse momentum k⊥ The integration over k⊥ runs from µ to s. As µ < k1 ⊥ < s, we have two regions to
consider: Region D1 were
µ < k1⊥ < k⊥ <
s (19)
and Region D2 were
µ < k⊥ < k1⊥ <
s (20)
Obviously, the softest particle in Region D1 is the emitted gluon, so it can be factorized as depicted in graphs (b,b’)
of Fig. 2.
On the contrary, the virtual gluon is the softest in Region D2 were its propagator is factorized as shown in graphs
(c,d,d’) of Fig. 2. Adding the Born contribution (graphs (a,a’) in Fig. 2) completes the IREE for F (1) depicted in
Fig. 2. Graphs (a-b’) do not depend on µ and vanish when differentiated with respect to µ. Blobs in graphs (c-d’) do
not depend on the longitudinal Sudakov variables, so integrations over α, β can be done like in the first loop. After
that the differential IREE for F (1) is
∂F (1)
CF ln
(2p2k1
(2p1k1
F (1) . (21)
Solving Eq. (21) and using that (2p1k1)(2p2k1) = sk
1⊥ leads to the expression
F (1) = exp
CF ln
(k21⊥
suggested in Ref. [9] and proved in Ref. [4] for any n. The IREE for the form factor G(n) was obtained and solved in
Ref. [8]. It was shown that
G(n) = −2∂F (n)/∂ρ . (23)
FIG. 2: The IREE for the inelastic quark form factor.
E. Exponentiation of Sudakov electroweak double-logarithmic contributions
The IREE -method was applied in Ref. [10] to prove exponentiation of DL correction to the electroweak (EW)
reactions in the hard kinematics. There is an essential technical difference between the theories with the exact gauge
symmetry (QED and QCD) and the EW interactions theory with the broken SU(2) ⊗ U(1) gauge symmetry: only
DL contributions from virtual photons yield IR singularities needed to be regulated with the cut-off µ whereas DL
contributions involving W and Z -bosons are IR stable because the boson masses MW and MZ act as IR regulators.
In Ref. [10] the difference between MW and MZ was neglected and the parameter
M & MW ≈ MZ (24)
was introduced, in addition to µ, as the second IR cut-off. It allowed to drop masses MW,Z . The IREE with two IR
cut-offs was composed quite similarly to Eq. (6), with factorizing one by one the softest virtual photon, Z-boson and
W -boson. As a result the EW Sudakov form factor FEW is
FEW = exp
− α(Q
ln2(s/µ2)−
SU(2)
(Y 21 + Y
− α(Q
ln2(s/M2)
where Q1,2 are the electric charges of the initial and final fermion (with W -exchanges accounted, they may be
different), Y1,2 are their hyper-charges and C
SU(2)
F = (N
2 − 1)/2N , with N = 2. We have used in Eq. (25) the
standard notations g and g′ for the SU(2) and U(1) -EW couplings. The structure of the exponent in Eq. (25) is quite
clear: the first, µ -dependent term comes from the factorization of soft photons like the exponent in Eq. (9) while
other terms correspond to the W and Z -factorization; the factor in the squared brackets is the sum of the SU(2)
and U(1) Casimirs, with the photon Casimir being subtracted to avoid the double counting. In the limit µ = M the
group factor in the exponent is just the Casimir of SU(2)⊗ U(1).
III. APPLICATION OF IREE TO THE POLARIZED DEEP-INELASTIC SCATTERING
Cross-sections of the polarized DIS are described by the structure functions g1,2. They appear from the standard
parametrization of the spin-dependent part Wµν of the hadronic tensor:
Wµν = ıǫµνλρqλ
Sρg1(x,Q
Sρ − pρ
g2(x,Q
where p, m and S are the momentum, mass and spin of the incoming hadron; q is the virtual photon momentum;
Q2 = −q2; x = Q2/2pq. Obviously, Q2 > 0 and 0 6 x 6 1.
Unfortunately, g1,2 cannot be calculated in a straightforward model-independent way because it would involve QCD
at long distances. To avoid this problem, Wµν is regarded as a convolution of Φq,g - probabilities to find a polarized
quark or gluon and the partonic tensors W̃
(q,g)
µν parameterized identically to Eq. (26). In this approach W̃
(q,g)
µν involve
only QCD at short distances, i.e. the Perturbative QCD while long-distance effects are accumulated in Φq,g. As Φq,g
are unknown, they are mimicked by the initial quark and gluon densities δq, δg. They are fixed aposteriori from
phenomenological considerations. So, the standard description of DIS is:
Wµν ≈ W (q)µν ⊗ δq +W (g)µν ⊗ δg . (27)
The standard theoretical instrument to calculate g1 is DGLAP[11] complemented with standard fits[12] for δq, δg.
We call it Standard Approach (SA). In this approach
g1(x,Q
2) = Cq(x/z)⊗∆q(z,Q2) + Cg(x/z)⊗∆g(z,Q2) (28)
where Cq, g are coefficient functions and ∆q(z,Q2), ∆g(z,Q2) are called the evolved (with respect to Q2)quark and
gluon distributions. They are found as solutions to DGLAP evolution equations
d lnQ2
Pqq∆q + Pqg∆g
d lnQ2
Pgq∆q + Pgg∆g
where Pab are the splitting functions. The Mellin transforms γab of Pab are called the DGLAP anomalous dimensions.
They are known in the leading order (LO) where they are ∼ αs and in the next-to-leading order (NLO), i.e. ∼ α2s.
Similarly, Cq,g are known in LO and NLO. Details on this topic can be found in the literature (e.g. see a review
[13]). Structure function g1 has the flavor singlet and non-singlet components, g
1 and g
1 . Expressions for g
1 are
simpler, so we will use mostly them in the present paper when possible. It is convenient to write g1 in the form of the
Mellin integral. In particular,
gNS DGLAP1 (x,Q
2) = (e2q/2)
CNS(ω)δq(ω) exp
γNS(ω, αs(k
where µ2 is the starting point of the Q2 -evolution; CNS and γNS are the non-singlet coefficient function and anomalous
dimension. In LO
γNS(ω,Q
ω(1 + ω)
+ S2(ω)
, (31)
CLONS(ω) = 1 +
2ω + 1
ω(1 + ω)
S1(ω) + S
1(ω)− S2(ω)
FIG. 3: The IREE for the non-singlet component of the spin structure function g1.
with Sr(ω) =
j=1 1/j
r . The initial quark and gluon densities in Eq. (30) are defined through fitting experimental
data. For example, the fit for δq taken from the first paper in Ref. [12] is
δq(x) = Nx−α
(1− x)β(1 + γxδ)
, (32)
with N being the normalization, α = 0.576, β = 2.67, γ = 34.36 and δ = 0.75.
DGLAP equations were suggested for describing DIS in the region
x . 1, Q2 ≫ µ2 (33)
(µ stands for a mass scale, µ ≫ ΛQCD) and there is absolutely no theoretical grounds to apply them in the small-x
region, however being complemented with the standard fits they are commonly used at small x. It is known that SA
provide a good agreement with available experimental data but the price is invoking a good deal of phenomenological
parameters. The point is that DGLAP, summing up leading lnk Q2 to all orders in αs, cannot do the same with
leading lnk(1/x). The later is not important in the region (33) where lnk(1/x) ≪ 1 but becomes a serious drawback
of the method at small x. The total resummation of DL contributions to g1 in the region
x ≪ 1, Q2 ≫ µ2 (34)
was done in Refs. [14]. The weakest point in those papers was keeping αs as a parameter, i.e. fixed at an unknown
scale. Accounting for the most important part of single-logarithmic contributions, including the running coupling
effects were done in Refs. [15]. In these papers µ2 was treated as the starting point of the Q2 -evolution and as the IR
cut-off at the same time. The structure function g1 was calculated with composing and solving IREE in the following
It is convenient to compose IREE not for g1 but for forward (with |t| . µ2) Compton amplitude M related to g1
as follows:
ℑM . (35)
It is also convenient to use for amplitude M the asymptotic form of the Sommerfeld-Watson transform:
ξ(−)(ω)F (ω,Q2/µ2) (36)
where ξ(−)(ω) = [e−ıπω − 1]/2 ≈ −ıπω/2 is the signature factor. The transform of Eq. (36) and is often addressed
as the Mellin transform but one should remember that it coincides with the Mellin transform only partly. IREE for
Mellin amplitudes F (ω,Q2) look quite simple.
For example, the IREE for the non-singlet Mellin amplitude FNS related to gNS1 by Eqs. (35,36) is depicted in
Fig. 3. In the Mellin space it takes the simple form:
[ω + ∂/∂y]FNS = (1 + ω/2)HNSF
NS (37)
where y = ln(Q2/µ2). Eq. (37) involves a new object (the lowest blob in the last term in Fig. 3): the non-singlet
anomalous dimension HNS accounting for the total resummaton of leading logarithms of 1/x. Like in DGLAP, the
anomalous dimension does not depend on Q2 but, in contrast to DGLAP, HNS can be found with the same method.
The IREE for it is algebraic:
ωHNS = A(ω)CF /8π
2 + (1 + ω/2)H2NS +D(ω)/8π
2 . (38)
The system of Eqs. (37,38) can be easily solved but before doing it let us comment on them. The left-hand sides
of Eqs. (37,38) are obtained with applying the operator −µ2∂/∂µ2 to Eq. (36). The Born contribution in Fig. 3
does not depend on µ and therefore vanishes. The last term in Fig. 3 (the rhs of Eq. (37)) is the result of a new, t
-channel factorization which does not exist in the hard kinematics defined in Eq. (1). In order to compose the IREE
for the Compton amplitude M , in accordance with the prescription in the previous section we should first introduce
the cut-off µ. Then Step 2 is to tag the softest particles. In the case under discussion we do not have soft external
particles. Had the softest particle been a gluon, it could be factorized in the same way like in Sect. II. However, the
only option now is to attach the softest propagator to the external quark lines and get ln(t/µ2) = 0 from integration
over β (cf Eq. (7)). So, the softest gluon does not yield DL contributions. The other option is to find a softest quark.
The softest t -channel quark pair factorizes amplitude M into two amplitudes (the last term in Fig. 3) and yield DL
contributions. The IREE for HNS is different:
(i) HNS does not depend on Q
2, so there is not a derivative in the lhs of Eq. (37).
(ii) The Born term depends on µ and contributes to the IREE (term A in Eq. (37))).
(iii) As all external particles now are quarks, the softest virtual particle can be both a quark and gluon. The case
when it is the t -channel quark pair, corresponds to the quadratic term in the rhs of Eq. (37). The case of the softest
gluon yields the term D, with
D(ω) =
dηe−ωη ln
(ρ+ η
[ ρ+ η
(ρ+ η)2 + π2
where b = (33− 2nf )/12π and η = ln(µ2/Λ2QCD).
The term A in Eq. (37) stands instead of αs. The point is that the standard parametrization αs = αs(Q
2) cannot
be used at x ≪ 1 and should be changed (see Ref. [16] for detail). It leads to the replacement αs by
A(ω) =
η2 + π2
dρe−ωρ
(ρ+ η)2 + π2
. (40)
Having solved Eqs. (37,38), we arrive at the following expression for gNS1 in the region (34):
gNS1 (x,Q
2) = (e2q/2)
(1/x)ωCNS(ω)δq(ω) exp
HNS(ω)y
where the coefficient function CNS(ω) is expressed through HNS(ω):
CNS(ω) =
ω −HNS(ω)
and HNS(ω) is the solution of algebraic equation (43):
HNS = (1/2)
ω2 −B(ω)
where
B(ω) = (4πCF (1 + ω/2)A(ω) +D(ω))/(2π
2) . (44)
It is shown in Ref. [17] that the expression for g1 in the region
x ≪ 1, Q2 . µ2 (45)
can be obtained from the expressions obtained in Refs. [15] for g1 in region (34) by the shift
Q2 → Q2 + µ20 (46)
where µ0 = 1 GeV for the non-singlet g1 and µ0 = 5.5 GeV for the singlet.
IV. COMPARISON OF EXPRESSIONS (30) AND (41) FOR gNS1
Eqs. (30) and (41) read that the non-singlet g1 is obtained from δq with evolving it with respect to x (using the
coefficient function) and with respect to Q2 (using the anomalous dimension). Numerical comparison of Eqs. (30)
and (41) can be done when δq is specified.
A. Comparison of small-x asymptotics, neglecting the impact of δq
In the first place let us compare the small-x asymptotics of for gNS DGLAP1 and g
1 , assuming that δq does not
affect them. In other words, we compare the differencee in the x-evolution at x → 0. Applying the saddle-point
method to Eqs. (30) and (41) leads to the following expressions:
gNS DGLAP1 ∼ exp
ln(1/x) ln ln(Q2/Λ2QCD)
gNS1 ∼ (1/x)∆NS(Q2/µ2)∆NS/2 (48)
where ∆NS = 0.42 is the non-singlet intercept
1. Expression (47) is the well-known DGLAP asymptotics. Obviously,
the asymptotics (48) is much steeper than the DGLAP asymptotics (30).
B. Numerical comparison between Eqs. (30) and (41), neglecting the impact of δq
A comparison between Eqs. (30) and (41) strongly depends on the choice of δq but also depends on the difference
between the coefficient functions and anomalous dimensions. To clarify the latter we choose the simplest form of δq:
δq(ω) = Nq . (49)
It corresponds to the evolution from the bare quark where δq(x) = Nqδ(1 − µ2/s). Numerical results for R =
[gNS1 − gNS DGLAP1 ]/gNS DGLAP1 with δq chosen by Eq. (49) manifest (see Ref. [19] for detail) that R increases when
x is decreases. In particular, R > 0.3 at x . 0.05. It means that the total resummation of leading ln
(1/x) cannot
be neglected at x . 0.05 and DGLAP cannot be used beyond x ≈ 0.05. On the other hand, it is well–known that
Standard Approach based on DGLAP works well at x ≪ 0.05. To solve this puzzle, we have to consider the standard
fit for δq in more detail.
C. Analysis of the standard fits for δq
There are known different fits for δq. We consider the fit of Eq. (32). Obviously, in the ω -space Eq. (32) is a sum
of pole contributions:
δq(ω) = Nη
(ω − α)−1 +
mk(ω + λk)
, (50)
with λk > 0, so that the first term in Eq. (50) corresponds to the singular term x
−α of Eq. (32) and therefore the
small-x asymptotics of fDGLAP is given by the leading singularity ω = α = 0.57 of the integrand in Eq. (50) so that
the asymptotics of gNS DGLAP1 (x,Q
2) is not given by the classic exponential of Eq. (47) but actually is the Regge-like:
gNS DGLAP1 ∼ C(α)(1/x)α
ln(Q2/Λ2)/ ln(µ2/Λ2)
)γ(α)/b
, (51)
with b = (33 − 2nf)/12π. Comparison of Eq. (48) and Eq. (51) demonstrates that both DGLAP and our approach
lead to the Regge behavior of g1, though the DGLAP prediction is more singular than ours. Then, they predict
1 The singlet intercept is much greater: ∆S = 0.86.
different Q2 -behavior. However, it is important that our intercept ∆NS is obtained by the total resummation of the
leading logarithmic contributions and without assuming singular fits for δq whereas the SA intercept α in Eq. (47)
is generated by the phenomenological factor x−0.57 of Eq. (32) which makes the structure functions grow when x
decreases and mimics in fact the total resummation2. In other words, the role of the higher-loop radiative corrections
on the small-x behavior of the non-singlets is, actually, incorporated into SA phenomenologically, through the initial
parton densities fits. It means that the singular factors can be dropped from such fits when the coefficient functions
account for the total resummation of the leading logarithms and therefore fits for δq become regular in x in this
case. They also can be simplified. Indeed, if x in the regular part N
(1 − x)β(1 + γxδ)
of the fit (32) is not large,
all x -dependent terms can be neglected. So, instead of the rather complicated expression of Eq. (32), δq can be
approximated by a constant or a linear form
δq(x) = N(1 + ax) . (52)
with 2 phenomenological parameters instead of 5 in Eq. (32).
V. CORRECTING MISCONCEPTIONS
The total resummation of lnk(1/x) allows to correct several misconceptions popular in the literature. We list and
correct them below.
Misconception 1: Impact of non-leading perturbative and non-perturbative contributions on the intercepts of g1
is large.
Actually: Confronting our results and the estimates of the intercepts in Refs. [18] obtained from fitting
available experimental data manifests that the total contribution of non-leading perturbaive and non-perturbative
contributions to the intercepts is very small, so the main impact on the intercepts is brought by the leading logarithms.
Misconception 2: Intercepts of g1 should depend on Q
2 through the parametrization of the QCD coupling
αs = α(Q
Actually: This is groundless from the theoretical point of view and appears only if the the parametrization of the
QCD coupling αs = α(k
) is kept in all ladder rungs. It is shown in Ref. [16] that this parametrization cannot be
used at small x and should be replaced by the parametrization of Eq. (40).
Misconception 3: Initial densities δq(x) and δg(x) are singular but they are defined at x not too small. Later,
being convoluted with the coefficient functions, they become less singular.
Actually: It is absolutely wrong: Eq. (50) proves that the pole singularity x−α in the fits does not become weaker
with the x-evolution.
Misconception 4: Fits for the initial parton densities are complicated because they mimic unknown non-
perturbative contributions.
Actually: Our results demonstrate that the singular factors in the fits mimic the total resummation of lnk(1/x)
and can be dropped when the resummation is accounted for. In the regular part of the fits the x -dependence is
essential for large x only, so impact of non-perturbative contributions is weak at the small-x region.
Misconception 5: Total resummations of lnk(1/x) may become of some importance at extremely small x but not
for x available presently and in a forthcoming future.
Actually: The efficiency of SA in the available small-x range is based on exploiting the singular factors in the
standard fits to mimic the resummations. So, the resummations have always been used in SA at small x in an
inexplicit way, through the fits, but without being aware of it.
2 We remind that our estimates for the intercepts ∆NS ,∆S were confirmed (see Refs. [18]) by analysis of the experimental data
VI. COMBINING THE TOTAL RESUMMATION AND DGLAP
The total resummaton of leading logarithms of x considered in Sect. IV is essential at small-x. When x ∼ 1, all
terms ∼ lnk(1/x) in the coefficient functions and anomalous dimensions cannot have a big impact compared to other
terms. DGLAP accounts for those terms. It makes DGLAP be more precise at large x than our approach. So,
there appears an obvious appeal to combine the DGLAP coefficient functions and anomalous dimensions with our
expressions in order to obtain an approach equally good in the whole range of x : 0 < x < 1. The prescription
for such combining was suggested in Ref. [19]. Let us, for the sake of simplicity, consider here combining the total
resummation and LO DGLAP. The generalization to NLO DGLAP can be done quite similarly. The prescription
consists of the following points:
Step A: Take Eqs. (31) and replace αs by A of Eq. (40), converting γNS into γ̃NS and C
NS into C̃
Step B: Sum up the obtained expressions and Eqs. (42,43):
c̃NS = C̃
NS +HS , h̃NS = γ̃NS +HNS . (53)
New expressions c̃NS , h̃NS combine the total resummation and DGLAP but they obviously contain the double count-
ing: some of the first–loop contributions are present both in Eqs. (31) and in Eqs. (42,43). To avoid the double
counting, let us expend Eqs. (42,43) into series and retain in the series only the first loop contributions3:
A(ωCF )
NS = 1 +
A(ωCF )
. (54)
Finally, there is Step C: Subtract the first-loop expressions (54) from Eq. (53)) to get the combined, or ”synthetic”
as we called them in Ref. [19], coefficient function cNS and anomalous dimension hNS :
cNS = c̃NS − C(1)NS , hNS = h̃NS −H
NS . (55)
Substituting Eqs. (55) in Eq. (41) leads to the expression for gNS1 equally good at large and small x. This description
does not require singular factors in the fits for the initial parton densities. An alternative approach for combining
DLA expression for g1 was suggested in Ref. [20]. However, the parametrization of αs in this approach was simply
borrowed from DGLAP, which makes this approach be unreliable at small x.
VII. CONCLUSION
We have briefly considered the essence of the IREE method together with examples of its application to different
processes. They demonstrate that IREE is indeed the efficient and reliable instrument for all-orders calculations in
QED, QCD and the Standard Model of EW interactions. As an example in favor of this point, let us just remind that
there exist wrong expressions for the singlet g1 in DLA obtained with an alternative technique and the exponentiation
of EW double logarithms obtained in Ref. [10] had previously been denied in several papers where other methods of
all-order summations were used.
VIII. ACKNOWLEDGEMENT
B.I. Ermolaev is grateful to the Organizing Committee of the Epiphany Conference for financial support of his
participation in the conference.
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[3] V.N. Gribov. Yad. Fiz. 5(1967)399.
3 For combining the total resummation with NLO DGLAP one more term in the series should be retained
[4] B.I. Ermolaev, L.N. Lipatov, V.S. Fadin. Yad. Fiz. 45(1987)817; B.I. Ermolaev. Yad. Fiz. 49(1989)546; M. Chaichian and
B. Ermolav. Nucl. Phys. B 451(1995)194.
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L.N.Lipatov, Sov. J. Nucl. Phys. 20 (1972) 95; Yu.L. Dokshitzer, Sov. Phys. JETP 46 (1977) 641.
[12] G. Altarelli, R.D. Ball, S. Forte and G. Ridolfi. Nucl. Phys. B496 (1997) 337; Acta Phys. Polon. B29(1998)1145; E. Leader,
A.V. Sidorov and D.B. Stamenov. Phys. Rev. D73 (2006) 034023; J. Blumlein, H. Botcher. Nucl. Phys. B636 (2002) 225;
M. Hirai at al. Phys. Rev. D69 (2004) 054021.
[13] W.L. Van Neerven. hep-ph/9609243.
[14] B.I. Ermolaev, S.I. Manaenkov and M.G. Ryskin. Z. Pyss. C 69(1996)259; J. Bartels, B.I. Ermolaev and M.G. Ryskin.
Z. Pyss. C 70(1996)273; Z. Pyss. C 72(1996)627.
[15] B.I. Ermolaev, M. Greco, S.I. Troyan. Nucl. Phys.B 571 (2000) 137; Nucl. Phys.B 594 (2001) 71; Phys.Lett.B 579 (2004)
[16] B.I. Ermolaev, M. Greco and S.I. Troyan. Phys.Lett.B 522(2001)57.
[17] B.I. Ermolaev, M. Greco and S.I. Troyan. hep-ph/0605133.
[18] J. Soffer and O.V. Teryaev. Phys. Rev.56( 1997)1549; A.L. Kataev, G. Parente, A.V. Sidorov. Phys.Part.Nucl 34(2003)20;
Nucl.Phys.A666(2000)184; A.V. Kotikov, A.V. Lipatov, G. Parente, N.P. Zotov. Eur.Phys.J.C26(2002)51; V.G. Krivohijine,
A.V. Kotikov, hep-ph/0108224; A.V. Kotikov, D.V. Peshekhonov hep-ph/0110229.
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http://arxiv.org/abs/hep-ph/9609243
http://arxiv.org/abs/hep-ph/0605133
http://arxiv.org/abs/hep-ph/0108224
http://arxiv.org/abs/hep-ph/0110229
http://arxiv.org/abs/hep-ph/9802386
Introduction
IREE for scattering amplitudes in the hard kinematics
IREE for the form factor f(q2) in QED
IREE for the form factor g(q2) in QED
e+e- -annihilation into a quark-antiquark pair
e+e- -annihilation into a quark-antiquark pair and gluons
Exponentiation of Sudakov electroweak double-logarithmic contributions
Application of IREE to the polarized Deep-Inelastic Scattering
Comparison of expressions (30) and (41) for g1NS
Comparison of small-x asymptotics, neglecting the impact of q
Numerical comparison between Eqs. (30) and (41), neglecting the impact of q
Analysis of the standard fits for q
Correcting misconceptions
Combining the total resummation and DGLAP
Conclusion
Acknowledgement
References
|
0704.0342 | Cofibrations in the Category of Frolicher Spaces. Part I | Cofibrations in the Category of Frölicher Spaces:
Part I
Brett Dugmore
Cadiz Financial Strategists (Pty) Ltd, Cape Town, South Africa
Email: [email protected]
Patrice Pungu Ntumba
Department of Mathematics and Applied Mathematics
University of Pretoria
Hatfield 0002, Republic of South Africa
Email: [email protected]
Abstract
Cofibrations are defined in the category of Frölicher spaces by weak-
ening the analog of the classical definition to enable smooth homotopy
extensions to be more easily constructed, using flattened unit intervals.
We later relate smooth cofibrations to smooth neighborhood deforma-
tion retracts. The notion of smooth neighborhood deformation retract
gives rise to an analogous result that a closed Frölicher subspace A of the
Frölicher space X is a smooth neighborhood deformation retract of X if
and only if the inclusion i : A →֒ X comes from a certain subclass of
cofibrations. As an application we construct the right Puppe sequence.
Subject Classification (2000): 55P05.
Key Words:Frölicher spaces, Flattened unit intervals, Smooth neighborhood de-
formation retracts, Smooth cofibrations, Cofibrations with FCIP, Puppe se-
quence.
1 Preliminaries
The purpose of this section is to survey brielfy the notion of Frölicher spaces.
Frölicher spaces arise naturally in physics, and do generalize the concept of
smooth manifolds. A Frölicher space, or smooth space as initially called by
Frölicher and Kriegl [7], is a triple (X, CX ,FX) consisting of a setX , and subsets
CX ⊆ XR, FX ⊆ RX such that
• FX ◦ CX = {f ◦ c| f ∈ FX , c ∈ CX} ⊆ C∞(R)
• ΦCX := {f : X → R| f ◦ c ∈ C∞(R) for all c ∈ CX} = FX
http://arxiv.org/abs/0704.0342v1
• ΓFX := {c : R → X | f ◦ c ∈ C∞(R)for all f ∈ FX} = CX
Frölicher and Kriegl [7], and Kriegl and Michor [10] are our main reference for
Frölicher spaces. The following terminology will be used in the paper: Given a
Frölicher space (X, CX ,FX), the pair (CX ,FX) is called a smooth structure; the
elements of CX and FX are called smooth curves and smooth functions respec-
tively. The topology assumed for a Frölicher space (X, CX ,FX) throughout the
paper is the initial topology TF induced by the set FX of functions. When there
is no fear of confusion, a Frölicher space (X, CX ,FX) will simply be denoted X .
The most natural Frölicher spaces are the finite dimensional smooth manifolds,
where if X is such a smooth manifold, then CX and FX consist of all smooth
curves R → X and smooth functions X → R. Euclidean finite dimensional
smooth manifolds Rn, when viewed as Frölicher spaces, are called Euclidean
Frölicher spaces. In the sequel, by Rn, n ∈ N, we mean the Frölicher space Rn,
equipped with its usual smooth manifold structure.
A Frölicher space X is called Hausdorff if and only if the smooth real-valued
functions on X are point-separating, i.e. if and only if TF is Hausdorff.
A Frölicher structure (CX ,FX) on a set X is said to be generated by a set
F0 ⊆ RX (resp. C0 ⊆ XR) if CX = ΓF0 and FX = ΦΓF0 (resp. FX = ΦC0
and CX = ΓΦC0 ). Note that different sets F0 ⊆ RX on the same set X may
give rise to a same smooth structure on X . A set mapping ϕ : X → Y between
Frölicher spaces is called a map of Frölicher spaces or just a smooth map if for
each f ∈ FY , the pull back f ◦ϕ ∈ FX . This is equivalent to saying that for each
c ∈ CX , ϕ ◦ c ∈ CY . For Frölicher spaces X and Y , C∞(X,Y ) will denote the
collection of all the smooth maps X → Y . The resulting category of Frölicher
spaces and smooth maps is denoted by FRL.
Some useful facts regarding Frölicher spaces can be gathered in the following
Theorem 1.1 The category FRL is complete (i.e. arbitrary limits exist ), co-
complete (i.e. arbitrary colimits exist), and Cartesian closed.
Given a collection of Frölicher spaces {Xi}i∈I , let X =
i∈I Xi be the set
product of the sets {Xi}i∈I and πi : X → Xi, i ∈ I, denote the projection map
(xi)i∈I 7→ xi. The initial structure on X is generated by the set
{f ◦ πi : f ∈ FXi}.
The ensuing Frölicher space (X,ΓF0, ϕΓF0) is called the product space of the
family {Xi}i∈I . Clearly,
ΓF0 = {c : R → X | if c(t) = (ci(t))i∈I , then ci ∈ CXi for every i ∈ I}.
Now, let
i∈I Xi be the disjoint union of sets {Xi}i∈I , and ιXi : Xi →
i∈I Xi
the inclusion map. Place the smooth final structure on
i∈I Xi corresponding
to the family {ιXi}i∈I . The resulting Frölicher space is called the coproduct of
{Xi}i∈I , and denoted
i∈I Xi, and
Xi = {f :
Xi → R| for each i ∈ I, f |Xi ∈ FXi}
is the collection of smooth functions for the coproduct.
Corollary 1.1 Let X, Y , and Z be Frölicher spaces. Then the following canon-
ical mappings are smooth.
• ev: C∞(X,Y )×X → Y , (f, x) 7→ f(x)
• ins:X → C∞(Y,X × Y ), x 7→ (y 7→ ins(x)(y) = (x, y))
• comp:C∞(Y, Z)× C∞(X,Y ) → C∞(X,Z), (g, f) 7→ g ◦ f
• f∗ : C∞(X,Y ) → C∞(X,Z), f∗(g) = f ◦ g, where f ∈ C∞(Y, Z)
• g∗ : C∞(Z, Y ) → C∞(X,Y ), g∗(f) = f ◦ g, where g ∈ C∞(X,Z).
Given Frölicher spaces X , Y , and Z; in view of the cartesian closedness of
the category FRL, the exponential law
C∞(X × Y, Z) ∼= C∞(X,C∞(Y, Z))
holds. Because FX = C∞(X,R), it follows by cartesian closedness of FRL that
the collection FX can be made into a Frölicher space on its own right.
Finally we would like to show how to construct smooth braking functions,
following Hirsch [8]. Smooth braking functions are tools that are behind most
results in this paper. In [11], it is shown that the function ϕ : R → R given by
ϕ(u) =
0 if u ≤ 0
u if u > 0
is smooth. Substituting x2 for u in the above function, one sees that the function
ψ : R → R, given by
ψ(x) =
0 if x ≤ 0
x2 if u > 0
is smooth. Now, let us construct a smooth function α : R → R with the following
properties. Let 0 ≤ a < b. α(t) should satisfy:
• α(t) = 0 for t ≤ a,
• 0 < α(t) < 1 for a < t < b,
• α is strictly increasing for a < t < b,
• α(t) = 1 for t ≥ b.
Define α : R → [0, 1] by
α(t) =
γ(x)dx
γ(x)dx
where γ(x) = ψ(x− a)ψ(b − x).
In the sequel, the notation αǫ, 0 < ǫ <
, will refer to a smooth braking
function with the following properties
• αǫ(t) = 0 for t ≤ ǫ,
• 0 < αǫ(t) < 1 for ǫ < t < 1− ǫ,
• α strictly increasing for ǫ < t < 1− ǫ,
• αǫ(t) = 1 for 1− ǫ ≤ t.
2 Basic Constructions of Homotopy Theory in
In this section, we define the fundamental notions of homotopy theory in the
category FRL, such as the homotopy relation and the mapping cylinder. We
begin with an overview of our approach to homotopy in FRL, and then discuss
alternate Frölicher structures on the unit interval which are used in this and
subsequent sections.
2.1 Our Approach to Homotopy Theory in FRL
One might begin investigating homotopy theory in FRL by simply following
the homotopy theory of topological spaces, replacing continuous functions with
smooth ones. One can certainly define the notion of a homotopy H : I×X → Y
between smooth maps H(0,−) and H(1,−) in this way (which we do). One can
even get as far as the left Puppe sequence (see [4]), but eventually difficulties
begin to arise.
Extending functions defined on a subspace of a Frölicher space tends to be
a little tricky, and so the definition of a cofibration in FRL is one that needs
careful consideration. We envisage to construct the right Puppe sequence in a
future paper. To do this we define a slightly weaker notion of cofibration than
the notion obtained from topological spaces. In addition, we define the mapping
cylinder of a smooth map f : X → Y using not the unit interval, but a modified
version called the weakly flattened unit interval, denoted I, which, as one
can show, is topologically homeomorphic to the unit interval. This modified
structure on the unit interval allows us to show that the inclusion of a space X
into the mapping cylinder of f : X → Y is a cofibration (in our weaker sense ).
The weakly flattened unit interval is useful, but it also has its drawbacks.
It would be ideal to have a single structure on the unit interval that can be
used throughout out homotopy theory, but the weakly flattened unit interval
is not suitable, because it has the rather restrictive property that a smooth
map f : I → I on the usual unit interval often does not define a smooth map
f : I → I unless the endpoints of the interval are mapped to the endpoints. This
restrictive property means that we only use the flattened unit intervals where
they are absolutely necessary.
In our future work, we will investigate whether with our modified notions of
cofibration and mapping cylinder, Baues’ cofibration axioms are satisfied.
2.2 Flattened Structures on the Unit Interval
We define two main Frölicher structures which we call the flattened unit in-
terval and the weakly flattened unit interval . Let (CI ,FI) be the subspace
structure induced on I by the inclusion I →֒ R.
Definition 2.1 The Frölicher space (I, CI,FI), where the structure (CI,FI) is
the structure generated by the set
F = {f ∈ FI| there exists 0 < ǫ < 14 with f(t) = f(0) for t ∈ [0, ǫ) and
f(t) = f(1) for t ∈ (1− ǫ, 1]},
is called the flattened unit interval.
It is easy to see that any continuous map c : R → [0, 1] defines a structure
curve on I if and only if it is smooth at every point t ∈ R, where c(t) ∈ (0, 1), .
We define the left (resp. right) flattened unit interval, denoted by I−
(resp. I+), to be the Frölicher space whose underlying set is the unit interval
[0, 1], and structure is the structure generated by the structure functions in FI
that are constant near 0 (resp. 1).
Definition 2.2 The Frölicher space (I, CI,FI), with the structure defined below
is called the weakly flattened unit interval. The underlying set is the unit
interval; the structure (CI,FI) is generated by the family
F = {f ∈ FI | lim
f(t) = 0, lim
f(t) = 0, n ≥ 1}.
We call the property, for all f ∈ F ,
f(t) = 0, lim
f(t) = 0, n ≥ 1,
the zero derivative property of f .
We shall prove that all structure functions on I have the zero derivative
property, in other words, FI = F . To that effect, we need the following lemma.
Lemma 2.1 Let c : R → R be a smooth real-valued function at t = t0, and let
f : R → R be a smooth real-valued function at t = c(t0). Then,
(f ◦ c)(t0) = f (n)(c(t0))(c′(t0))n + terms of the form
af (k)(c(t0))(c
′(t0))
m1(c′′(t0))
m2 . . . (c(n−1)(t0))
mn−1 ,
where k < n and a ∈ R. In addition, if a 6= 0 then at least one ofm2,m3, . . . ,mn−1
is also non-zero.
Proof. The proof is done by induction. For the sake of brevity, we call the
term f (n)(c(t0))(c
′(t0))
n the primary term for n, and the terms of the form
af (k)(c(t0))(c
′(t0))
m1(c′′(t0))
m2 . . . (c(n−1)(t0))
mn−1 the lower order terms for n.
The statement is true for n = 1 and for n = 2. Suppose the result is true for
n = k. To show that the result holds for n = k + 1, since
dtk+1
(f ◦ c)(t0) =
(f (k)(c(t0))(c
′(t0))
+terms of the form d
(af (j)(c(t0))(c
′(t0))
m1(c′′(t0))
m2 . . . (c(k−1)(t0))
mk−1),
where j < k + 1 and a ∈ R, we need only show that
(af (j)(c(t0))(c
′(t0))
m1(c′′(t0))
m2 . . . (c(k−1)(t0))
mk−1)
gives rise to lower terms for n = k + 1, which is by the way straightforward. �
Theorem 2.1 FI = {f ∈ FI | limt→0+ d
f(t) = 0 = limt→1−
f(t)} =: F
Proof. That F ⊆ FI is evident. We must show the reverse inequality. Let
0 < ǫ < 1
, and 0 < M < 1. Consider the function cM : R → R, given by
cM (t) = (1− αǫ(|t|))βM (t) + αǫ(|t|),
where αǫ : R → R is a smooth braking function as defined in the Preliminaries,
and βM : R → R is given by
βM (t) =
−Mt if t ≤ 0
t if t > 0
It is easily seen that cM is continuous over all R, and smooth over all R except
at t = 0. Also note that 0 < cM (t) < 1 for all t ∈ R, and cM (t) = βM (t) = 0 for
all 0 ≤ t < ǫ. Now,
cM (t) =
βM (t) = −M, for −ǫ < t < 0
cM (t) =
βM (t) = 1, for 0 < t < ǫ
For n > 1, we have
cM (t) =
βM (t) = 0, for t ∈ (−ǫ, 0) ∪ (0, ǫ).
We now show that for cM ∈ ΓF . To this end, let f ∈ F . To show that
f ◦ cM : R → R is smooth, it is obvious that we need only concentrate on the
point t = 0, because f ◦ c is smooth at every t 6= 0. It follows for t 6= 0, and
n ∈ N that Lemma 2.1 applies. But as t → 0, cM (t) → 0+, and so, letting
s = cM (t), we have
f (j)(cM (t)) = lim
f (j)(s) = 0,
for all j ∈ N, by the zero derivative property of f . Thus, as t approaches
the value 0, the primary term and all the lower order terms of d
(f ◦ cM )(t)
vanish, and we have shown that f ◦ cM is smooth at t = 0. This implies that
f ◦ cM ∈ C∞(R,R) for all f ∈ F . It follows that cM ∈ ΓF .
We are now ready to show that FI ⊆ F . To this end, suppose that we
are given a structure function f ∈ FI. We shall show that this f has the zero
derivative property, and is thus an element of F .
Since f ∈ FI, we know that f ◦ c is a smooth real-valued function for every
c ∈ ΓF . In particular, f ◦ cM is smooth for all 0 < M < 1. Thus, for any n ∈ N,
(f ◦ cM )(t) = lim
(f ◦ cM )(t).
As t→ 0−, cM (t) → 0+; let us consider the lower order terms for n. Each term
of the form
af (k)(cM (t))(c
M (t))
m1(c′′M (t))
m2 . . . (c
(n−1)
M (t))
has some term (c
(t))mi , for some i > 1, with mi 6= 0. But limt→0− c
(t) = 0,
if i > 1, and so
af (k)(cM (t))(c
M (t))
m1(c′′M (t))
m2 . . . (c
(n−1)
(t))mn−1 = 0.
So all the lower order terms fall away, therefore
limt→0−
(f ◦ cM )(t) = limt→0− f (n)(cM (t))(c′M (t))n
= limt→0− f
(n)(cM (t))(−M)n
= lims→0+ f
(n)(s)(−M)n,
where s = cM (t). In a similar way one shows that
(f ◦ cM )(t) = lim
f (n)(s).
But f◦cM is smooth, therefore lims→0+ f (n)(s)(−M)n = lims→0+ f (n)(s), which
implies that lims→0+ f
(n)(s) = 0.
We have shown that the zero derivative property of f holds for the left
endpoint of the unit interval. To show that the zero derivative property of f
holds for the right endpoint of f , note that dM : R → R, dM (t) = 1− cM (t), is
a smooth real-valued function with d(0) = 1, and 0 ≤ dM (t) ≤ 1 for all t ∈ R.
One can follow a similar procedure to the above, using dM instead of cM to
show that lims→1− f
(n) = 0. �
2.3 Some Properties of Smooth Functions between the
Flattened Unit Intervals
One has to be careful when dealing with the various flattened unit intervals. A
smooth function f : I → I from the R- Frölicher subspace unit interval I to
itself need not define a smooth function f : I → I, for example. Conversely,
not every smooth function f : I → I defines a smooth function f : I → I. In
particular, we need to be aware of the fact that addition and multiplication of
functions when defined between the various flattened unit intervals does not
preserve smoothness, as is the case with the usual unit interval.
Example 2.1
The function f : I → I, f(t) = 1
t is clearly smooth, but the corresponding
function f : I → I, given by the same formula, is not smooth. To see this, let
α : R → R be a smooth braking function with the properties that
• α(t) = −1, for t < − 3
• α(t) = t, for − 1
< t < 1
• α(t) = 1, for t > 3
Define c : R → I by c(t) = 1− |α(t)|. The curve c is smooth everywhere except
at t = 0, where c(0) = 1. However, every generating function f on I is constant
near 1, and so the composite f ◦ c is smooth. Thus c is a structure curve on I.
Now, f ◦ c : R → I is given by (f ◦ c)(t) = 1
(1 − |α(t)|). Let h : I → R be a
structure function with the properties that
• h(s) = 0, for s < 1
• h(s) = s, for 1
< s < 3
• h(s) = 1, for 7
Then (h ◦ f ◦ c)(t) = 1
(1− |α(t)|) for t near 0, and is not smooth at t = 0. Thus
f does not define a smooth function from I to I.
Example 2.2
The function f : I → I, f(t) =
t, is smooth, but the corresponding f : I → I,
given by the same formula, is not smooth. This follows from the fact that f is
smooth on the open interval (0, 1), and a generating function g on I is constant
near 0 and 1. On the side, f : I → I is not smooth, because if c : R → I is a
structure curve with c(t) = t2 near t = 0, then (f ◦ c)(t) = |t| near t = 0, which
is not smooth on I at t = 0.
Example 2.3
The functions f, g : I− → I−, given by f(t) = 1
t and g(t) = 1
are both
smooth, but the sum f(t) + g(t) = 1
is not smooth.
The following lemma follows from the definition of the Frölicher structures
on the various flattened unit intervals.
Lemma 2.2 Let f : I → I be a smooth function with the properties that f(0) =
0 and f(1) = 1. Then the following maps are smooth:
• f : I → I±,
• f : I → I,
• f : I± → I,
• f : I → I,
• f : I → I.
The function defined in the following example is for later reference.
Example 2.4
Let H : I × I− → I− be given by H(t, s) = (1 − α(t))s, where α : R → R is a
smooth braking function with the properties that
• α(t) = 0 for t < 1
• 0 ≤ α(t) ≤ 1 for all t ∈ R,
• α(t) = 1 for t > 3
We show that H is smooth. To see this, let f : I− → R be a generating function
on I−. So f is constant near 0. Now, let c : R → I × I− be a structure curve,
given by c(v) = (t(v), s(v)). The curve t is a structure curve on I, and so is
a smooth real-valued function for all v ∈ R, except possibly when t(v) = 0 or
t(v) = 1. Similarly, the curve s is a structure curve on I−, and so is smooth
for all v ∈ R except possibly when s(v) = 0. Now consider the composite
H ◦ c : R → I−. Clearly, α(t(v)) is smooth for all v, since the only possible
points for non-smoothness occur when t(v) = 0 or t(v) = 1, and α(t(v)) is
locally constant near these points. Consequently, H ◦ c is smooth everywhere
except possibly when s(v) = 0. Now, let’s consider f ◦H ◦ c : R → R; the only
possible points for non-smoothness are those in which s is 0, i.e. H◦ = 0. But f
is a structure generating function on I−, and so is locally constant near 0. This
shows that f ◦H ◦ c is smooth for all v ∈ R, and thus H is smooth.
2.4 Homotopy in FRL and Related Objects
Definition 2.3 (1) Let X be a Frölicher space, and x0, x1 ∈ X. We say that
x0 is smoothly path-connected to x1 if there is a smooth path c : I → X such
that c(0) = x0 and c(1) = x1. We write x0 ≃ x1. The relation ≃ is called
smooth homotopy when it is applied to hom-sets.
(2) Let f : X → Y be a map of Frölicher spaces. f is called a smooth
homotopy equivalence provided there exists a smooth map g : Y → X such that
f ◦ g ≃ 1Y and g ◦ f ≃ 1X .
One can show that smooth homotopy is a congruence in RFL. In practice, we
say that smooth maps f, g : X → Y are smoothly homotopic if there exists a
smooth map H : I ×X → Y with H(0,−) = f and H(1,−) = g. If A ⊆ X is
subspace of X , then we say that H is a smooth homotopy (rel A) if the map
H has the additional property that H(t, a) = a for each t ∈ I and a ∈ A. See
Cherenack [5] and Dugmore [6] for more detail regarding smooth homotopy.
The notion of deformation retract is fundamental to topological homotopy
theory. The following definitions are adapted for smooth homotopy, and will be
needed at a later stage.
Definition 2.4 Let A ⊆ X be a subspace of a Frölicher space X, and let i :
A →֒ X denote the inclusion map. Then
• We say that A is a retract of X if there exists a smooth map r : X → A
such that ri = 1A. We call r a retraction.
• We call A a weak deformation retract of X if the inclusion i is a smooth
homotopy equivalence.
• The subspace A is called a deformation retract of X if there exists a re-
traction r : X → A such that ir ≃ 1X .
• The subspace A is called a strong deformation retract of X if there exists
a retraction r : X → A such that ir ≃ 1X(relA).
Definition 2.5 The mapping cylinder If of f : X → Y is defined by the fol-
lowing pushout
I ×X // If
where i1 : X → I ×X is given by i1(x) = (1, x), for any x ∈ X. We denote the
elements of If by [t, x] or [y], where (t, x) ∈ I ×X and y ∈ Y .
Replacing I ×X in the above pushout diagram by I×X or I×X, we obtain
the flattened mapping cylinder If and weakly flattened mapping cylinder If of f
respectively. We use the same notation for elements of these flattened mapping
cylinders as described above for the mapping cylinder.
There is also a map i0 : X → I ×X , defined by i0(x) = (0, x) for x ∈ X . This
induces an inclusion map i′0 : X → If , which identifies X with the Frölicher
subspace i′0(X) of If . An inclusion is induced in a similar way for the flattened
mapping cylinders. If one identifies {0}×X to a point in the mapping cylinder
If of a map f : X → Y , then one obtains the mapping cone Tf of the
map f . In a similar fashion, we define the flattened mapping cone Tf and
weakly flattened mapping cone Tf of a smooth map f : X → Y .
2.5 Cofibrations in FRL
A cofibration is a map i : A→ X for which the problem of extending functions
from i(A) to X is a homotopy problem. In other words, if a map f : i(A) → Z
can be extended to a map f∗ : X → Z, then so can any map homotopic to f . For
topological spaces, the usual definition is phrased in a slightly more restrictive
way. The extension of a map g ≃H f , for some homotopy H : I × i(A) → Z, is
required to exist at every level of the homotopy simultaneously. In other words,
one requires each H(t,−) to be extendable in such a way that the resulting
homotopy H∗ : I ×X → Z is continuous.
We weaken this definition somewhat, to enable smooth homotopy extensions
to be more easily constructed using a flattening at the endpoints of the homo-
topy. This enables us to characterize smooth cofibrations in terms of a flattened
unit interval, and then later to relate smooth cofibrations to smooth neigh-
borhood deformation retracts. Our definition of smooth cofibration, though
different from from Cap’s definition, see [1], leads to several classical results as
does Cap’s. As pointed out by Cap, the analogue of the classical definition of
cofibration would not allow even {0} →֒ I to be a smooth cofibration. So, we
have the following
Definition 2.6 A smooth map i : A → X is called a smooth cofibration if,
corresponding to to every commutative diagram of the form
(0,1A)
f // Z
66mmmmmmmmmmmmmm
there exists a commutative diagram in FRL of the form
(0,1X )
::tttttttttt
where G′ : I × A → Z is given by G′(t, a) = G(αǫ(t), a) for some 0 < ǫ < 12 ,
and each t ∈ I, a ∈ A.
The problem of extending a map smoothly from a subspace of a Frölicher
space to the whole space is a more difficult problem than simply extending con-
tinuously. It is mainly for this reason that the definition of smooth cofibration
differs somewhat from the corresponding definition of a topological cofibration.
Lemma 2.3 Let i : A → X be a smooth cofibration, then i is an initial mor-
phism in FRL. In addition, if A is Hausdorff, then i is injective. So in this
case A can be regarded as a subspace of X.
Proof. Let us show that every smooth map f : A→ R factors through i, that
is for every f ∈ FA, there exists f̃ ∈ FX such that f = f̃ ◦ i. To this end,
consider the smooth map G : I × A → R, given by H(t, a) = tf(a). Clearly,
0|A = G(0,−), where 0 : X → R is the constant map 0. It follows that there is
map F : I ×X → R such that F ◦ (1× i) = G′. Then, clearly f̃ := F (1,−) has
the desired property.
The remaining part of the proof of Proposition 3.3, in [1], holds verbatim
here as well. �
In this paper, we are interested only in cofibrations that are injective. Hence-
forth, all cofibrations are assumed to be injective.
All topological cofibrations are inclusions, and this result is true for smooth
cofibrations too. The proof of the following lemma is essentially the same as
the proof given by James [9] for the topological result, although James’s proof
is in some sense dual to ours, using path-spaces in place of cartesian products
and the adjoint versions of our homotopies.
Lemma 2.4 A cofibration
i // X
is a smooth inclusion.
Proof. Let Ii be a mapping cylinder of i, and let j : X → Ii be the standard
inclusion map. Consider the smooth map γ : I → I, γ(t) = 1 − t, for all t ∈ I,
and the quotient map q : (I ×A) ⊔X → Ii; we have the following commutative
diagram
(0,1A)
j // Ii
66mmmmmmmmmmmmmm
where G(t, a) = [(1 − t, a)]. Notice that the map G is smooth. Since i is a
cofibration, we have the commutative diagram
(0,1X )
::uuuuuuuuuu
where G′(t, a) = G(αǫ(t), a) for some 0 < ǫ <
. Define U : X → Ii by
U(x) = F (1, x). We have U ◦ i = G′(1,−), where G′(1, a) = [(0, a)], for every
a ∈ A. Thus the assignment a 7→ G′(1, a) defines the usual inclusion of A into
the mapping cylinder. From this we deduce that U ◦ i is an inclusion, and hence
i is an inclusion. �
There is an equivalent formulation of definition 2.6, given in the following
lemma.
Lemma 2.5 A smooth map
i // X
is a cofibration if and only if, for every smooth map h : (0×X)∪(I−×i(A)) → Z,
the following diagram
(0×X) ∪ (I− × i(A)) h //
I− ×X
77oooooooooooooo
where j is the evident inclusion, exists in FRL.
Proof. Suppose that the inclusion A //
i // X is a smooth cofibration, and
suppose that h : (0 × X) ∪ (I− × i(A)) → Z is a smooth map. We have the
diagram
(0×B) ∪ (I− × i(A)) h //
I− ×X
We need to fill in a smooth map G : I− × X → Z which makes the resulting
diagram commute. To do this, notice that h|I−×i(A) is smooth, and thus the
corresponding map h|I × i(A), using the usual unit interval, is also smooth. We
have the following diagram
(0,1A)
h|0×X // Z
66mmmmmmmmmmmmmm
where h|0×X(0,−) : X → Z is denoted as h|0×X . The fact that i is a smooth
cofibration yields the following FRL-commutative diagram:
h0×X //
(0,1A)
::tttttttttt
where (h|I−×A)′(t, a) = h|I−×A(αǫ(t), a), for some 0 < ǫ < 12 . Now, chose a
smooth braking function β : R → R with the following properties.
• α(t) = 0 for t < ǫ
• α(t) = t for ǫ < t.
F may not be smooth on I− × A due to the flattening requirements of the left
flattened unit interval. To correct this, set G(t, a) = F (β(t), a). Notice that the
insertion of this braking function does not affect the commutativity conditions
of G, since the only adjustments to F occur in the first coordinate where the
map (h|I−×X)′ is constant.
Now, assume the converse, i.e. to every smooth map h : (0 × X) ∪ (I− ×
i(A)) → Z, corresponds a commutative diagram
(0×X) ∪ (I− × i(A)) h //
I− ×X
77oooooooooooooo
We wish to show that the inclusion i : A → X is a cofibration; so assume we
have the following diagram
(0,1A)
f // Z
66mmmmmmmmmmmmmm
There exists the diagram
(0,1A)
f // Z
66mmmmmmmmmmmmmm
where G′(t, a) = G(αǫ(t), a). Our hypothesis allows us to construct the diagram
(0×X) ∪ (I− × i(A))
f∪G′ //
I− ×X
77oooooooooooooo
Note that f ∪ G′ is smooth since αǫ(t) is constant near 0. Since H is smooth
on I− ×X it defines a smooth map on I ×X . One can verify that the diagram
(0,1X)
::tttttttttt
commutes as required. �
3 Smooth Neighborhood Deformation Retracts
This section is concerned with the formulation of a suitable notion of smooth
neighborhood deformation retract. For topological spaces, the statement that a
closed subspace A of X is a neighborhood deformation retract of X is equivalent
to the statement that the inclusion i : A →֒ X is a closed cofibration. We show
that in the category of Frölicher spaces there is a notion of smooth neighborhood
deformation retract that gives rise to an analogous result that a closed Frölicher
subspace A of the Frölicher space X is a smooth neighborhood deformation
retract of X if and only if the inclusion i : A →֒ X comes from a certain subclass
of cofibrations. As an application, we construct the right Puppe sequence.
3.1 SNDR pairs and SDR pairs
The definition of ‘smooth neighborhood deformation retract’ that we adopt in
this paper is similar to the definition of ‘R-SNDR pair’suggested in [6], but we
have modified the definition in order to retain only the essential aspects of ‘first
coordinate independence’ defined in [6].
We begin by defining the ‘first coordinate independence property’ of a func-
tion on a product of a Frölicher space with I (or I−, I+).
Definition 3.1 Let i : A → X be a smooth map, and c : R → X a structure
curve on X. Define
Λ(c, i) = {t∗ ∈ c−1(i(A))| there exists a sequence {tn} of real numbers
with limn→∞ tn = t∗ and each tn ∈ c−1(X − i(A))}.
The points in Λ(c, i) are those values in R where the curve ‘enters’ i(A) from
X − i(A), or ‘touches’ a point in i(A) whilst remaining in X − i(A) nearby.
Now, we are ready to define the ‘first coordinate independence property’ for a
structure function on a product.
Definition 3.2 Let i : A→ X be a smooth map and suppose f : I×X → R is
a structure function on I ×X. Let c : R → I ×X, given by c(s) = (t(s), x(s))
have the following properties
• The map x(s) is a structure curve on X.
• For all ǫ > 0, t(s) is a smooth real-valued function on R−∪s∗∈Λ(x,i)[s∗ −
ǫ, s∗ + ǫ].
If, for every such map c, the composite f ◦ c is a smooth real-valued function,
then we say that f : I×X → R has the first independence property (FCIP) with
respect to i.
Extending the definition, we say that a map g : I × X → Y has the FCIP
with respect to i if the composite h ◦ g : I ×X → R has the FCIP with respect
to i for every h ∈ FY .
Notice that we can formulate a similar definition of the FCIP if we replace
I throughout by I− or I+, leaving the rest of the definition unchanged. We will
have occasion to use this type of first coordinate independence property in the
later part of this work.
Note. Let i : A→ X , and suppose that we are given a map g : I×X → Y . Let
f : Y → R be a structure function on Y , and suppose that f ◦ g : I ×X → R
has the FCIP with respect to i for any such f . Then, given a smooth map
h : Y → Z, the composite f ′ ◦ h ◦ g : I×X → R has the FCIP with respect to
i for any structure function f ′ on Z.
The above note applies equally well if g : I− ×X → Y or g : I+ ×X → Y
has the FCIP with respect to i when composed with a smooth function h on Y .
Example 3.1
1. For any i : A→ X , the projection onto the second coordinate πX : I×X → X
has the FCIP.
2. Let α : R → R be a smooth braking function with the properties that
• α(t) = 0 if t < 1
• 0 < α(t) < 1 if 1
≤ t ≤ 3
• α(t) = 1 if 3
Consider 0 →֒ I−. Let H : I× I− → I− be given by H(t, s) = (1−α(t))s. Then,
f ◦H : I× I− → R has the FCIP with respect to the inclusion 0 →֒ I−, for any
f ∈ FI− .
Definition 3.3 Consider a smooth inclusion i : A →֒ X. Suppose that there
exists a smooth map u : X → I, with u−1(0) = i(A). If there exists a smooth
map H : I×X → X that satisfies the following properties:
• H has the FCIP with respect to i.
• H(0, x) = x for all x ∈ X.
• H(t, x) = x for all (t, x) ∈ I× i(A).
• H(1, x) ∈ i(A) for all x ∈ X with u(x) < 1,
then the pair (X,A) is called a smooth neighborhood deformation retract pair,
or SNDR pair for short.
If, in addition, H is such that H(1 × X) ⊂ i(A), then the pair (X,A) is
called a smooth deformation retract pair, or an SDR pair for short.
The subspace A is called a smooth neighborhood deformation retract or smooth
deformation retract of X if (X,A) is an SNDR pair or SDR pair, respectively.
The pair (u,H) is called a representation for the SNDR (or SDR) pair.
Example 3.2
1. The pair (X, ∅) is an SNDR pair. A representation is u(x) = 1, H(t, x) = x,
for each t ∈ I and x ∈ X .
2. The pair (X,X) is an SNDR pair. A representation is u(X) = 0, H(t, x) = x,
for each t ∈ I and x ∈ X .
Lemma 3.1 The pair (I−, 0) is an SDR pair.
Proof. Let α : R → R be the smooth braking function of Examples 3.1. A
representation for (I−, 0) as an SDR pair is (u,H), where u : I− → I and
H : I× I− → I− are given by u(s) = s, and H(t, s) = (1 − α(t))s. Clearly, the
identity u : I− → I is smooth. And the map H , as shown in Example 2.4, is
smooth and clearly has the FCIP with respect to the inclusion, since whenever
v approaches a value for which s(v) = 0, one has
g((1− α(t(v)))s(v)) = g(0)
for v in a neighborhood of this value and g ∈ FI− . �
Lemma 3.2 The pair (I, {0, 1}) is an SNDR pair.
Proof. A representation (u,H) for the SNDR pair can be given as follows.
Define u : I → I to be a bump function such that
• u(t) = 0 for t = 0 or t = 1,
• u(t) = 1 for t ∈ [ 1
• 0 < u(t) < 1 otherwise,
and let β : I → I be a braking function with the properties that β(s) = 0 for
0 ≤ s ≤ 1
, and β(s) = 1 for 3
≤ s ≤ 1. Let 0 < ǫ 1
, and define H : I× I → I by
H(t, s) = (1− αǫ(t))s+ αǫ(t)β(s). It is clear that H(0, s) = s, H(t, 0) = 0, and
H(t, 1) = 1. Suppose that u(s) < 1. Then, s ∈ [0, 1
) ∪ (3
, 1]. This implies that
β(s) = 0 or β(s) = 1. We then have H(1, s) = 0 or H(1, s) = 1, which means
that H(1, s) ∈ {0, 1} if u(s) < 1.
To see that H is smooth, let f : I → R be a generating function for the
flattened unit interval. The only possible points of non-smoothness are points
where t = 0, 1 and s = 0, 1. The braking function αǫ ensures that H is locally
constant in the tb variable whenever t is near 0 or 1, so no problem arises from
the t component. When s is near s = 0, we have H(t, s) near 0, and so the
generating function f is locally constant. Similarly, when s is near s = 1, we
have H(t, s) near 1, and the generating function f is again locally constant. �
We now show that the product of SNDR pairs is again an SNDR pair.
Theorem 3.1 Let i : A →֒ X and j : B →֒ Y be inclusion mappings. If (X,A)
and (Y,B) are SNDR pairs, then so is
(X × Y, (X ×B) ∪ (A× Y )).
If one of (X,A) or (Y,B) is an SDR pair, then so is the pair
(X × Y, (X ×B) ∪ (A× Y )).
Proof. Let α : R → I be a smooth braking function with the properties that
α(t) = 0 for t ≤ 1
, and α(t) = 1 for t ≥ 3
, and let β : R → R be a smooth
increasing braking function with the properties that β(t) = t for t ≤ 1
, and
β(t) = 1 for t ≥ 3
. Suppose that (u,H) and (v, J) are representations for the
SNDR pairs (X,A) and (Y,B), respectively. Let u : X → I, and v : Y → I be
given by u(x) = β(u(x)) and v(y) = β(v(y)) respectively. Define w : X×Y → I
by w(x, y) = u(x)v(y). The braking function β ensures smoothness of u and
v, and consequently of w. We have w−1(0) = (X × B) ∪ (A × Y ), as required.
Define Q : I×X × Y → X × Y as follows .
Q(t, x, y) =
(H(α(t), x), J(α(t), y)) if u(x) = v(y) = 0
(H(α(t), x), J(α(
)α(t), y)) if v(y) ≥ u(x), v(y) > 0,
(H(α(
)α(t), x), J(α(t), y)) if u(x) ≥ v(y), u(x) > 0.
We must show that Q is a smooth map, with the first coordinate independence
property with respect to the inclusion (X × B) ∪ (A × Y ) →֒ X × Y . We first
consider each part of the definition of Q separately. The first part is clearly
smooth. Let us verify that Q is smooth on the second part of its definition; the
third part is similar.
We need only focus on the component J(α(
)α(t), y). Each function
making up J(α(
)α(t), y) is smooth individually, so we need only pay extra
attention to those parts that involve flattened unit intervals, remembering that
addition and multiplication on the flattened unit interval need not preserve
smoothness, as is the case for the usual unit interval.
So let us consider α(
); it is smooth except possibly when
approaches
0 or 1, since it is here that structure curves on the flattened unit interval need
not be smooth in the usual sense. Clearly, if u(x) approaches 0 and v(y) does
not approach 0, then the braking function α ensures that
= 0 near such
points. If v(y) approaches 0, then u(x) must approach 0 too. This situation is
dealt with later.
Thus, Q, in part two of the definition, is smooth, and one can show similarly
that Q in the third part of the definition is smooth as well.
Let us now consider the overlaps of the three parts of the definition of Q.
Observe that if u(x) is in a sufficiently small neighborhood of v(y), with u(x) 6= 0
and v(y) 6= 0, then we have α(u(x)
) = 1, and so the second and third
parts of the definition of Q coincide here. Thus, it remains only to show that Q
is smooth as u(x) and v(y) both approach 0.
If Q is smooth in each of its coordinates then it is smooth, so consider the
coordinate involving the map J . Let c : R → I×X × Y be a structure that is
given by c(s) = (t(s), x(s), y(s)). Then, the map c1 : R → I× Y , given by
c1(s) =
(α(t(s)), y(s)) if u(x(s)) = v(y(s)) = 0
u(x(s))
v(y(s))
)α(t(s)), y(s)) if v(y(s)) ≥ u(x(s)), v(y(s)) > 0
(α(t(s)), y(s)) if u(x(s)) ≥ v(y(s)), u(x(s)) > 0
is a map satisfying the conditions of Definition 3.2, since its second coordinate is
smooth, but its first coordinate may be singular as v(y(s)) ( and hence u(x(s)))
approaches 0. Since J has the first coordinate independence property, the map
(Joc1)(s) =
J(α(t(s)), y(s)) if u(x(s)) = v(y(s)) = 0
u(x(s))
v(y(s))
)α(t(s)), y(s)) if v(y(s)) ≥ u(x(s)), v(y(s)) > 0
J(α(t(s)), y(s)) if u(x(s)) ≥ v(y(s)), u(x(s)) > 0
is smooth. Thus, Q ◦ c is smooth, and since c is arbitrary, Q is smooth. In a
similar way, the coordinate of Q involving H can be shown to be smooth.
We now verify that Q satisfies the required boundary conditions. When t =
0, all three lines defining Q reduce to (H(0, x), J(0, y)) = (x, y). Let x ∈ A and
y ∈ B; then u(x) = v(y) = 0. Therefore, Q reduces to (H(α(t), x), J(α(t), y)) =
(x, y). If x ∈ A and y /∈ B, then Q is given by the second part of its definition,
which reduces to (H(α(t), x), J(0, y)). The case when x /∈ A and y ∈ B is
similar. If t = 1 and 0 < w(x, y) < 1 then either 0 < u(x) < 1 or 0 < v(y) < 1.
Suppose that 0 < u(x) < 1. Then either u(x) ≤ v(y) or v(y) < u(x). If
u(x) ≤ v(y), then Q is given by the second part of its definition, which reduces
to (H(1, x), J(α(
, y)) ∈ i(A)× Y . If v(y) < u(x), then the third part of the
definition of Q applies and Q reduces to (H(α(
), x), J(1, y)) ∈ X × j(B).
Finally, we must show that for any f ∈ FX×Y , f ◦Q has the first coordinate
independence property with respect to the inclusion (X×B)∪(A×Y ) →֒ X×Y .
To this end, consider a map c : R → I×X×Y , given by c(s) = (t(s), x(s), y(s)).
Let {sn} be a sequence of real numbers converging to s∗ with c(sn) ∈ (X×Y )−
((A× Y ) ∪ (X ×B)), and c(s∗) ∈ (A× Y ) ∪ (X ×B). There are three cases to
consider.
• Suppose that c(s∗) ∈ A×B. Then x(s∗) ∈ A and y(s∗) ∈ B. The fact that
H and J have the first coordinate independence property with respect to
i and j respectively means that each coordinate of Q is smooth, and so Q
is smooth.
• Suppose that c(s∗) ∈ A × Y , and that y(s∗) /∈ B. Then at each of the
points c(sn), (Q ◦ c)(sn) is given by the second part of the definition of
Q, for n large enough. Since x(s∗) ∈ A, the component of Q involving H
is smooth, since H has the first coordinate independence property. For
any s in a neighborhood of s∗, α(
u(x(s))
v(y(s))
) = 0. Thus, the component of
Q involving J is constant for s in a neighborhood of s∗, and so is smooth
there.
• The case with c(s∗) ∈ X ×B, and x(s∗) /∈ A is similar to the second case
above.
For the last part of the theorem, suppose that (u,H) represent (X,A) as an
SDR pair. If we replace u by u′ = 1
u, then (u′, H) also represent (X,A) as an
SDR pair. Making the above constructions now with u′ in place of u, it follows
that w(x, y) < 1 for all (x, y) and so Q(1, x, y) ∈ (X × B) ∪ (A × Y ). This
completes the proof. �
4 Cofibrations
In this section, we show that for a subspace A ⊆ X that is closed in the under-
lying topology, the inclusion i : A → X is a cofibration if and only if (X,A) is
an SNDR pair.
Definition 4.1 Let i : A→ X be a cofibration. We call i a cofibration with
FCIP if any homotopy extension can be chosen to have the FCIP with respect
to i.
Using the equivalent formulation of the notion of cofibration, given by Lemma
2.5, we may restate Definition 4.1 as follows: A cofibration i : A → X is a
cofibration with the FCIP if and only if the map G that we may fill in to
complete the commutative diagram
(0×X) ∪ (I− ×A) h //
I− ×X
may be chosen to have the FCIP with respect to the inclusion i.
We have the following result, which corresponds to a similar topological
result.
Lemma 4.1 A smooth map i : A → X is a cofibration (with the FCIP) if
and only if (0 × X) ∪ (I− × A) is a retract of I− × X, (where the retraction
r : I− ×X → (0×X) ∪ (I− ×A) has the FCIP ).
Proof. In the one direction, suppose that (0 × X) ∪ (I− × A) is a retract of
I− ×X . We wish to complete the following diagram:
(0×X) ∪ (I− ×A) h //
I− ×X
By hypothesis, there exists r : I−×X → (0×X)∪ (I− ×A) such that r ◦ j = 1.
Define G = h ◦ r. If r has the FCIP, then so does h ◦ r.
Conversely, suppose that i : A → X is a cofibration (with the FCIP). We
may find a map r such that the diagram
(0×X) ∪ (I− ×A) 1//
(0 ×X) ∪ (I− ×A)
I− ×X
commutes. Thus, r ◦ j = 1. If i is cofibration with the FCIP with respect to i,
then r can be chosen to have the FCIP. �
The next theorem shows the relationship between cofibrations, retracts and
SNDR pairs.
Theorem 4.1 Let i : A → X be an inclusion, with A closed in the underlying
topology of X. Then the following are equivalent.
(1) The pair (X,A) is an SNDR pair.
(2) There is a smooth retraction r : I− × X → (0 ×X) ∪ (I− × A) with the
FCIP.
(3) The map i : A→ X is a cofibration with the FCIP.
Proof. To show that (1) and (2) are equivalent, note that the pair (I−×X, (0×
X) ∪ (I− × A)) is an SDR pair, as a consequence of Lemma 3.1 and Theorem
3.1. Let (w,Q) be a representation for the pair (I− ×X, (0×X)∪ (I− ×A)) as
an SDR pair, and let Q be constructed as in Theorem 3.1. Define
r : I− ×X → (0×X) ∪ (I− ×A)
by r(t, x) = Q(1, t, x), where (t, x) ∈ I− ×X . We observe that r has the FCIP,
since Q has this property, and Q has this property since each of its components
has this property.
The equivalence of (2) and (3) is Lemma 4.1.
We need only show that (2) implies (1). Let r : I−×X → (0×X)∪ (I−×A)
be a retraction with the FCIP with respect to i. Define H : I × X → X
by H(t, x) = (πX ◦ r)(α(t), x), where πX is the projection onto the second
coordinate, and α : R → R is a braking function with the following properties:
α(t) = 0 for t ≤ 0, α(t) = 1 for t ≥ 3
, and 0 < α(t) < 1 for 0 < t < 3
. This
braking function is necessary to ensure smoothness at the right endpoint of the
flattened unit interval I. Smoothness at the left endpoint is already taken care
of by the fact that r is defined in terms of the left flattened unit interval. The
map H satisfies the following properties:
• H has the FCIP since r has this property.
• H(0, x) = (πX ◦ r)(0, x) = x, for x ∈ X .
• H(t, x) = (πX ◦ r)(α(t), x) = x, for x ∈ A.
We now construct u : X → I. Let πI : I×X → I denote the projection onto I.
Define a smooth function β : R → R by
β(t) =
0 if t ≤ 0
t2 if t > 0.
Now, define u : X → I by
u(x) =
β(α(t) − (πI ◦ r)(1, x)(πI ◦ r)(α(t), x))dt
β(α(t))dt
It is clear that u is a smooth mapping.
We now verify that (u,H) represents (X,A) as an SNDR pair.
(1) Let x ∈ A. Clearly, (πI ◦ r)(1, x) = 1 and πI ◦ r)(α(t), x) = α(t), and so
β(α(t)− (πI ◦ r)(1, x)(πI ◦ r)(α(t), x))dt = 0. Thus, u(x) = 0, for all x ∈ A.
(2) Suppose that x ∈ X−A. Since 0×(X−A) is open in the underlying topology
on (0×X)∪ (I− ×A), we may choose an open neighborhood W ⊆ 0× (X −A)
of (0, x). Since r is continuous, there is a neighborhood V ⊆ I− ×X such that
r(V ) ⊆W ⊆ 0× (X −A). Now, consider the mapping qx : I → I×X , given by
qx(t) = (α(t), x), for each x ∈ X . This is clearly smooth. Thus, there exists a
neighborhood U ⊆ I− such that qx(U) ⊆ V . In other words, U × {x} ⊆ V . So,
we have (πI ◦ r)(α(t), x) = 0, for all t ∈ U . Thus, we have
u(x) =
β(α(t) − (πI ◦ r)(1, x)(πI ◦ r)(α(t), x))dt +
β(α(t))dt
β(α(t))dt
Combining this with part (1), we deduce that u−1(0) = A.
(3) Suppose that x is such that u(x) < 1. There must be a neighborhood U of I
such that (πI◦r)(1, x)(πI ◦r)(α(t), x) > 0, for t ∈ U . Thus (πI◦r)(1, x) > 0, but
this implies that r(1, x) ∈ I×A, and hence H(1, x) ∈ A. The proof is complete.
5 The Mapping Cylinder
In this section we show that the inclusion of X into the flattened mapping
cylinder If of a map f : X → Y is a cofibration with the FCIP.
Theorem 5.1 Let f : X → Y be a smooth map. Then, the pair (If , X) is an
SNDR pair.
Proof. Let α : I → R be a smooth braking function with the following proper-
ties: α(t) = 0 if 0 ≤ t ≤ 1
, α(t) = 1 if 3
≤ t ≤ 1, 0 < α(t) < 1, otherwise. Define
two more braking functions α1, α2 : I → R as follows: α1(0) = 0, 0 < α1(t) < 1
if 0 < t < 3
, α1(t) = 1 if
≤ t ≤ 1, and α2(t) = 0 if 0 ≤ t ≤ 34 , α2(t) = 1
≤ t ≤ 1. Now, define u : If → I by u([t, x]) = α1(t) and u([y]) = 1, for
(t, x) ∈ I×X and y ∈ Y . Define H : I× If → If by
H(s, [t, x]) = [(1 − α(s))t+ α(s)α2(t), x] if (t, x) ∈ I×X
H(s, [y]) = [y] if y ∈ Y .
That u is smooth comes from the fact that it is smooth when restricted to
each component of the coproduct (I×X)⊔Y ; it is thus smooth on the quotient
To see that the map H : I × If → If is smooth, note that since we are
working in a cartesian closed category, products commute with quotients, i.e. if
q is quotient, then so is 1× q, where 1 is an identity map. Thus, we may think
of H as being defined on the space
(I× I×X) ⊔ (I× Y )
where ∼ is the identification (t, 1, x) = (t, f(x)) for t ∈ I, and x ∈ X . Since H is
smooth when restricted to each component of the coproduct (I×I×X)⊔(I×Y ),
H is smooth on the quotient I× If .
We now verify that (u,H) is a representation for (If , X) as an SNDR pair.
• u−1(0) = [0, x] = i0(X).
• H(0, [t, x]) = [t, x] and H(0, [y]) = [y].
• H(s, [0, x]) = [0, x].
• If u[t, x] < 1, then t < 3
and so α2(t) = 0. Thus, H(1, [t, x]) = [0, x].
This completes the proof. �
Finally, we have the following important corollary.
Corollary 5.1 Given any smooth map f : X → Y , the inclusion X →֒ If is a
cofibration with the FCIP.
6 The Exact Sequence of a Cofibration
Our aim in this section is to show how one can use SNDR pairs to prove the
existence of the right exact Puppe sequence. We state the result in Theorem
6.1 and break the proof of the result up into a number of lemmas. We follow
the method used by Whitehead [12] for the topological case.
Throughout this section we work in the category FRL∗ of pointed Fr—’olicher
spaces, and basepoint preserving smooth maps.
Theorem 6.1 Let W be an object in FRL∗, and suppose that i : A →֒ X is a
cofibration in FRL∗. For any basepoint x0 ∈ A ⊆ X there is a sequence
. . . // [
A,W ]
Ti,W ]
X,W ]
A,W ] // . . .
. . . // [
A,W ]
// [Ti,W ]
// [X,W ]
// [A,W ]
which is an exact sequence in SETS∗, where j : X → Ti is the inclusion dis-
cussed in Paragraphe 2.4 and k : Ti →
A is the quotient map defined below.
It is, in fact, possible to prove that the sequence above is an exact sequence
of groups as far as
A,W ] and that the morphisms to this point are group
homomorphisms, but we shall not do so here.
The reduced(flattened)suspension of a pointed Frölicher space X is de-
fined as
X = (I/{0, 1}) ∧X,
where the reduced join is defined as for topological spaces with the identified
set taken as basepoint, and with 0 the basepoint of I.
In this section, whenever we refer to the suspension of a space , we mean
the reduced flattened suspension defined above.
Lemma 6.1 If (x,A) is an SNDR pair and p : X → X/A the quotient map,
then the sequence
i // X
p // X/A
is right exact.
Proof. To show that the given sequence is right exact we must show that for
any Frölicher space W the following sequence is exact in SETS:
[X/A,W ]
// [X,W ]
// [A,W ] .
It is easy to see that im p∗ ⊆ ker i∗. To see the reverse inclusion, let g : X →
W be an element of [X,W ], with g|A ≃ w0 (rel w0), where w0 ∈ W . Since
i // X is an SNDR pair, the map i is a cofibration, and so we may extend
w0 to a smooth map g
′ : X → W such that g′ ≃ g. But g′ is constant on A,
and so there exists a smooth map g1 : X/A → W such that p∗(g1) = g′. This
shows that ker i∗ ⊂ im p∗. �
Lemma 6.2 For any smooth map f : X → Y , the sequence
f // Y
l // Tf
is right exact, where l is the usual inclusion of Y into the mapping cone; i.e.
y 7→ [y] ∈ Tf .
Proof. One can show that there is a homotopy commutative diagram
i ��@
// Tf
where i, j, and l are the usual inclusions, and p is the quotient map that collapses
away {0} ×X to a point. Since, by Theorem 5.1, (If , X) is an SNDR pair, it
follows from Lemma 6.1 that the sequence
i // If
p // Tf
is right exact. It is fairly easy to show that j : Y → If is a homotopy equivalence.
Therefore, the sequence
f // Y
l // Tf
is right exact. �
Lemma 6.3 For any smooth map i : A → X, there is an infinite right exact
sequence
i // X
// Ti
// . . . i
// Tin−2
// Tin−1
// . . .
where in, n ≥ 1, are inclusion maps.
Proof. The pair (Ti, X) is an SNDR pair. The representation for the pair
(If , X) in Theorem 5.1 can be adapted to show this. One iterates the procedure
of Lemmas 6.1 and 6.2. �
One can easily see that there is an isomorphism between Ti/X and
Define q : Ti →
A to be the map which identifies X ⊂ Ti to a point, followed
by the isomorphism Ti/X →
Lemma 6.4 The sequence
// Ti
is right exact.
Proof. As noted above the pair (Ti, X) is an SNDR pair. We have the com-
mutative diagram
// Ti
where p : Ti → Ti/X is the identification map, and q0 : Ti/X →
A is an
isomorphism. The top line of the diagram is right exact, by Lemma 6.1, and so
the sequence
// Ti
is right exact. �
There is a commutative diagram
// Ti
where q1 is a homotopy equivalence. ( See Whitehead [12] for more details of
this map. ) Using commutative diagrams of this form, one can now proceed
almost exactly as one does in the topological situation, as in Whitehead [12] for
example, to get the following infinite right exact sequence:
i // X
// Ti
// . . .
. . . //
// . . .
The definition of right exactness now gives us the exact sequence of Theorem
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Preliminaries
Basic Constructions of Homotopy Theory in FRL
Our Approach to Homotopy Theory in FRL
Flattened Structures on the Unit Interval
Some Properties of Smooth Functions between the Flattened Unit Intervals
Homotopy in FRL and Related Objects
Cofibrations in FRL
Smooth Neighborhood Deformation Retracts
SNDR pairs and SDR pairs
Cofibrations
The Mapping Cylinder
The Exact Sequence of a Cofibration
|
0704.0343 | Experimental observation of structural crossover in binary mixtures of
colloidal hard spheres | Experimental observation of structural crossover in binary mixtures of colloidal hard
spheres
Jörg Baumgartl1,∗, Roel P.A. Dullens1, Marjolein Dijkstra2, Roland Roth3 and Clemens Bechinger1
12. Physikalisches Institut, Universität Stuttgart, 70550 Stuttgart, Germany
2Soft Condensed Matter Group, Utrecht University, 3584 CC Utrecht, The Netherlands
3Max-Planck-Institut für Metallforschung, 70569 Stuttgart, Germany and
Institut für Theoretische und Angewandte Physik, Universität Stuttgart, 70569 Stuttgart, Germany
Using confocal-microscopy we investigate the structure of binary mixtures of colloidal hard spheres
with size ratio q = 0.61. As a function of the packing fraction of the two particle species, we observe
a marked change of the dominant wavelength in the pair correlation function. This behavior is in
excellent agreement with a recently predicted structural crossover in such mixtures. In addition, the
repercussions of structural crossover on the real-space structure of a binary fluid are analyzed. We
suggest a relation between crossover and the lateral extension of networks containing only equally
sized particles that are connected by nearest neighbor bonds. This is supported by Monte-Carlo
simulations which are performed at different packing fractions and size ratios.
PACS numbers: 82.70.Dd, 61.20.-p
Most systems in nature and technology are mixtures
of differently sized particles. Each distinct particle size
introduces another length scale and its competition gives
rise to an exceedingly rich phenomenology in compari-
son with single-component systems. Already the simplest
conceivable multi-component system, i.e. a binary mix-
ture of hard spheres, exhibits interesting and complex
behavior. Just a few examples include entropy driven
formation of binary crystals [1, 2, 3], frustrated crys-
tal growth [4], the Brazil nut effect [5], glass-formation
[6, 7] and entropic selectivity in external fields [8]. Al-
though interaction potentials in atomic systems are more
complex than those of hard spheres, the principle of vol-
ume exclusion is ubiquitous and thus always dominates
the short-range order in liquids [9]. Accordingly, hard
spheres form one of the most important and successful
model systems in describing fundamental properties of
fluids and solids. It has been demonstrated that many
of their features can be directly transferred to atomic
systems where fundamental mechanisms are often ob-
structed by additional material specific effects [10]. Bi-
nary hard sphere systems are fully characterized by their
size ratio q = σS/σB with σi the diameters of the small
(S) and big (B) spheres and the small and big sphere
packing fractions ηS , ηB , respectively.
The pair-correlation functions, gij(r), are the central
measure of structure in fluids; they describe the probabil-
ity of finding a particle of size i at distance r from another
particle of size j. It is well known that all pair-correlation
functions in any fluid mixture with short-ranged inter-
actions (not just hard spheres) exhibit the same type
of asymptotic decay, which can be either purely (mono-
tonic) exponential or exponentially damped oscillatory
([11] and references therein). This prediction, which is
valid in all dimensions, suggests that all pair-correlation
functions decay with a common wavelength and decay
length in the asymptotic limit. For binary hard-sphere
mixtures where ηB � ηS or ηS � ηB , this is rather obvi-
ous since the system is dominated by either big or small
particles. The pair-correlation functions will asymptot-
ically oscillate with a wavelength determined either by
σB (ηB � ηS) or σS (ηS � ηB). Rather surprising is
that the above statement is also valid for all other rela-
tive packing fractions where the system is not dominated
by particles of a single size ([11, 12]). Accordingly, in the
asymptotic limit the (ηS , ηB) phase diagram is divided
by a sharp crossover line where the decay lengths of the
contributions to gij(r) with the two wavelengths become
identical. Below and above this line, however, the pair-
correlation function is either determined by the diameter
of the small spheres or that of the big spheres [13].
Despite the generic character of structural crossover
and the close relationship between structural and me-
chanical properties, this effect has not been observed in
experiments as the asymptotic limit is difficult to reach
in scattering experiments on atomic and molecular liq-
uids. However, recent calculations suggest that struc-
tural crossover is already detectable at relatively small
distances [12]. Because colloidal particles are directly
accessible in real space, such systems provide an oppor-
tunity to explore the structure of binary fluids and to
investigate structural crossover experimentally.
As colloidal suspension we used an aqueous binary
mixture of small melamin particles (σS = 2.9µm) and
big polystyrene spheres (σB = 4.8µm). Addition of
salt screens residual electrostatic interactions thus lead-
ing to an effective hard sphere system. Since melamin
has a higher density (ρM = 1.51g/cm3) than polystyrene
(ρP = 1.05g/cm3) the sedimentation velocities are sim-
ilar and, hence, we obtain a homogeneous system after
mixing. The suspension was contained in a cylindrical
sample cell with a silica bottom plate to allow optical
imaging with an inverted confocal microscope in reflec-
tion mode (Leica TCS SP2). From the images, particle
positions were obtained with digital video microscopy
[14]. Strong layering at the bottom wall allowed us to
image only the first two-dimensional bottom layer of the
three-dimensional system. We define the packing fraction
Figure 1: Different paths with constant total packing fraction
η = ηS + ηB in the (ηS , ηB)- plane. Experimental data (open
symbols: η = 0.72, q = 0.61) are sorted into ten bins. The
bin size is indicated by the ’error bars’. Closed symbols cor-
respond to the MC-simulations (N: η = 0.62, q = 0.4) and
(•: η = 0.57, q = 0.5). For convenience all samples are la-
beled with numbers increasing in the direction indicated by
the arrows.
as ηi = πσ2i /4, with ρi the number density of component
i. Variation of the relative packing fractions of the par-
ticles was achieved by addition of small particles to a
suspension of big spheres (Fig.1). Thus, the total pack-
ing fraction in the two-dimensional bottom layer remains
constant for all samples: η = 0.72. In the following we
will refer to the different samples by the sample numbers
(No.) as given in Fig.1.
Typical snapshots of the system for different packing
fractions of big and small particles are shown in Figs.2A-
C. The images demonstrate how the structure of the bot-
tom layer changes from being rich in small particles (No.
1, Fig.2A) to being rich in big particles (No. 10, Fig.2C).
Fig.2B (No. 5) corresponds to about the same number
density of small and big spheres. In order to analyze the
samples for a possible structural crossover, we calculated
the pair correlation function from the determined particle
positions. To minimize statistical noise we did not distin-
guish between big and small spheres. This is justified be-
cause the crossover has been predicted to be visible in all
pair-correlation functions and thus also in any linear com-
bination [11, 12]. The dominating wavelength in the os-
cillations is identified by computing the total correlation
function htot(r) =
i,j xixjhij(r) =
ij xixj [gij(r)−1],
with the mole fraction xi = ρi/
i ρi of component i [12].
Fig.2D exemplarily shows ln |htot(r)| for samples No. 1,5,
and 9. Note that in this representation the oscillation
wavelength is halved. The correlation functions of sam-
ples No.1 and 9 clearly oscillate with a single wavelength,
respectively, given by ≈ σB/2 and ≈ σS/2. In contrast,
sample 5 does not show a dominating wavelength but
an interference of different length scales which is typical
near the structural crossover. It is important to mention,
Figure 2: A-C) Typical snapshots of the bottom layer of a
binary mixture observed with a confocal microscope used in
reflection mode. The mixtures correspond to sample 10 (A), 5
(B) and 1 (C). The field of view is 40×40µm2. D) Logarithmic
plot of the total correlation functions htot(r) for the experi-
mental binary mixtures with η = 0.72 ± 0.04. Correlation
functions are plotted for sample numbers 1,5 and 9 (compare
Fig.1) and are shifted in vertical direction for clarity. The hor-
izontal bars correspond to σB/2 and σS/2, respectively. E)
Fourier-transforms of htot(r) for the experimental data points
(compare Fig. 1). Vertical lines indicate the wave vectors k
corresponding to the diameters of the small (S) and big par-
ticles (B), respectively. (color online).
that this intermediate behavior is only observed for sam-
ples No.5 and 6, i.e. only for about 10% of the entire
range over which ηB and ηS was varied. The experimen-
tally identified crossover-region is in excellent agreement
with the theoretically calculated value of ηS ≈ 0.3 at
those size ratios, which were determined from the decay
of the pair correlation functions calculated within density
functional theory in the test particle limit [15]. Fig.2E
Figure 3: Visualization of the different bond-types as determined by a Delaunay triangulation: big-big (black), big-small
(yellow) and small-small (red). Different plots correspond to the sample numbers as indicated in Fig.1. The field of view is
180× 180µm2.
shows the Fourier transforms of htot(r) for all samples
where the rather sudden change of the dominating wave-
length is seen more clearly [16]. At small and high pack-
ing fractions, the correlations are clearly dominated by
frequencies corresponding to either small or large parti-
cles (vertical lines) while around sample No.5 hardly any
dominating frequency is observed. This experimentally
confirms structural crossover as well as its occurrence at
finite particle distances.
So far, structural crossover has been discussed in terms
Figure 4: Averaged radii of gyration 〈Rig〉 (normalized to L/2
with L2 the size of the field of view) of networks formed by
large (solid symbols) and small particles (open symbols) as a
function of the sample number for A) the experimental data,
B) the MC-simulations at η = 0.57 and q = 0.5 and, C) the
MC-simulations at η = 0.62 and q = 0.4. The correspond-
ing packing fraction of small particles ηS is indicated as well.
The grey area and the dashed line respectively indicate the
crossover as inferred from the correlation functions and from
density functional theory. (color online).
of pair correlation functions, i.e. spatially averaged quan-
tities. Since our experiments naturally provide detailed
structural information, we investigate what the reper-
cussions are of the structural crossover on the real-space
structure. We first subjected a Delaunay triangulation to
the set of particle centers and identified nearest-neighbor
bonds between big-big (black), big-small (yellow), and
small-small (red) particles, respectively (see Fig.3). As
observed in Fig.3, sample 1 predominantly consists of big-
big bonds which form a large network spreading across
the entire field of view. With increasing sample No., i.e.
increasing ηS , the number of small-small bonds increases,
which leads to fragmentation of the big-big network into
smaller, randomly distributed patches. At large sample
numbers, the role of big and small particles is inverted
and small-small bonds form a network spanning the en-
tire area (No.10). Having distinguished between differ-
ent bond-types, a natural and well-known measure of the
spatial extend of a network formed by ni particles of size
i at positions ~xik (k = 1 . . . n
i) is given by the radius of
gyration Rig =
k=1(~x
k − ~R
2, with ~Ri0 the cen-
troid position of the network. Computing this quantity
for all, say N iC , networks formed by connected particles
of size finally yields a weighted averaged radius of gyra-
tion 〈Rig〉 =
m=1 ni(m)R
g(m) where N
i denote the
total number of particles i. We calculated 〈Rig〉 for net-
works consisting of connected big or small particles and
plotted these values for our experimental data in Fig.4A
as a function of the sample number. At small and high
sample numbers the quantities saturate while a relatively
sharp transition with an intersection point occurs around
sample 6. This location is indeed in very good agreement
with the crossover transition as determined from the cor-
relation functions in Fig.2 and density functional theory
(also indicated in Fig. 4A). This suggests that the struc-
tural crossover corresponds to a competition between the
sizes of networks consisting of connected big or small par-
ticles, respectively.
As structural crossover is also predicted for other size
ratios and packing fractions, we use Monte-Carlo (MC)
simulations to test our findings for more dilute systems
with size ratios q = 0.5 and q = 0.4. The corresponding
paths through the phase diagram (see closed symbols in
Fig.1) were obtained from 2-dimensional simulations with
a fixed number of particles of about 0 < N < 3000 for
both species and box areas of about 1500σ2B employing
periodic boundary conditions. From the configurational
snapshots we first determined the region of crossover by
analyzing htot(r) (the correlation functions are sampled
using 104 MC cycles per particle). Then, we performed
the above described Delaunay triangulation to calculate
〈Rig〉 for networks of connected big or small particles, re-
spectively. The corresponding radii of gyration are plot-
ted in Fig.4B and C and show a similar behavior as in the
experiment. Again, the intersection points are consistent
with the crossover region as inferred from the correlation
functions and DFT calculations. Note that the crossover
region sensitively depends on the size ratio and packing
fractions. Both the experiment and Monte-Carlo sim-
ulations show that structural crossover is accompanied
by a pronounced change in the typical size of networks
consisting of connected big and small particles. By in-
troducing small particles into a system of big spheres,
connections between big particles are broken and, at the
same time, connections between small particles are made.
This sensitively affects the typical size of networks con-
taining connected, equally-sized particles and thereby the
chance of finding another particle with the same size at a
relatively large distance. Consequently, the change from
〈RBg 〉 > 〈RSg 〉 to 〈RSg 〉 > 〈RBg 〉 (and vice versa) provides
a simple real-space argument why the oscillation wave-
length of the gij(r) in the asymptotic limit is either set
by σB or σS .
We have experimentally demonstrated the structural
crossover in a binary colloidal hard sphere system. Fur-
thermore, we show that structural crossover is strongly
coupled to the size of networks containing connected
equally-sized particles only. Going across the structural
crossover, the size ratio of such networks comprised by
either connected big or small particles is reversed. We be-
lieve this real-space configurational picture of structural
crossover is not just applicable to binary hard spheres,
as structural crossover is a generic feature of mixtures
with competing length scales. Moreover, it shows inter-
esting similarities with force chains in granular matter
[17] and glassy systems [6, 7, 18] of dissimilar sized parti-
cles. Therefore, our finding may help to gain more insight
into structure-related properties in binary systems at an
universal level.
∗Electronic address: [email protected]
stuttgart.de
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References
|
0704.0344 | The Blazar Spectral Sequence and GLAST | The Blazar Spectral Sequence and GLAST
L. Maraschi, G. Ghisellini and F. Tavecchio
INAF-Osservatorio Astronomico di Brera, Milano, Italy
Abstract. The present status and understanding of the "spectral sequence" of blazars is discussed in the perspective of the
upcoming GLAST launch. The vast improvement in sensitivity will allow to i) determine more objectively the "average"
gamma-ray properties of classes objects ii) probe more deeply the ratio between accretion power and jet power in different
systems.
Keywords: Gamma-rays - Relativistic jets - Galaxies: active
PACS: 95.85.Pw; 98.54.Cm
INTRODUCTION
The spectral sequence of blazars (Fossati et al.,1998 al.,1998) was constructed merging three complete blazar samples
(two radio selected, one X-ray selected: 2 Jy FSRQ, Wall & Peacock 1985, 1Jy BL Lac, Kuhr et al. 1981, and Slew
Survey BL Lac, Elvis et al. 1992), grouping all the objects in radio luminosity bins and averaging monochromatic
luminosities of objects within each radio-luminosity bin. The procedure is thus prone to various biases (Maraschi &
Tavecchio 2001), in particular the gamma-ray data were largely incomplete.
The resulting "sequence" shows that the blazar SEDs are double humped and that the two peaks shift to higher
energies with decreasing luminosity. Systematic modelling of the SEDs of individual objects (Ghisellini et al. 1998)
yields basically uniform beaming factors and jet parameters varying along the sequence in the sense of an increasing
energy density and decreasing electron critical energy at higher luminosities. Thus the "sequence" offers a suggestive
indication that the basic spectral properties of blazar jets could be related to the different powers involved and possibly
represent an evolutionary sequence in cosmic history (Boettcher and Dermer 2002; Cavaliere and D’Elia 2002).
The validity of the sequence concept has been questioned on the basis of deeper and larger blazar surveys (e.g.
Giommi et al. 2005, Padovani 2007) which however lack until now the very important gamma-ray data.
Here we wish to address two points. The first concerns the validity of the original claim within the presently known
bright blazar SEDs, the second concerns an anticipation of the types of blazars that may be detected by GLAST.
NEW DATA / NEW SOURCES
Given the limited space we will illustrate our points schematically, commenting few representative figures. All the
figures will have in the background the double humped lines interpolating the blazar spectral sequence. The latter are
just polinomial expressions connecting the average monochromatic luminosities obtained as described above.
The SED of a new high redshift FSRQ serendipitously discovered by SWIFT (BAT) J0746+2548 (z=2.979)
(Sambruna et al. 2006) is shown in Fig. 1a. Clearly J0746 is extremely luminous and conforms well to the sequence,
possibly suggesting a gamma-ray peak at Mev energies. The spectral shape in the gamma-ray band that will be
measured by GLAST for a large number of blazars will provide an essential information to constrain the position
of the high energy peak of blazar SEDs thus probing the sequence concept.
3C 454.3 is a highly variable FSRQ (z=0.859) already detected in gamma-rays by EGRET. The data for a "normal"
state (Tavecchio et al. 2007) are shown in Fig. 1b. This source could be detected with GLAST at 1% the intensity
level shown in the figure which is the average of EGRET measurements. The source underwent a strong outburst
recently and was observed by SWIFT (BAT) and INTEGRAL up to more than 100 keV (Pian et al. 2006, Giommi et
al. 2006). In the latter state the expected gamma-ray flux could have been an order of magnitude brighter than detected
by EGRET. A source with an intrinsically similar jet could then be detected in gamma-rays even if the jet was at a
larger angle to the line of sight. The thick lines in Fig. 2 represent the model used to describe the "normal" state of 3C
454.3, computed for different viewing angles. The gamma-ray emission could be detected by GLAST up to an angle
http://arxiv.org/abs/0704.0344v1
FIGURE 1. Spectral Energy Distribution of the blazars J0746+2548 (left, from Sambruna et al. 2006) and 2251+158 (right, from
Tavecchio et al. 2007) overimposed on the curves interpolating the blazar sequence. For 2251+158 we also report the model used
to reproduce the data (upper black curve) and the emission expected for a misaligned jet with angles of respectively 6, 8 and 10
degrees (from top to bottom).
of 10 degrees to the jet axis. In this case the SED would be significantly different than expected from the sequence,
simply because the jet emission is less beamed and less prominent with respect the SED of the accretion disk, included
here as a blackbody component plus a Seyfert like X-ray component. The sequence is not expected to extend to objects
with jets seen at intermediate angles. The different Doppler factor causes only a linear shift of the peak position but a
dramatic change in luminosiy.
Fig. 2 is devoted to blazars with lower luminosities. This part of the sequence is populated exclusively by BL
Lac objects defined as HBLs due to their SEDs peaking at high energies, in the X-ray and TeV bands. In Fig. 2a
the data for the "normal" state of PKS 2155-304 are plotted in green. They are well consistent with the sequence.
The multifrequency data obtained during the exceptional TeV flare observed from this source in July August 2006
are also shown (see Foschini et al. 2007). During the outburst the two emission peaks do not appear to shift much
in frequency but the luminosities increase by a large factor (for a short time) especially in the TeV band. Thus the
high state SED deviates remarkably from the sequence expectations. For these objects, though relatively weak at GeV
energies, GLAST observations will be important to define the shape of the high energy peak and its possible evolution
during outbursts.
Finally, in Fig. 2b we show the data for 1629+4008 (z=0.272), a blazar with an emission peak between the UV and
the X-ray band discovered within a survey aimed at finding objects with anomalous properties (Padovani et al. 2002).
The SED of this source complies reasonably well with the sequence expectation for an HBL, however this object
shows emission lines which is not the case for HBLS. In fact the sequence included only X-ray selected BL Lacs, but
no X-ray selected radio-loud objects with emission lines, as no such complete sample was available at the time (see
Wolter & Celotti 2001).
This source indicates that jets with SEDs peaking at high energies can occur in emission line AGNs. This is a new
result, which however does not break the correlations inferred from the sequence, as it occurs in the low luminosity
range. The question then is: what distinguishes HBLs from objects like 1629? Why emission lines are completely
absent in HBLs but present in 1629 whose jet is of comparable luminosity? According to our ideas (Maraschi 2001,
Maraschi & Tavecchio 2003) HBL should accrete at highly subEddington rates, therefore in the radiatively inefficient
accretion (RIAF) regime, while 1629, which shows emission lines, should be in the “standard” accretion disk regime,
therefore near to its Eddington limit. This in turn implies that this source contains a central black hole of relatively
modest mass. From the accretion luminosity, assuming that it corresponds to 0.1 the Eddington luminosity we can infer
a mass of 6× 107 solar masses. More direct estimates of the black hole mass are needed to confirm this prediction.
FIGURE 2. SEDs of the blazars 2251-304 (left, Foschini et al. 2007) and 1629+4008 (Padovani et al. 2002) overimposed on the
blazar sequence interpolations. For PKS 2155-304 a normal state is shown together with optical/X-ray and TeV data during the
exceptional outburst of July-August 2006
CONCLUSIONS
The few examples discussed above are meant to indicate how the concept of a spectral sequence for blazars, based
on averages over limited samples involving only the brightest objects of each class, may be probed by GLAST. In
particular, strong emphasis has been put in the past on BL Lac objects, neglecting the X-ray selected counterparts of
FSRQ which may also be gamma-ray emitters. GLAST is expected to produce extraordinary advances in this field.
It will increase by orders of magnitude the number of objects with measured gamma-ray flux (see Dermer these
proceedings) thus allowing to study deeper and differently selected samples. These will certainly contain "mixed"
objects in which the jet emission is less prominent in comparison to other AGN properties. The new gamma-ray
populations should carry great potential for understanding the link between accretion power and the production of jets
in extragalactic objects.
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http://arxiv.org/abs/astro-ph/0610545
http://arxiv.org/abs/astro-ph/0703359
Introduction
New data / new sources
Conclusions
|
0704.0345 | A High Robustness and Low Cost Model for Cascading Failures | epl draft
A High Robustness and Low Cost Model for Cascading Failures
Bing Wang and Beom Jun Kim
Department of Physics, BK21 Physics Research Division, and Institute of Basic Science, Sungkyunkwan University,
Suwon 440-746, Korea
PACS 89.75.Hc – Networks and genealogical trees
PACS 05.10.-a – Computational methods in statistical physics and nonlinear dynamics
PACS 89.20.Hh – World Wide Web, Internet
PACS 89.75.Fb – Structures and organization in complex systems
Abstract. - We study numerically the cascading failure problem by using artificially created
scale-free networks and the real network structure of the power grid. The capacity for a vertex
is assigned as a monotonically increasing function of the load (or the betweenness centrality).
Through the use of a simple functional form with two free parameters, revealed is that it is indeed
possible to make networks more robust while spending less cost. We suggest that our method to
prevent cascade by protecting less vertices is particularly important for the design of more robust
real-world networks to cascading failures.
The network robustness has been one of the most central
topics in the complex network research [1]. In scale-free
networks, the existence of hub vertices with high degrees
has been shown to yield fragility to intentional attacks,
while at the same time the network becomes robust to
random failures due to the heterogeneous degree distribu-
tion [2–5]. On the other hand, for the description of dy-
namic processes on top of networks, it has been suggested
that the information flow across the network is one of the
key issues, which can be captured well by the betweenness
centrality or the load [6].
Cascading failures can happen in many infrastructure
networks, including the electrical power grid, Internet,
road systems, and so on. At each vertex of the power
grid, the electric power is either produced or transferred
to other vertices, and it is possible that from some reasons
a vertex is overloaded beyond the given capacity, which is
the maximum electric power the vertex can handle. The
breakdown of the heavily loaded single vertex will cause
the redistribution of loads over the remaining vertices,
which can trigger breakdowns of newly overloaded ver-
tices. This process will go on until all the loads of the
remaining vertices are below their capacities. For some
real networks, the breakdown of a single vertex is suffi-
cient to collapse the entire system, which is exactly what
happened on August 14, 2003 when an initial minor distur-
bance in Ohio triggered the largest blackout in the history
of United States in which millions of people suffered with-
out electricity for as long as 15 hours [7]. A number of as-
pects of cascading failures in complex networks have been
discussed in the literature [8–16], including the model for
describing cascade phenomena [8], the control and defense
strategy against cascading failures [9, 10], the analytical
calculation of capacity parameter [11], and the modelling
of the real-world data [12]. In a recent paper [16], the cas-
cade process in scale-free networks with community struc-
ture has been investigated, and it has been found that
a smaller modularity is easier to trigger cascade, which
implies the importance of the modularity and community
structure in cascading failures.
In the research of the cascading failures, the following
two issues are closely related to each other and of signif-
icant interests: One is how to improve the network ro-
bustness to cascading failures, and the other particularly
important issue is how to design manmade networks with
a less cost. In most circumstances, a high robustness and a
low cost are difficult to achieve simultaneously. For exam-
ple, while a network with more edges are more robust to
failures, in practice, the number of edges is often limited
by the cost to construct them. In brevity, it costs much to
build a robust network. Very recently, Schäfer et. al. pro-
posed a new proactive measure to increase the robustness
of heterogeneous loaded networks to cascades. By defin-
ing the load dependent weights, the network turns to be
more homogeneous and the total load is decreased, which
means the investment cost is also reduced [15]. In the
present Letter, for simplicity, we try to find a possible way
of protecting networks based on the flow along shortest-
http://arxiv.org/abs/0704.0345v1
B. Wang B.J. Kim
l/lmax
This work
ML model in Ref.[8]
Fig. 1: The capacity c is assigned as c = λ(l)l with the initial
load l. The step function λ(l) = 1 + αΘ(l/lmax − β) with two
free parameters α and β is used in our model. For comparison,
the curve for the Motter-Lai (ML) capacity model in Ref. [8],
where λ(l) = constant, is also shown.
hop path, first proposed by Motter-Lai [8]. Through the
use of our improved capacity model, we numerically exam-
ine the cascades in scale-free networks and the electrical
power grid network. Since for heterogeneously loaded net-
works, overload avalanches can be triggered by the failure
of only one of the most loaded vertices, the following re-
sults are all based on the removal of one vertex with the
highest load. Our results suggest that networks can indeed
be made more robust while spending less cost.
We first construct the Barabási-Albert (BA) scale-free
network [17] of the size N = 5000 with the average degree
〈k〉 ≈ 4 to study the cascading failures. The BA network is
characterized by the degree distribution p(k) ∼ k−γ with
the degree exponent γ = 3, and it has been shown that the
load distribution also exhibits the power-law behavior [6],
which means that there exist a few vertices with very large
loads.
The betweenness centrality for each vertex, defined as
the total number of shortest paths passing through it, is
used as the measure of the load and computed by using
the efficient algorithm [18]. The capacity cv for the vertex
v is assigned as
cv = λ(lv)lv, (1)
where lv is the initial load without failed vertices. Al-
though it should be possible to find, via a kind of the
variational approach, the optimal functional form of λ(lv)
which gives rise to the lower cost and the higher robust-
ness (see below for the definitions of the two) we in this
work simplify λ(lv) as shown in Fig. 1:
λ(lv) = 1 + αΘ(lv/lmax − β), (2)
where Θ(x) = 0(1) for x < 0(> 0) is the Heaviside step
function, lmax = maxv lv, and we use α ∈ [0,∞) and β ∈
[0, 1] as two control parameters in the model. In Ref. [8]
a constant λ has been used (see Fig. 1 for comparison),
which corresponds to the limiting case of β = 0 with the
identification λ = 1 + α in our model.
At the initial time t = 0, the vertex with the highest
load is removed from the network, and then new loads
for all other vertices are recomputed.1 We then check the
failure condition cv < lv(t) for each vertex, and remove
all overloaded vertices to get the network at t + 1. The
above process continues until all existing vertices fulfill the
condition cv > lv(t), and the size of the giant component
N ′ at the final stage is measured. The relative size of the
cascading failures is conveniently captured by the ratio [8]
, (3)
which we call the robustness from now on. For networks
of homogeneous load distributions, the cascade does not
happen and g ≈ 1 has been observed [8]. Also for net-
works of scale-free load distributions, one can have g ≈ 1
if randomly chosen vertices, instead of vertices with high
loads, are destroyed at the initial stage [8].
In general, one can split, at least conceptually, the to-
tal cost for the networks into two different types: On the
one hand, there should be the initial construction cost to
build a network structure, which may include e.g., the cost
for the power transmission lines in power grids, and the
cost proportional to the length of road in road networks.
Another type of the cost is required to make the given
network functioning, which can be an increasing function
of the amount of flow and can be named as the running
cost. For example, we need to spend more to have big-
ger memory sizes and faster network card and so on for
the computer server which delivers more data packets. In
the present Letter, we assume that the network structure
is given, (accordingly the construction cost is fixed), and
focus only on the running cost which should be spent in
addition to the initial construction cost.
Without consideration of the cost to protect vertices,
the cascading failure can be made never to happen by
assigning extremely high values to capacities. However,
in practice, the capacity is severely limited by cost. We
expect the cost to protect the vertex v should be an in-
creasing function of cv, and for convenience define the cost
λ(lv)− 1
/N. (4)
It is to be noted that for a given value of α, the original
Motter-Lai (ML) capacity model in Ref. [8] has always a
higher value of the cost than our model (see Fig. 1). Al-
though e = 0 at β = 1, it should not be interpreted as
a costfree situation; we have defined e only as a relative
measure in comparison to the case of λ(l) = 1 for all ver-
tices. For a given network structure, the key quantities
to be measured are g(α, β) and e(α, β), and we aim to in-
crease g and decrease e, which will eventually provide us
1In real situations of failures, the initial breakdown can happen
at any vertex in the network. However, the eventual scale of dam-
ages must be greater when a heavily loaded vertex is broken, and
accordingly we in this work restrict ourselves to the worst case when
the vertex with the highest load is initially broken.
A High Robustness and Low Cost Model for Cascading Failures
0.002 0.003 0.004
α =1.00
=0.30
=0.25
=0.20
=0.15
=0.10
0.002 0.003 0.004
(b) α =0.30
=0.25
=0.20
=0.15
=0.10
0 0.2 0.4 0.6 0.8 1
α =0.30
=0.25
=0.20
=0.15
=0.10
Fig. 2: Cascading failures in the BA network of the size N =
5000 and the average degree 〈k〉 ≈ 4, triggered by the removal
of a single vertex with the highest load. The robustness g
and the cost e in Eqs. (3) and (4) are shown in (a) and (b),
respectively, as functions of β at various α values [see Fig. 1 for
α and β, the two parameters in the function λ(l) in Eq. (2)].
(c) The relation between e and g at different α’s. Compared
with the ML model in Ref. [8], it is clearly shown that the
network can be made more robust but with less cost.
a way to achieve the high robustness and the low cost at
the same time.
In Fig. 2(a), we report the robustness g for the BA net-
work of the size N = 5000 with the average degree 〈k〉 ≈ 4
as a function of β at α = 0.10, 0.15, 0.20, 0.25, 0.30, and
1.0 (from bottom to top). As β increases further beyond
the region in Fig. 2(a), the robustness g is found to de-
crease toward zero (not shown here), which is as expected
since the larger β makes vertices with larger loads less pro-
tected (see Fig. 1). We also skip in Fig. 2 small values of
β below approximately 0.001: If β < lmin/lmax, with the
minimum load lmin, all vertices are given λ(l) = 1 + α,
equivalent to the ML model corresponding to β = 0. It
is shown in Fig. 2(a) that for α . 0.30, g first increases
and then decreases as β is increased, exhibiting a well-
developed maximum gmax at β = β
∗. This is a partic-
ularly interesting observation since the network becomes
more robust (larger g) by protecting less vertices (larger
β). In more detail, the curve for α = 0.20 in Fig. 2(a)
shows the maximum gmax ≈ 0.62 (at β
∗ ≈ 0.00133), which
is about 3.5 times bigger than g ≈ 0.175 (at β = 0). In
other words, the network can be made much more robust
by assigning smaller capacities to vertices with less loads.
For larger values of α, on the other hand, it is found that
gmax occurs at β = 0, which indicates that the above find-
ing, i.e., possibility of making network more robust by
protecting less vertices, does not hold, as exemplified by
the curve for α = 1 in Fig. 2(a).
The above observation is closely related with Ref. [9],
where it has been found that in order to reduce the size
of cascades (or to have a larger g), some of less loaded
vertices should be removed just after the initial attack. In
reality, however, we believe that the direct application of
this strategy of intentional breakdowns is not easy, for cas-
cading failures usually propagate across the whole network
very soon just after the initial breakdown. In contrast, we
propose in this work a way to make the network better
prepared to breakdowns, by protecting less vertices.
In order to look at the cost benefit of protecting less
vertices in a more careful way, we plot in Fig. 2(b) the
cost e in Eq. (4) versus β at various values of α. As is
expected from Fig. 1, the cost e is shown to be a mono-
tonically decreasing (increasing) function of β (α) at fixed
α (β). Take again the case with α = 0.20 as an exam-
ple with e(β∗) ≈ 0.153 and e(β = 0) = 0.2: It is then
concluded that for α = 0.2 one can make the network
3.5 (≈ 0.62/0.175) times more robust while spending only
76.5% (≈ 0.153/0.2) of the original cost.
In Fig. 2(c), we use the same data as in Fig. 2(a) and
(b), and show the relation between the robustness and
the cost for α = 0.10, · · · , 0.30 from bottom to top. For
comparison, the values (g,e) for β = 0, corresponding to
the ML model, are also displayed as symbols at the end
of curves. It is clearly shown that for a given α, one can
achieve the higher robustness and the lower cost by tuning
β toward the right-most point on each curve. We can also
use Fig. 2(c) to choose the most efficient way to get a
given robustness g: For example, suppose that g = 0.6
is the required robustness. The vertical line for g = 0.6
crosses several different curves, and one can choose the
crossing point which has the lowest cost.
We next study the cascading failures in the real net-
work structure of the North American power grid of the
size N = 4941 [19]. Although the electrical power grid
network is a very homogeneous network in terms of the
degree distribution, the load distribution, in a sharp con-
trast, shows a strong heterogeneity as shown in Fig. 3. In
other words, the degree distribution is more like an ex-
ponential one, while the load distribution is similar to the
power-law form. The broad load distribution can be one of
the reasons of the fragility of the power grid to cascading
failures [8].
We then apply, the same method as we used above, to
the power grid, and obtain g and e as functions of β for
B. Wang B.J. Kim
104 105 106
0 5 10 15 20
Fig. 3: The cumulative load distribution of power grid network
P (l) in log-log scale. The inset shows the cumulative degree
distribution P (k) of the power grid in linear-log scale.
0 0.01 0.02 0.03 0.04 0.05
α =1.0
=0.8
=0.4
=0.2
=0.1
0 0.01 0.02 0.03 0.04 0.05
α =1.0
=0.8
=0.4
=0.2
=0.1
0 0.2 0.4 0.6 0.8
α =1.0
=0.8
=0.4
=0.2
=0.1
Fig. 4: Cascading failures in the electrical power grid of the size
N = 4941. (Compare with Fig. 2 for the corresponding plots
for the BA network.) The robustness g and the cost e versus
β at various α values are shown in (a) and (b), respectively,
while (c) is for the relation between e and g. Again, it is
shown that one can achieve the higher robustness and the less
cost simultaneously, by choosing the right-most point in (c).
0.001 0.002 0.003 0.004
=0.0
=0.2
=0.4
=0.6
=0.8
=1.0
Fig. 5: Cascading failures in the BA network of the size N =
5000 and the average degree 〈k〉 ≈ 4, triggered by the removal
of a single vertex with the highest load. Each vertex’s capacity
is disturbed with probability ε for α = 0.2. The data are
averaged over 20 runs.
given values of α. Figure 4 for the cascading failures of
the power grid is in parallel to Fig. 2 for the BA network:
Fig. 4(a) for g versus β, (b) for e versus β, and (c) for e
versus g. There are some quantitative differences between
curves for the power grid and the BA network. However,
qualitatively speaking, both networks are shown to ex-
hibit the following common features: (i) For a given α,
the robustness has a maximum gmax at β = β
∗, (ii) e is a
monotonically decreasing function of β at a given α, and
(iii) there exists a lob-like structure in the g-e plane, which
indicates that one can make the network exhibit a higher
robustness and a lower cost at the same time than the cor-
responding values for the ML model. It is worth mention-
ing that the power grid in Fig. 4 can be made to show the
higher g and the lower e than the ML model in a broader
region of α: Even at α = 1, the power grid can have much
better robustness and much less cost in comparison to the
ML model. Specifically, at α = 1.0 the ML model has
g ≈ 0.40 and e = 1.0 while our model can yield g ≈ 0.73
and e ≈ 0.26 (at β ≈ 0.00583) [see Fig. 4(c)], which occurs
when only 26% of vertices are given the higher capacity
λ(l) = 2, and the other remaining 74% of vertices have
the lower capacity λ(l) = 1. In other words, by assigning
lower capacities to 74% of vertices, the network becomes
much more robust.
In reality, it is also interesting to observe the effect of
noise on the dynamical process. In Ref. [20], when noise
is introduced into the nonlinear dynamical system, it has
been shown that noise changes the singularity at a special
time to a statistical time distribution and shows various in-
teresting behaviors. In the present work, we are interested
in how the presence of noise influences the final cascading
failure behavior within our scheme. Here, we introduce
A High Robustness and Low Cost Model for Cascading Failures
effects of noise as an erroneous assignment of the capac-
ity function. In detail, at a given error probability ε, the
vertex v is assigned the capacity c′v instead of its correct
c′v = cv(1 + r), (5)
where r is the uniform random variable with zero mean
(r ∈ [−1, 1]). We believe that this erroneous behavior
is plausible in reality, since the perfect knowledge for the
true value of the load for each vertex may not be available,
which may cause an erroneous assignment of the capacity
on a vertex. In the limiting case of ε = 0, we recover our
error-free results presented above. In Fig. 5, we report the
results at α = 0.2 for the robustness g for the BA network
as a function of β for different error probability ε [see
Fig.2(a) for comparison]. It is seen that for small ε, the
overall behavior is qualitatively the same as in Fig. 2(a),
i.e., the existence of a well-developed robustness peak and
gradual decrease as β is increased. The peak height of the
robustness is found to decrease as ε is increased, indicating
the negative effect of the noise. An interesting observation
in Fig. 5 is that as ε becomes larger there exits a region
of β in which the robustness is actually higher than the
error-free case of ε = 0.
In summary, we have suggested a new capacity model to
cascading failures, by improving the existing ML capacity
model in Ref. [8]. The main idea in our model is the same
as in existing studies: In a highly heterogeneous network
with a broad load distribution, vertices with large loads
should be more protected by assigning large capacities.
Different from other studies in which the capacity is as-
signed in proportion to the load, i.e., c = λl, we generalize
the model so that the proportionality constant λ is now
changed to an increasing function λ(l) of l. In more detail,
we use the Heaviside step function for λ(l) characterized
by two parameters, the step height α, and the step posi-
tion β. By applying this capacity model to the artificial
BA network as well as the real network of the power grid,
we have clearly shown that it is indeed possible to make
the network more robust, while at the same time the cost
to assign capacities is drastically reduced. We believe that
our suggested model to assign capacities to vertices should
be practically useful in designing infrastructure networks
in an economic point of view. As a final remark, it needs
to be pointed out that the model proposed in this work
should be considered as only the first step to find the op-
timal functional form λ(l) of the capacity as a function of
the load. As a future work, we are planning to apply a
sort of variational method to find the optimal functional
form of λ(l).
B.J.K. was supported by grant No. R01-2005-000-
10199-0 from the Basic Research Program of the Korea
Science and Engineering Foundation.
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http://arxiv.org/abs/cond-mat/0503615
|
0704.0346 | Diffuse X-ray Emission from the Carina Nebula Observed with Suzaku | Diffuse X-ray Emission from the Carina Nebula Observed with
Suzaku
Kenji Hamaguchi1,2, the Suzaku η Carinae team and the Carinae D-1 team
1CRESST and X-ray Astrophysics Laboratory NASA/GSFC, Greenbelt, MD 20771
2Universities Space Research Association, 10211 Wincopin Circle, Suite 500,
Columbia, MD 21044
A number of giant HII regions are associated with soft diffuse X-ray emission. Among
these, the Carina nebula possesses the brightest soft diffuse emission. The required plasma
temperature and thermal energy can be produced by collisions or termination of fast winds
from main-sequence or embedded young O stars, but the extended emission is often observed
from regions apart from massive stellar clusters. The origin of the X-ray emission is unknown.
The XIS CCD camera onboard Suzaku has the best spectral resolution for extended
soft sources so far, and is therefore capable of measuring key emission lines in the soft band.
Suzaku observed the core and the eastern side of the Carina nebula (Car-D1) in 2005 Aug and
2006 June, respectively. Spectra of the south part of the core and Car-D1 similarly showed
strong L-shell lines of iron ions and K-shell lines of silicon ions, while in the north of the core
these lines were much weaker. Fitting the spectra with an absorbed thin-thermal plasma
model showed kT∼0.2, 0.6 keV and NH∼1−2×10
21 cm−2 with a factor of 2-3 abundance
variation in oxygen, magnesium, silicon and iron. The plasma might originate from an old
supernova, or a super shell of multiple supernovae.
§1. Extended X-ray Emission from the Star Forming Region
Soft X-ray emission nebulae with kT∼0.1–0.8 keV, log LX∼33-35 ergs s
−1, and
size of ∼1–103 pc accompany a number of giant HII regions (see Table 4 of Ref. 6).
Chandra observations of extended emission in a few star forming clusters indicate
that the emission may arise from the fast O star stellar winds thermalized either by
wind-wind collisions or by a termination shock. However, the emission is often found
outside of the massive stellar clusters, so that another origin, such as an otherwise
unrecognized supernova remnant, cannot be ruled out.
In principle, the origin of the diffuse emission can be determined by measuring
its composition. For example, the plasma should be overabundant in nitrogen and
neon if it originates from winds from nitrogen-rich Wolf-Rayet stars (WN), while it
would be overabundant in oxygen if it arises from a Type II SNR. The temperature
of the plasma, typically a few million degrees, makes soft X-ray band studies highly
desirable, because of the presence in this band of strong lines from these elements,
plus carbon, silicon and iron.
The Carina Nebula, which contains several evolved and main-sequence massive
stars such as η Car, WR 25 and massive stellar clusters such as Trumpler 14 (Tr
14), emits soft diffuse X-rays 10–100 times stronger than any other Galactic giant
HII region (LX ∼10
35 ergs s−1).4) The high surface brightness made possible the
discovery of the diffuse emission by the Einstein Observatory in the late 1970’s.
The Einstein observations revealed that the diffuse emission tends to be associated
typeset using PTPTEX.cls 〈Ver.0.9〉
http://arxiv.org/abs/0704.0346v1
2 K. Hamaguchi et al.
with optically bright regions containing massive stars. Recent Chandra observations
provided a point source free measurement of the diffuse flux,1) and suggested the
presence of a north-south Fe and Ne abundance gradient.5)
The X-ray CCD cameras (XISs: X-ray Imaging Spectrometer) onboard the
Suzaku observatory have the best spectral resolution for extended soft X-ray emission
and thus they provide good diagnostics of emission lines especially below ∼1 keV.
§2. Suzaku and XMM-Newton Observations of the Carina Nebula
Figure 1 shows a mosaic image of the Carina nebula between 0.4−7 keV created
from 32 XMM-Newton observations. The image depicts several bright X-ray point
sources: η Car (an LBV), WR25, WR22 (Wolf-Rayet stars), HD 93250, HD 93043
(O3 stars), and Tr 14, Tr 16 (massive stellar clusters). The image also clearly shows
apparently extended emission toward the east-west direction. In a color image (e.g.
Figure 1 of Ref. 2), XMM-Newton Image Gallery∗)) the emission is softer between
Tr 14, WR 25 and η Car.
We analyzed the Suzaku data of the core and the eastern side (named Car-D1)
of the Carina nebula taken on 2005 Aug. 29 and 2006 June 5. The XIS FOVs
of these observations are shown in Figure 1 with dotted lines. To investigate the
color variation in detail, we divided the core region into two and thus extracted
three spectra from two Suzaku observations (core-north, core-south and Car-D1).
The background was reproduced with the night earth data. The spectra showed
strong emission between 0.3 and 2 keV, which is probably dominated by soft diffuse
emission associated with the Carina nebula, while the spectra above 2 keV may be
explained with CXB, Galactic Ridge X-ray Emission, X-ray point sources resolved
with Chandra and unresolved pre-main-sequence stars.
Figure 2 shows an overlay of the BI spectra between 0.3–2 keV. The left panel
compares spectra of the core-north region with the core-south region. A strong differ-
ence is seen between 0.7 keV and 1.2 keV, which apparently is the source of the two
colors of diffuse emission. The band in which the difference is found is dominated by
emission lines from the iron L-shell complex. Additionally, the core-south spectrum
shows a stronger Si line. The Car-D1 spectrum shows similar intensity in the Si and
Fe lines to the core-south spectrum (right panel of Figure 2) while it shows relatively
strong magnesium and oxygen lines. All these spectra look similar except for these
emission lines. This suggests that the differences represent an elemental abundance
variation, and not a temperature difference.
This is supported by spectral fits of the individual spectra. All three spectra
between 0.3−2 keV were reproduced by an absorbed 2T thin-thermal plasma models
although the best-fit models are not formally acceptable. The plasma tempera-
tures of all three regions are ∼0.2 and ∼0.6 keV, and their column densities are
∼3×1021 cm−2, which is consistent with extinction toward the Carina nebula.3) The
abundances of some elements show a factor of 2−4 variations: the core-north region
has a factor of 2 lower silicon abundance and a factor of 4 lower iron abundance
∗) http://xmm.esac.esa.int/external/xmm science/gallery/public
Diffuse X-rays from the Carina nebula 3
Fig. 1. Mosaic image (∼90′×60′) of the Carina nebula between 0.4−7 keV created from 32 XMM-
Newton observations. The image is created with the ESAS package, divided by the exposure
map and smoothed with the adaptive smoothing technique. The dotted lines show the XIS
FOVs of the Suzaku observations of η Car (right) and the Car-D1 field (left). The solid lines
show source extraction regions for the spectral analysis.
than the core-south region, while the Car-D1 region has a factor of 2 higher oxygen
and magnesium abundances. On the other hand, spectral fits of the core region
with higher sensitivity around 0.5 keV gave small upper-limits (.0.02 solar) of the
nitrogen abundance.
§3. Origin of the Diffuse Plasma
The N/O abundance ratio inferred from the spectral fits is .0.4, over 20 times
less than around η Car. The abundance distribution is totally contrary to that
expected from stellar winds from evolved massive stars, unless the winds somehow
heat the interstellar matter without enriching it, thus leaving the X-ray plasma with
abundances typical of interstellar matter. At the same time, the X-ray luminosity
of the Carina Nebula is about two orders of magnitude higher than that of other
Galactic star forming regions, but the number of early O stars is only an order of
magnitude higher (see Table 4 in Ref. 6). These results suggest an additional energy
source is needed to power the X-ray emission in the Carina Nebula.
An obvious possibility is one or more core-collapse supernovae (i.e. Type Ib,c
or II), mentioned as a possibility by Ref. 6). The regions vary strongly in oxygen,
magnesium, silicon, and iron abundances. These elements are products of core-
collapse supernovae, and young SNRs such as Cas A and Vela show strong abundance
4 K. Hamaguchi et al.
Fe L complex Fe L complex
Fig. 2. Comparison of the XIS1 spectra between the fields – left: the core-north region (black) and
the core-south region (grey), right: the Car-D1 field (black) and the core-south region (grey).
The above labels demonstrate energies of emission lines detected (black) or concerned (grey)
with this result. Emission lines with the solid lines showed variation in their line intensity. Low
count rates of the Car-D1 spectrum below 1 keV is caused by degradation of soft response by
progressive contamination on the XIS.
variation from location to location. The total energy content in the hot gas of
∼2×1050 ergs is a modest fraction of the ∼1051 ergs of kinetic energy produced by a
canonical supernova, while assuming an iron abundance of 0.30 solar, the total iron
mass in the diffuse gas requires at least 3-5 supernovae.
Acknowledgements
K. H. is financially supported by a US Chandra grant No. GO3-4008A and US
Suzaku grant.
References
1) N. R. Evans, F. D. Seward, M. I. Krauss, T. Isobe, J. Nichols, E. M. Schlegel, and S. J.
Wolk, Astrophysical Journal 2003 (589), 509
2) K. Hamaguchi, R. Petre, H. Matsumoto, M. Tsujimoto, S. S. Holt, Y. Ezoe, H. Ozawa,
Y. Tsuboi, Y. Soong, S. Kitamoto, A. Sekiguchi, and M. Kokubun. Publication of Astro-
nomical Society of Japan 2007 (59), 151
3) M. A. Leutenegger, S. M. Kahn, and G. Ramsay. Astrophysical Journal 2003 (585), 1015
4) F. D. Seward and T. Chlebowski. Astrophysical Journal 1982 (256), 530
5) L. K. Townsley. Proceeding of the STScI May Symposium, ”Massive Stars: From Pop III
and GRBs to the Milky Way, 2006, (astro–ph/0608173)
6) L. K. Townsley, E. D. Feigelson, T. Montmerle, P. S. Broos, Y.-H. Chu, and G. P. Garmire.
Astrophysical Journal 2003 (593), 874
http://arxiv.org/abs/astro--ph/0608173
Extended X-ray Emission from the Star Forming Region
Suzaku and XMM-Newton Observations of the Carina Nebula
Origin of the Diffuse Plasma
|
0704.0349 | The Colin de Verdi\`ere number and graphs of polytopes | The Colin de Verdière number
and graphs of polytopes
Ivan Izmestiev ∗
Institut für Mathematik
Technische Universität Berlin
Str. des 17. Juni 136
10623 Berlin, Germany
[email protected]
July 25, 2008
Abstract
The Colin de Verdière number µ(G) of a graph G is the maximum
corank of a Colin de Verdière matrix for G (that is, of a Schrödinger
operator on G with a single negative eigenvalue). In 2001, Lovász gave
a construction that associated to every convex 3-polytope a Colin de
Verdière matrix of corank 3 for its 1-skeleton.
We generalize the Lovász construction to higher dimensions by in-
terpreting it as minus the Hessian matrix of the volume of the polar
dual. As a corollary, µ(G) ≥ d if G is the 1-skeleton of a convex
d-polytope.
Determination of the signature of the Hessian of the volume is based
on the second Minkowski inequality for mixed volumes and on Bol’s
condition for equality.
1 Introduction
1.1 The Colin de Verdière number
At the end of 80’s, Yves Colin de Verdière introduced a graph parameter
µ(G) based on spectral properties of certain matrices associated with the
graph G.
Definition 1.1 Let G be a graph with n vertices. A Colin de Verdière
matrix for G is a symmetric n × n matrix M = (Mij) with the following
properties.
Research for this article was supported by the DFG Research Unit 565 “Polyhedral
Surfaces”.
http://arxiv.org/abs/0704.0349v3
(M1) M is a Schrödinger operator on G, that is
< 0, if ij is an edge of G;
= 0, if ij is not an edge of G and i 6= j.
(M2) M has exactly one negative eigenvalue, and this eigenvalue is simple.
(M3) If X is a symmetric n × n matrix such that MX = 0 and Xij = 0
whenever i = j or ij is an edge of G, then X = 0.
The set of all Colin de Verdière matrices for graph G is denoted by MG.
The Colin de Verdière number µ(G) is defined as the maximum corank of
matrices from MG:
µ(G) := max
dimkerM.
A Colin de Verdière matrix of maximum corank is called optimal.
Basically, the Colin de Verdière number is the maximum multiplicity of
the second least eigenvalue λ2 of a discrete Schrödinger operator M satis-
fying a certain stability assumption (M3). By replacing M with M − λ2Id,
we can make the second eigenvalue zero (M2), so that multiplicity be-
comes corank. Definition 1.1 was motivated by the study of Schrödinger
and Laplace operators associated with degenerating families of Riemannian
metrics on surfaces.
The parameter µ(G) turned out to be interesting on its own. In partic-
ular, it posesses the minor monotonicity property: if a graph H is a minor
of G, then µ(H) ≤ µ(G). By the Robertson-Seymour theorem this implies
that graphs with µ(G) ≤ n can be characterized by a finite set of forbidden
minors. For n up to four such characterizations are known and allow nice
topological reformulations: e. g. µ(G) ≤ 3 iff G is planar (that is doesn’t
have K5 or K3,3 as minors), and µ(G) ≤ 4 iff G is linklessly embeddable in
3 (that is doesn’t have any graph of the Petersen family as a minor). An
overview of results and open problems on the Colin de Verdière number can
be found in [4], [14], and [5]. The book [4] deals also with other spectral
invariants arising from discrete Schrödinger and Laplace operators.
1.2 Nullspace representations and Steinitz representations
Let M be a Colin de Verdière matrix for graph G with dimkerM = d.
Choose a basis (u1, . . . , ud) for kerM ⊂ R
n, fix a coordinate system in Rn,
and read off the coordinates of (uα):
(u1, . . . , ud) = (v1, . . . , vn)
The map that associates to every vertex i of G the vector vi ∈ R
d is called
a nullspace representation of the graph G.
Nullspace representations were studied in [11]. In a subsequent paper
[10] Lovász showed that, for a 3-connected planar G, the nullspace repre-
sentation with properly scaled vectors (vi) realizes G as the skeleton of a
convex 3-polytope. Lovász provided also an inverse construction that as-
sociated to every convex 3-polytope with 1-skeleton G a Colin de Verdère
matrix of corank 3. The proof that the constructed matrix had an appro-
priate signature was indirect, and a more geometric approach was desirable.
1.3 Hessian matrix of the volume as a Colin de Verdière
matrix
In this paper we relate the Lovász construction (that of a matrix from a poly-
tope) to the mixed volumes. Our approach allows a straightforward gener-
alization to higher dimensions. That is, we associate to every d-dimensional
convex polytope with 1-skeleton G a Colin de Verdière matrix for G of
corank d.
As a consequence, the graph of a convex d-dimensional polytope has
Colin de Verdière number at least d. This result is not really new, since it
follows from the minor monotonicity of µ, from the fact that the graph of a
d-polytope has Kd+1 as a minor [8], and from µ(Kd+1) = d.
Our result is based on the following observation. Take a convex d-
polytope P and deform it by shifting every facet parallelly to itself. Then
the Hessian matrix of the volume of P , where partial derivatives are taken
with respect to the distances of the shifts, has corank d and exactly one
positive eigenvalue. Besides, the mixed partial derivative
∂2vol(P )
∂xi∂xj
is positive
if the ith and the jth facets are adjacent, and vanishes otherwise. Thus
the negative of the Hessian matrix satisfies conditions (M1) and (M2) from
Definition 1.1. The condition (M3) follows quite easily, too.
The signature of the Hessian of the volume is encoded in the second
Minkowski inequality for mixed volumes together with Bol’s characteriza-
tion of the case of equality. For simple polytopes, the determination of the
signature of the Hessian is an essential part in the proof of the Alexandrov-
Fenchel inequality.
1.4 Plan of the paper
In Section 2.1 we recall the Lovász construction of a Colin de Verdière matrix
for the skeleton of a convex 3-polytope Q.
After inroducing some terminology and notation in Section 2.2, we show
in Section 2.3 that the Lovász matrix is minus the Hessian matrix of the
volume of the polar dual polytope Q∗.
In Section 2.4, dealing with 3-polytopes, we point out an interesting
identity (first found and used elsewhere [2]) between the Hessian matrix of
vol(Q∗) and the Hessian matrix of another geometric quantity associated
with Q. This gives another interpretation of the Lovász matrix M and
relates the equality dimkerM = 3 with the infinitesimal rigidity of the
polytope Q.
In Section 3.1 we discuss the (im)possibility of inverting the construction,
that is of finding a convex polytope whose Hessian matrix of the volume
equals to a given Colin de Verdière matrix.
In Section 3.2 we give an estimate of the negative eigenvalue (and thus
of the spectral gap) for the Hessian matrices of the volume.
Finally, in the Appendix we derive the signature of the Hessian from the
second Minkowski inequality and Bol’s condition. Although this seems to
be a folklore knowledge in narrow circles, we failed to find a written account
on this subject.
1.5 Acknowledgements
I am grateful to the organizers of the 2006 Oberwolfach conference “Discrete
Differential Geometry”, where the idea of this paper was born. I also thank
Ronald Wotzlaw for pointing me out a mistake in a preliminary version.
2 From a convex polytope to a Colin de Verdière
matrix
2.1 Lovász construction
Let us recall the Lovász construction of an optimal Colin de Verdière matrix
associated with a polytopal representation of a graph in R3.
Let Q ⊂ R3 be a convex polytope containing the coordinate origin in
its interior. Let G be the 1-skeleton of Q. We denote the vertices of G
by i, j, . . ., and the corresponding vertices of Q by vi, vj , . . .. Let Q
∗ be the
polar dual of Q. The vertices of Q∗ are denoted by wf , wg, . . ., where f, g, . . .
are faces of Q.
For ij ∈ G, consider the edge vivj of Q and the dual edge wfwg of Q
see Figure 1. It is easy to show that the vector wf − wg is orthogonal to
both vectors vi and vj, hence parallel to their cross product vi × vj. Thus
we have
wf − wg = Mij(vi × vj), (1)
with Mij < 0 (we agree to choose the labeling of wf and wg so that we get
the correct sign).
Further, consider the vector
v′i =
Mijvj,
vi × vj
Figure 1: To the definition of the matrix M .
where the sum extends over all vertices of G adjacent to i. From (1) it is
easy to see that vi × v
i = 0. Thus there exists a real number Mii such that
v′i = −Miivi. (2)
Putting Mij = 0 for distinct non-adjacent vertices i and j of G, we complete
the construction of the matrix M .
Theorem 2.1 (Lovász, [10]) The matrix M is a Colin de Verdière matrix
for the graph G.
The equation (2) can be rewritten as
Mijvj = 0. (3)
Thus M has corank at least 3. Since µ(G) ≤ 3 for planar graphs, M is an
optimal Colin de Verdière matrix for G.
The proof of Theorem 2.1 goes through a deformation argument, using
the fact that the space of convex 3-polytopes with a given graph is connected.
2.2 Polytopes with a given set of normals
Here we fix some terminology and notation needed in the subsequent sec-
tions.
All polytopes in this paper are assumed to be convex. A facet of a
d-dimensional polytope is a (d− 1)-dimensional face of it.
We will study families of polytopes with fixed facet normals. Let v1, . . . , vn
be vectors in Rd such that the coordinate origin lies in the interior of their
convex hull. Consider a d× n matrix formed by row vectors v⊤i :
V = (v1, . . . , vn)
Definition 2.2 Denote by P(V ) the set of all convex polytopes with the
outer facet normals v1, . . . , vn.
Every polytope in P(V ) is the solution set of a system of linear inequal-
ities:
P (x) = {p ∈ Rd |V p ≤ x},
where x = (xi)
i=1 ∈ R
n. Denote by Fi(x) the facet of P (x) with the outer
normal vi. We have
Fi(x) = {p ∈ P (x) | v
i p = xi}.
The numbers xi are called the support parameters of the polytope P (x).
The map P (x) 7→ x embeds P(V ) into Rn as an open convex subset. The
support parameter xi is proportional to the signed distance from 0 to the
affine hull of the facet Fi(x):
xi = ‖vi‖ · hi.
By vold we denote the volume of a d-dimensional polytope. We use the
subscript because both vold(P ) and vold−1(Fi) will occur in our formulas.
We omit the subscript at vol, when it seems reasonable to do so.
2.3 Interpreting and generalizing the Lovász construction
By definition of the polar dual, we have
Q∗ = {p ∈ R3 | v⊤i p ≤ 1 for all i}.
Thus Q∗ can be viewed as an element of the set P(V ) of polytopes with
facet normals (vi)i∈G. In terms of Section 2.2, Q
∗ = P (1, . . . , 1). Let’s vary
the support parameters of Q∗ and look how does this change its volume.
Lemma 2.3 Let M be the matrix constructed in Section 2.1. Then we have
Mij = −
∂2vol(P (x))
∂xi∂xj
x=(1,...,1)
where P (x) is as in Section 2.2.
Proof . Let Fi(x) be the facet of P (x) with the normal vi. It is not hard to
show that
∂vol3(P (x))
vol2(Fi(x))
Further, for i 6= j we have
∂vol2(Fi(x))
vol1(Fij(x))
‖vj‖ sin θij
∆vol1(Fj)
∆vol2(P )
Figure 2: Partial derivatives of the volume with respect to the support
parameters.
if faces Fi(x) and Fj(x) are adjacent; otherwise this derivative is zero. Here
Fij(x) is the common edge of Fi(x) and Fj(x), and θij is the angle between
the vectors vi and vj (i. e. the outer dihedral angle at the edge Fij). The
equations are illustrated in Figure 2 in one dimension lower and for ‖vi‖ = 1.
Thus at x = (1, . . . , 1) we have
∂2vol(P (x))
∂xi∂xj
vol1(Fij(x))
‖vi‖‖vj‖ sin θij
‖wf − wg‖
‖vi × vj‖
= −Mij (4)
for all i 6= j.
To deal with the case i = j, differentiate the well-known identity
vol2(Fj(x))
with respect to xi. This gives
∂2vol(P (x))
j 6=i
∂2vol(P (x))
∂xi∂xj
vj = 0. (5)
In view of (3) and (4), we have
∂2vol(P (x))
|x=(1,...,1) = −Mii. �
Lemma 2.3 suggests the following generalization of the Lovász construc-
tion.
Theorem 2.4 Let
P (x0) = {p ∈ Rn | v⊤i p ≤ x
i for all i}
be a convex polytope with outer facet normals vi and support parameters
x0i , i = 1, . . . , n. Let G be the dual 1-skeleton of P (x
0). Then the matrix M
defined by
Mij = −
∂2vol(P (x))
∂xi∂xj
is a Colin de Verdiére matrix for the graph G.
The corank of M is equal to d. In particular, µ(G) ≥ d for every graph
G that can be realized as the 1-skeleton of a d-dimensional polytope.
Proof . Similarly to Lemma 2.3, for adjacent facets Fi and Fj we have
∂2vold(P (x))
∂xi∂xj
vold−2(Fij(x))
‖vi‖‖vj‖ sin θij
where Fij is their common (d− 2)-face, and θij is the angle between vi and
vj. For non-adjacent Fi and Fj this derivative is zero. Therefore matrix M
satisfies property (M1) from Definition 1.1.
The proof of property (M2) is the most interesting part of the theo-
rem. The signature of the Hessian of the volume is encoded in the second
Minkowski inequality for mixed volumes enhanced by Bol’s condition for
equality.
Theorem A.10 in Section A states in particular that the matrix M has
corank d. The kernel of M is easy to identify: due to the equation (5) it
consists of the vectors ξ ∈ Rn such that ξi = v
i p for some vector p ∈ R
Assuming this description of kerM , let us prove that the matrix M
satisfies property (M3). If MX = 0, then there are vectors p1, . . . , pn ∈ R
such that Xij = v
i pj for all i, j. Fix j. Then by assumption on X we have
pj ⊥ vj and pj ⊥ vi for all ij ∈ G. But the normal vj to the face Fj and the
normals to the neighboring faces span the space Rd. Thus we have pj = 0
for all j, which implies X = 0.
As for the last sentence of the theorem, if G is the dual 1-skeleton of a
convex polytope P , then G is the skeleton of the polar (P − p)∗, where p is
any interior point of P . �
2.4 Case d = 3 and infinitesimal rigidity of convex polytopes
In the case d = 3 there is another interpretation of the matrix M . As in
Section 2.1, let Q be a convex polytope that has skeleton G and contains
0 in the interior. Triangulate the faces of Q by diagonals and cut Q into
pyramids with apices at 0 and triangles of the triangulation as bases. De-
note by ri the length of the edge that joins 0 to the vertex vi of Q. Now
deform the pyramids by changing the lengths ri and leaving the lengths of
boundary edges constant. During such deformation, the dihedral angles of
the pyramids change, and the total angle ωi around the i-th edge can be-
come different from 2π. By computing the derivatives of ωi explicitly, we
obtain ([2], Theorem 3.11)
vol1(Fij)
sin θij
= ‖vi‖‖vj‖ ·Mij, (7)
where we use the notations from Section 2.3. If we change the variables xi to
hi = ‖vi‖ · xi, so that hi is the distance of 0 from aff (Fi), then the equation
(7) takes a particularly nice form
∂vol2(Fi)
By (7), the matrix (∂ωi
) is obtained from the matrix M by multiplying
the i-th row and the i-th column with ‖vi‖, for all i. This implies
Corollary 2.5 The matrix (∂ωi
) is an optimal Colin de Verdière matrix for
graph G.
The fact that the matrix (∂ωi
) has corank 3 is equivalent to the infinitesi-
mal rigidity of the polytope Q. Indeed, every infinitesimal deformation (dri)
such that dωi = 0 for all i gives rise to an infinitesimal isometric deformation
of Q. The resulting deformation is trivial iff it is produced by moving the
apex 0 inside Q.
Another interesting fact is that the matrix (∂ωi
) is the Hessian matrix
of a geometric quantity related to the polytope Q (deformed by varying ri).
Namely, put
S(r) =
riκi +
ℓijθij,
where κi = 2π − ωi is the “curvature” along the i-th radial edge, and ℓij =
vol1(Fij) is the length of the edge vivj. Then the Schläfli formula implies
= κi.
Hence
∂2S(Q)
∂ri∂rj
∂2vol(Q∗)
∂hi∂hj
and both matrices are equal to the negative of the Lovász matrix M , up to
scaling the rows and columns by ‖vi‖.
3 Concluding remarks
3.1 What fails in the inverse construction
Let M be a Colin de Verdère matrix for the graph G. Is there a convex
polytope P such that M arises from P as a result of the construction de-
scribed in Section 2.3? Of course, in general the answer is no, because G
must be the dual skeleton of P , and P must have dimension d = dimkerM .
In particular, all vertices of G must have degrees at least d. But, due to the
minor monotonicity of µ, there exist trivalent graphs with µ(G) arbitrarily
large.
Nevertheless, it is worth looking at what fails when we try to reconstruct
the polytope P from matrix M .
Let u1, . . . , ud ∈ R
n be a basis of kerM . Let v⊤i be the i-th row in the
matrix (u1, . . . , ud). Then we have
Mijvj = 0 (8)
for all i. Therefore, the vectors v1, . . . , vn ∈ R
d are good candidates for the
outer normals to the faces of the polytope P . At this point we can already
fail, if the following assumptions aren’t fulfilled:
1. vi 6= 0 for all i, and vi 6= vj for all i 6= j;
2. for every i, the projections vij of vj on v
i for ij ∈ G satisfy the
previous assumption and span v⊥i .
We proceed assuming that these conditions hold. Codimension 2 faces
Fij of P must be in 1-to-1 correspondence with the edges of G, and their
volumes are determined by the matrix M :
vold−2(Fij) = Aij := −Mij‖vi‖‖vj‖ sin θij,
where θij is the angle between vi and vj .
Lemma 3.1 For every i, there exists a convex (d−1)-dimensional polytope
Fi ⊂ v
i with outer facet normals vij and facet volumes Aij , ij ∈ G.
Proof . By projecting the equation (8) on v⊥i , we obtain
Mij · vij = 0. (9)
Due to ‖vij‖ = ‖vj‖ sin θij, it follows that
Aij ·
‖vij‖
By Minkowski’s theorem [13, Section 7.1], this implies the existence of a
polytope Fi as stated in the lemma. �
The polytopes Fi in Lemma 3.1 should become facets of the polytope
P . But here is the second point where the reconstruction can fail: the j-th
facet Fij of Fi might be different from the i-th facet Fji of Fj ; the only thing
we know is vold−2(Fij) = Aij = vold−2(Fji).
In the case d = 3, however, this suffices: Fi are convex polygons and
fit together along their edges to form a polytope P . Conditions 1. and 2.
above hold if we assume that G is a 3-connected planar graph [11]. Thus for
3-connected planar graphs every Colin de Verdière matrix corresponds to a
polytope. This is one of the results of [10].
The following example shows that even for highly connected graphs the
number µ(G) can be bigger than the maximum dimension of a polytope
with 1-skeleton G.
Example Let Gn = K2,2,...,2 be the multipartite graph on 2n vertices (the
graph of an n-dimensional cross-polytope). By [9], µ(Gn) = 2n − 3 for
n ≥ 3. For n = 3, 4 the graph Gn can also be represented as the skeleton of
a (2n − 3)-dimensional convex polytope: for n = 3 this is the octahedron,
for n = 4 the join of two convex quadrilaterals in general position in R5.
For n ≥ 5, however, there is no (2n − 3)-dimensional convex polytope with
skeleton Gn. Indeed, by studying the Gale diagram [16, Lecture 6] of a
d-polytope with d + 3 vertices, one can show that the complement to the
graph of such polytope cannot have more than 4 edges.
Note that the equation (9) is reminiscent of the definition of a (d − 2)-
weight in [12].
3.2 Negative eigenvalue
Theorem 3.2 Let λ1 be the negative eigenvalue of the matrix (6). Then
the following inequality holds:
λ1 ≤ −d(d− 1) ·
vold(P (x
‖x0‖2
The equality takes place iff
x0i = c ·
vold−1(Fi(x
for all i and some constant c.
Proof . By induction on d, it is easy to show that the function vold(P (x)) is a
degree d homogeneous polynomial in x as long as the combinatorics of P (x)
does not change. For different combinatorics, the polynomials have different
coefficients. However, since vold(P (x)) is twice differentiable, we can apply
Euler’s homogeneous function theorem twice at the point x0, independently
on how generic the combinatorics of P (x0) is. This yields
(x0)⊤Mx0 = −d(d− 1) · vold(P (x
Since λ1 = min‖ξ‖=1 ξ
⊤Mξ, the inequality follows.
Since λ1 is the unique negative eigenvalue of M , the inequality turns
into equality iff Mx0 = λx0 for some λ. We have
j = −
∂vold−1(Fi(x
x0j = −(d− 1) ·
vold−1(Fi(x
Thus Mx0 = λx0 is equivalent to x0i = c ·
vold−1(Fi(x
, and the theorem is
proved. �
The number λ2 − λ1 is called the spectral gap. In our case λ2 = 0 by
definition. Thus Theorem 3.2 provides an estimate on the spectral gap of
the matrix M .
Usually, one seeks to make the spectral gap as large as possible, but in
order this to make sense for Colin de Verdère matices, one has to choose a
matrix norm, [4, Chapter 5.7]. The norm of the matrix (6) is a function of
its coefficients, which have a geometric meaning. Thus, as soon as the choice
of a matrix norm is made, one can try to solve the problem of the spectral
gap by geometric means (at least for 3-connected planar graphs, for which
every optimal Colin de Verdière matrix can be realized through a polytope).
A The second Minkowski inequality for mixed vol-
umes and the signature of the matrix
∂2vol
∂xi∂xj
The goal of this appendix is to prove Theorem A.10 that describes the sig-
nature of the matrix (6). The theorem is derived from the second Minkowski
inequality for mixed volumes and Bol’s condition for equality.
The relation between the theory of mixed volumes and infinitesimal rigid-
ity (as we know, the rank of matrix (6) accounts for the infinitesimal rigidity
of the dual polytope, see Section 2.4) was noticed long ago [1, 15]. In the
decades thereafter this phenomenon seemed to be forgotten. Quite recently,
Carl Lee and Paul Filliman [7] discovered it again.
A.1 The second Minkowski inequality and Bol’s condition
Definition A.1 Let P,Q ⊂ Rd be convex bodies. A mixed volume of P and
Q is a coefficient in the expansion
vol(λP + µQ) =
vol(Q, . . . , Q
︸ ︷︷ ︸
, P, . . . , P
︸ ︷︷ ︸
)λd−kµk (10)
with λ, µ > 0, where A + B for A,B ⊂ Rd denotes the Minkowski sum. In
particular,
vol(P, . . . , P ) = vol(P ).
In a similar way one defines the mixed volume of more than two convex
bodies. It turns out that the mixed volume is polylinear with respect to the
Minkowski addition and multiplication with positive scalars. A proof that
the expansion (10) takes place and more information on mixed volumes can
be found in [6, 13].
Theorem A.2 Let P,Q ⊂ Rd be convex bodies. Then the following holds:
1. (The second Minkowski inequality)
vol(Q,P, . . . , P )2 ≥ vol(P ) · vol(Q,Q,P, . . . , P ). (11)
2. (Bol’s condition) Assume that dimQ = d. Then equality holds in (11)
if and only if either dimP < d − 1 or P is homothetic to a (d − 2)-
tangential body of Q.
For a proof see [13, Theorem 6.2.1, Theorem 6.6.18]. Bol’s condition was
conjectured by Minkowski but proved only decades later by Bol, [3].
Definition A.3 If P ⊂ Q ⊂ Rd are d-dimensional convex polytopes, then
Q is called a p-tangential body of P iff P has a non-empty intersection with
every face of Q of dimension at least p.
A.2 Mixed volumes as derivatives of the volume
By substituting in (10) λ = 1 and µ = t, we obtain
vol(P + tQ) = vol(P ) + tdvol(Q,P, . . . , P )
d(d− 1)
vol(Q,Q,P, . . . , P ) + · · · (12)
for all t > 0, which can be seen as the Taylor expansion of vol. We will look
at it in the case when P and Q are polytopes with the same sets of facet
normals.
The space P(V ) of all polytopes with outer facet normals v1, . . . , vn is
defined in Section 2.2. We want to study the partial derivatives of the volume
of P (x) ∈ P(V ) with respect to the support parameters x. For brevity, let’s
use the notation
vol(x) := vol(P (x)).
Similarly, the mixed volume of polytopes from P(V ) will be written as a
function of the support parameters:
vol(x1, . . . , xd) := vol(P (x1), . . . , P (xd)).
Now we would like to compute vol(x+ ty) with the help of (12). This is
not as straightforward as it seems, because the support parameters behave
not quite linearly under the Minkowski addition. We have P (ty) = tP (y)
for t > 0. Also we have P (x) + P (y) ⊂ P (x + y), but the equality doesn’t
always hold. To describe the cases in which we do have the equality, we
need a new definition.
Definition A.4 The normal cone N(F,P ) of the face F of a polytope P ⊂
d is the set of vectors w ∈ Rd such that
(w⊤x) = max
(w⊤x).
The normal fan N(P ) is the decomposition of Rd into the normal cones of
the faces of P . If the normal fan N(Q) subdivides the normal fan N(P ),
then we write N(Q) > N(P ).
Note that the normal fan of a polytope P ∈ P(V ) has the rays R+vi as
1-dimensional cones. The higher-dimensional cones of the normal fan deter-
mine the combinatorics of P . Therefore polytopes with equal normal fans
are sometimes called strongly isomorphic.
We denote the normal fans of the polytopes from P(V ) by N(x) :=
N(P (x)). The following lemma is classical.
Lemma A.5 If N(y) > N(x), then P (x) + P (y) = P (x+ y).
Now we are ready to prove
Lemma A.6 Let y ∈ P(V ) be such that N(y) > N(x). Then
∇yvol(x) = d · vol(y, x, . . . , x),
∇2yvol(x) = d(d− 1) · vol(y, y, x, . . . , x),
where ∇y denotes the directional derivative along y.
Proof . Due to Lemma A.5 we have P (x + ty) = P (x) + tP (y). By substi-
tuting P = P (x) and Q = P (y) in (12), we obtain
vol(x+ ty) = vol(x) + tdvol(y, x, . . . , x)
d(d− 1)
vol(y, y, x, . . . , x) + · · · ,
which implies the lemma. �
Remark. For polytopes with the same normal fan (“strongly isomorphic
polytopes”), there is the following description of mixed volumes. Denote
P∆(V ) = {P (x) ∈ P(V ) |N(x) = ∆}.
By induction on d, it is easy to show that there exists a homogeneous poly-
nomial V∆ of degree d in n variables such that
vol(P (x)) = V∆(x),
for all x ∈ P∆(U). If we use the same symbol V∆ to denote the associated
symmetric polylinear form, then we have
vol(P (x(1)), . . . , P (x(d))) = V∆(x
(1), . . . , x(d))
for all x(1), . . . , x(d) ∈ P∆(V ).
A.3 From the second Minkowski inequality to the signature
of the Hessian of the volume
By geometric arguments similar to those in the proof of Lemma 2.3, the
function vol is twice continuously differentiable on P(V ). Therefore the
following definition makes sense.
Definition A.7 Let x ∈ P(V ). Define a symmetric bilinear form Φ on Rn
Φ(ξ, η) = ∇η∇ξvol(x).
Let y ∈ P(V ) be such that N(y) > N(x). By combining Euler’s homo-
geneous function theorem and Lemma A.6, we obtain
Φ(x, x) = d(d− 1) vol(x, . . . , x),
Φ(x, y) = d(d− 1) vol(y, x, . . . , x),
Φ(y, y) = d(d− 1) vol(y, y, x, . . . , x).
Lemma A.8 Let L ⊂ Rn be a 2-dimensional vector subspace such that
x ∈ L. Then the restriction of the form Φ to L has signature (+,−) or
(+, 0).
Proof . Let y ∈ P(V ) be such that N(y) > N(x). The second Minkowski
inequality (11) applied to P = P (x) and Q = P (y) can be rewritten as
Φ(x, x) Φ(x, y)
Φ(x, y) Φ(y, y)
≤ 0. (13)
Since, moreover, Φ(x, x) = d(d−1) vol(P ) > 0, it follows that the restriction
of Φ to span {x, y} has signature (+, 0) or (+,−).
It remains to show that every 2-subspace L ∋ x can be represented as
span {x, y} with N(y) > N(x). This is true since x is an interior point of
the set {y ∈ P(V ) |N(y) > N(x)}. (When we perturb x, we can create new
faces, but cannot destroy old ones.) �
Lemma A.9 The form Φ has corank d.
Proof . Let us exhibit a d-dimensional subspace of ker Φ. Associate with
every point p ∈ Rd a vector p ∈ Rn with coordinates
pi = 〈vi, p〉.
The polytope P (x+p) is the translate of P (x) by p. Therefore the directional
derivative ∇pvol(x) vanishes for all x, which implies Φ(p, η) = 0 for all η.
Thus we have
p ∈ kerΦ
for all p ∈ Rd.
Let ξ ∈ ker Φ. We need to show that ξ = p for some p ∈ Rd. Denote
the span of x and ξ by L. Then, by Lemma A.8, the restriction Φ|L has
signature (+, 0) and hence
Rξ = L ∩ ker Φ. (14)
Choose y ∈ L such that N(y) > N(x), and x and y are linearly independent.
Then the degeneracy of Φ|L means that we have an equality in (13) and
thus also in the Minkowski inequality for P = P (x) and Q = P (y). By
Bol’s condition, see Theorem A.2, this happens if and only if the polytope
P (x) is homothetic to a (d − 2)-tangential body of the polytope P (y). By
studying Definition A.3, we see that in P(V ) it is equivalent to P (x) being
homothetic to P (y). If P (x) is homothetic to P (y), then x = λy + p for
some p ∈ Rd, thus p ∈ L. Since p ∈ kerΦ, it follows that
Rp = L ∩ ker Φ.
By comparing this to (14), we conclude that ξ = µp = µp for some µ ∈ R.
Thus the kernel of Φ is confined to the vectors of the form p. �
Theorem A.10 The form Φ has corank d and exactly one positive eigen-
value, which is simple.
Proof . The corank of Φ is computed in Lemma A.9.
The form Φ has at least one positive eigenvector since Φ(x, x) > 0.
Assume that it has more than one. Then there exists a 2-subspace of Rn
on which Φ is positively definite. The subgroup of GL(Rn) that preserves Φ
acts transitively on the cone of positive directions. Thus there is a positive
2-subspace L that passes through x. This contradicts Lemma A.8. Theorem
is proved. �
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Introduction
The Colin de Verdière number
Nullspace representations and Steinitz representations
Hessian matrix of the volume as a Colin de Verdière matrix
Plan of the paper
Acknowledgements
From a convex polytope to a Colin de Verdière matrix
Lovász construction
Polytopes with a given set of normals
Interpreting and generalizing the Lovász construction
Case d=3 and infinitesimal rigidity of convex polytopes
Concluding remarks
What fails in the inverse construction
Negative eigenvalue
The second Minkowski inequality for mixed volumes and the signature of the matrix (2 volxi xj)
The second Minkowski inequality and Bol's condition
Mixed volumes as derivatives of the volume
From the second Minkowski inequality to the signature of the Hessian of the volume
|
0704.0350 | Visible spectroscopic and photometric survey of Jupiter Trojans: final
results on dynamical families | Visible spectroscopic and photometric survey
of Jupiter Trojans: final results on dynamical
families. ∗
Fornasier S.1,2, Dotto E.3, Hainaut O.4, Marzari F.5,
Boehnhardt H.6, De Luise F.3, Barucci M.A.2
October 22, 2018
1 University of Paris 7, France
2 LESIA – Paris Observatory, France.
3 INAF – Osservatorio Astronomico di Roma, Italy;
4 European Southern Observatory, Chile;
5 Dipartimento di Fisica, Università di Padova, Italy;
6 Max-Planck Institute for Solar System Research, Katlenburg-Lindau, Germany
Submitted to Icarus: December 2006
e-mail: [email protected]
fax: +33145077144, phone: +33145077746
Running head: Investigation of Dynamical Families of Jupiter Trojans
∗Based on observations carried out at the European Southern Observatory (ESO), La
Silla, Chile, ESO proposals 71.C-0650, 73.C-0622, 74.C-0577
http://arxiv.org/abs/0704.0350v1
Send correspondence to:
Sonia Fornasier
LESIA-Observatoire de Paris
Batiment 17
5, Place Jules Janssen
92195 Meudon Cedex
France
e-mail: [email protected]
fax: +33145077144
phone: +33145077746
Abstract
We present the results of a visible spectroscopic and photometric
survey of Jupiter Trojans belonging to different dynamical families.
The survey was carried out at the 3.5m New Technology Telescope
(NTT) of the European Southern Observatory (La Silla, Chile) in
April 2003, May 2004 and January 2005. We obtained data on 47
objects, 23 belonging to the L5 swarm and 24 to the L4 one. These
data together with those already published by Fornasier et al. (2004a)
and Dotto et al. (2006), acquired since November 2002, constitute a
total sample of visible spectra for 80 objects.
The survey allows us to investigate six families (Aneas, Anchises, Mis-
enus, Phereclos, Sarpedon, Panthoos) in the L5 cloud and four L4 fam-
ilies (Eurybates, Menelaus, 1986 WD and 1986 TS6). The sample that
we measured is dominated by D–type asteroids, with the exception of
the Eurybates family in the L4 swarm, where there is a dominance of
C– and P–type asteroids.
All the spectra that we obtained are featureless with the exception
of some Eurybates members, where a drop–off of the reflectance is
detected shortward of 5200 Å. Similar features are seen in main belt
C–type asteroids and commonly attributed to the intervalence charge
transfer transition in oxidized iron.
Our sample comprises fainter and smaller Trojans as compared to the
literature’s data and allows us to investigate the properties of objects
with estimated diameter smaller than 40–50 km. The analysis of the
spectral slopes and colors versus the estimated diameters shows that
the blue and red objects have indistinguishable size distribution, so
any relationship between size and spectral slopes has been found.
To fully investigate the Trojans population, we include in our anal-
ysis 62 spectra of Trojans available in literature, resulting in a total
sample of 142 objects. Although the mean spectral behavior of L4
and L5 Trojans is indistinguishable within the uncertainties, we find
that the L4 population is more heterogeneous and that it has a higher
abundance of bluish objects as compared to the L5 swarm.
Finally, we perform a statistical investigation of the Trojans’s spectra
property distributions as a function of their orbital and physical pa-
rameters, and in comparison with other classes of minor bodies in the
outer Solar System. Trojans at lower inclination appear significantly
bluer than those at higher inclination, but this effect is strongly driven
by the Eurybates family. The mean colors of the Trojans are similar
to those of short period comets and neutral Centaurs, but their color
distributions are different.
Keywords: Trojan Asteroids – Photometry – Spectroscopy – – Asteroids
families
1 Introduction
Jupiter Trojans are small bodies of the Solar System located in the Jupiter
Lagrangian points L4 and L5. Up to now more than 2000 Trojans have been
discovered, ∼ 1150 belonging to the L4 cloud and ∼ 950 to the L5 one.
The number of L4 Trojans with radius greater than 1 km is estimated to be
around 1.6 ×105 (Jewitt et al., 2000), comparable with the estimated main
belt population of similar size.
The debate about the origin of Jupiter Trojans and how they were trapped
in librating orbits around the Lagrangian points is still open to several possi-
bilities. Considering that Trojans have orbits stable over the age of the Solar
System (Levison et al, 1997, Marzari et al. 2003) their origin must date back
to the early phase of the solar system formation. Some authors (Marzari &
Scholl, 1998a,b; Marzari et al., 2002) suggested that they formed very close
to their present location and were trapped during the growth of Jupiter.
Morbidelli et al. (2005) suggested that Trojans formed in the Kuiper belt
and were subsequently captured in the Jupiter L4 and L5 Lagrangian points
during planetary migration, just after Jupiter and Saturn crossed their mu-
tual 1:2 resonances. In this scenario, Jupiter Troians would give important
clues on the composition and accretion of bodies in the outer regions of the
solar nebula.
Several theoretical studies conclude that Jupiter Trojan clouds are at
least as collisionally evolved as main belt asteroids (Shoemaker et al., 1989;
Binzel & Sauter, 1992; Marzari et al., 1997; Dell’Oro et al., 1998). This
result is supported by the identification of several dynamical families, both
in the L4 and L5 swarm (Shoemaker et al., 1989, Milani, 1993, Beaugé and
Roig, 2001).
Whatever the Trojan origin is, it is plausible to assume that they formed be-
yond the frost line and that they are primitive bodies, are possibly composed
of anhydrous silicates and organic compounds, and possibly still contain ices
in their interior. Several observations of Trojans in the near infrared region
(0.8-2.5 µm) have failed to clearly detect any absorption features indicative
of water ice (Barucci et al, 1994; Dumas et al, 1998; Emery & Brown, 2003,
2004; Dotto et al., 2006). Also in the visible range Trojan spectra appear
featureless (Jewitt & Luu, 1990; Fornasier et al., 2004a, Bendjoya et al.,
2004; Dotto et al., 2006). Up to now only 2 objects (1988 BY1 and 1870
Glaukos) show the possible presence of faint bands (Jewitt & Luu, 1990).
However, these bands are comparable to the peak to peak noise and are not
yet confirmed.
Recently, mineralogical features have been detected in emissivity spectra of
three Trojan asteroids measured by the Spitzer Space Telescope. These fea-
tures are interpreted as indicating the presence of fine-grained silicates on
the surfaces (Emery et al. 2006).
Several questions about Jupiter Trojans’ dynamical origin, physical prop-
erties, composition and link with other groups of minor bodies such as outer
main belt asteroids, cometary nuclei, Centaurs and KBOs are still open.
In order to shed some light on these questions, we have carried out a spectro-
scopic and photometric survey of Jupiter Trojans at the 3.5m New Technol-
ogy Telescope (NTT) of the European Southern Observatory (La Silla, Chile)
and at the 3.5m Telescopio Nazionale Galileo (TNG), La Palma, Spain. In
this paper we present new visible spectroscopic and photometric data, ob-
tained during 7 observing nights, carried out at ESO-NTT on April 2003,
May 2004, and January 2005, for a total of 47 objects belonging to the L5
(23 objects) and L4 (24 objects) swarms. Considering also the results already
published in Fornasier et al. (2004a) and Dotto et al. (2006), obtained in
the framework of the same project, we collected a total sample of 80 Jupiter
Trojan visible spectra, 47 belonging to the L5 clouds and 33 to the L4. This
is the largest homogeneous data set available up to now on these primitive
asteroids.
The principal aim of our survey was the investigation of Jupiter Trojans
belonging to different dynamical families. In fact, since dynamical families
are supposed to be formed from the collisional disruption of parent bodies,
the investigation of the surface properties of small and large family members
can help in understanding the nature of these dynamical groups and might
provide a glimpse of the interior structure of the larger primordial parent
bodies.
We also present an analysis of the visible spectral slopes for all the data in
our survey along with those available in the literature, for a total sample of
142 Trojans.
This enlarged sample allowed us to carry out a significant statistical investi-
gation of the Trojans’ spectral property distributions, as a function of their
orbital and physical parameters, and in comparison with other classes of mi-
nor bodies in the outer Solar System. We also discuss the spectral slope
distribution within the Trojan families.
2 Observations and data reduction
[HERE TABLE 1 AND 2]
The data were obtained in the visible range during 3 different observing
runs at ESO-NTT: 10 and 11 April 2003 for the spectroscopic and pho-
tometric investigation of 6 members of the 4035 1986 WD and 1 member
of 1986 TS6 families; 25 and 26 May 2004 for a spectroscopic survey of L4
Eurybates family; 17, 18, and 19 January 2005 for the spectroscopic and pho-
tometric investigation of 5 Anchises, 6 Misenus, 5 Panthoos, 2 Cloanthus, 2
Sarpedon and 3 Phereclos family members (L5 swarm).
We selected our targets from the list of Jupiter Trojan families provided by
Beaugé and Roig (2001 and P.E.Tr.A. Project at www.daf.on.br/ froig/petra/).
The authors have used a cluster-detection algorithm called Hierarchical Clus-
tering Method (HCM, e.g. Zappalà et al., 1990) to find asteroid families
among Jupiter Trojans starting from a data–base of semi-analytical proper
elements (Beaugé & Roig, 2001). The identification of families is performed
by comparing the mutual distances with a suitable metric in the proper el-
ements’ space. The clustering chain is halted when the mutual distance,
measuring the incremental velocity needed for orbital change after the puta-
tive parent body breakup, is larger than a fixed cut-off value. A lower cutoff
implies a higher statistical significance of the family. Since families in L4 are
on average more robust than those around L5 (Beaugé and Roig, 2001), we
prefer to adopt a cutoff of 100 m/s for the L4 cloud and of 150 m/s for L5.
For the very robust Eurybates family we decided to limit our survey to those
family members defined with a cutoff of 70 m/s.
All the data were acquired using the EMMI instrument, equipped with a
2x1 mosaic of 2048×4096 MIT/LL CCD with square 15µm pixels. For the
spectroscopic investigation during May 2004 and January 2005 runs we used
the grism #1 (150 gr/mm) in RILD mode to cover the wavelength range
4100–9400 Å with a dispersion of 3.1 Å/px (200 Å/mm) at the first order,
while on April 2003 we used a different grism, the #7 (150 gr/mm), covering
the spectral range 5200–9500 Å, with a dispersion of 3.6 Å/px at the first
order. April 2003 and January 2005 spectra were taken through a 1 arcsec
wide slit, while during May 2004 we used a larger slit (1.5 arcsec). The slit
was oriented along the parallactic angle during all the observing runs in order
to avoid flux loss due to the atmospheric differential refraction.
For most objects, the total exposure time was divided into several (usually
2-4) shorter acquisitions. This allowed us to check the asteroid position in
the slit before each acquisition, and correct the telescope pointing and/or
tracking rates if necessary. During each night we also recorded bias, flat–
field, calibration lamp (He-Ar) and several (6-7) spectra of solar analog stars
measured at different airmasses, covering the airmass range of the science
targets. During 17 January 2005, part of the night was lost due to some
technical problems and only 2 solar analog stars were acquired. The ratio
of these 2 stars show minimal variations (less than 1%) in the 5000–8400 Å
range, but higher differences at the edges of this range. For this reason we
omit the spectral region below 4800 Å for most of the asteroids acquired that
night.
The spectra were reduced using ordinary procedures of data reduction as
described in Fornasier et al. (2004a). The reflectivity of each asteroid was
obtained by dividing its spectrum by that of the solar analog star closest
in time and airmass to the object. Spectra were finally smoothed with a
median filter technique, using a box of 19 pixels in the spectral direction for
each point of the spectrum. The threshold was set to 0.1, meaning that the
original value was replaced by the median value if the median value differs
by more than 10% from the original one. The obtained spectra are shown
in Figs. 1–5. In Table 1 and Table 2 we report the circumstances of the
observations and the solar analog stars used respectively for the L5 and L4
family members.
[TABLE 3]
The broadband color data were obtained during the April 2003 and Jan-
uary 2005 runs just before the Trojans’ spectral observation. We used the
RILD mode of EMMI for wide field imaging with the Bessell-type B, V, R,
and I filters (centered respectively at 4139, 5426, 6410 and 7985Å). The ob-
servations were carried out in a 2 × 2 binning mode, yielding a pixel scale
of 0.33 arcsec/pixel. The exposure time varied with the object magnitude:
typically it was about 12-90s in V, 30-180s in B, 12-70s in R and I filters.
The CCD images were reduced and calibrated with a standard method (For-
nasier et al., 2004a), and absolute calibration was obtained through the ob-
servations of several Landolt fields (Landolt, 1992). The instrumental mag-
nitudes were measured using aperture photometry with an integrating radius
typically about three times the average seeing, and sky subtraction was per-
formed using a 5-10 pixels wide annulus around each object.
The results are reported in Table 3. From the visual inspection and the radial
profiles analysis of the images, no coma was detected for any of the observed
Trojans.
On May 2004, as the sky conditions were clear but not photometric, we did
not perform photometry of the Eurybates family targets.
3 Results
[TABLE 4 AND 5]
For each Trojan we computed the slope S of the spectral continuum using
a standard least squared technique for a linear fit in the wavelength range
between 5500 and 8000 Å. The choice of these wavelength limits has been
driven by the spectral coverage of our data. We choose 5500 Å as the lower
limit because of the different instrumental setup used during different ob-
serving runs (with some spectra starting at wavelength ≥ 5200 Å), while
beyond 8000 Å our spectra are generally noisier due to a combination of the
CCD drop-off in sensitivity and the presence of the strong atmospheric water
bands.
The computed slopes and errors are listed in Table 4 and 5. The reported er-
ror bars take into account the 1σ uncertainty of the linear fit plus 0.5%/103Å
attributable to the use of different instruments and solar analog stars (esti-
mated from the different efficiency of the grism used, and from flux losses
due to different slit apertures). In Table 4 and 5 we also report the taxo-
nomic class derived following the Dahlgren & Lagerkvist (1995) classification
scheme.
In the L5 cloud we find 27 D–, 3 DP–, 2 PD–, and 1 P–type objects. In
the L4 cloud we find 10 C–type and 7 P–type objects inside the Eurybates
family, while for the Menelaus, 1986 TS6 and 1986 WD families, including
the data published in Dotto et al. (2006), we get 9 D–, 3 P–, 3C–, and 1
DP–type asteroids.
The majority of the spectra are featureless, although some of the observed
Eurybates’ members show weak spectral absorption features (Fig. 5). These
features are discussed in the following section.
We derived an estimated absolute magnitude H by scaling the measured
V magnitude to r = ∆ = 1 AU and to zero phase assuming G=0.15 (Bowell
et al., 1989). The estimated H magnitude of each Trojan might be skewed
uncertain rotational phase, as the lightcurve amplitudes of Trojans might
vary up to 1 magnitude. In order to investigate possible size dependence in-
side each family, and considering that IRAS diameters are available for very
few objects, we estimate the size using the following relationship:
1329× 10−H/5
where D is the asteroid diameter, p is the geometric albedo, and H is the abso-
lute magnitude. We use H derived from our observations when available, and
from the ASTORB.DAT file (Lowell observatory) for the Eurybates mem-
bers, for which we did not carry out visible photometry. We evaluated the
diameter for an albedo range of 0.03–0.07, assuming a mean albedo of 0.04
for these dark asteroids (Fernandez et al., 2003). The resulting D values are
reported in Tables 4 and 5.
3.1 Dynamical families: L5 swarm
3.1.1 Anchises
[FIGURE 1]
We investigated 5 of the 15 members of the Anchises family (Fig. 1): 1173
Anchises, 23549 1994 ES6, 24452 2000 QU167, 47967 2000 SL298 and 124729
2001 SB173 on 17 January 2005. For 4 out of 5 observed objects we omit
the spectral range below 4800Å due to low S/N ratio and problems with the
solar analog stars. The spectral behavior is confirmed by photometric data
(see Table 3). All the obtained spectra are featureless.
The Anchises family survives at a cutoff corresponding to relative veloc-
ities of 150 m/s. The biggest member, 1173 Anchises, has a diameter of
126 km (IRAS data) and has the lowest spectral slope (3.9 %/103Å) among
the investigated family members. It is classified as P–type, while the other
4 members are all D–types. Anchises was previously observed in the 4000-
7400Å region by Jewitt & Luu (1990), who reported a spectral slope of 3.8
%/103Å, in perfect agreement with the value we found. The three 19-29 km
sized objects have a steeper spectral slope (7.4-9.2 %/103Å), while the small-
est object, 2001 SB173 (spectral slope = 14.78±0.99 %/103Å) is the reddest
one (Table 4).
Even with the uncertainties in the albedo and diameter, a slope–size rela-
tionship is evident among the observed objects, with smaller–fainter members
redder than larger ones (Fig. 7).
3.1.2 Misenus
[FIGURE 2]
For this family we investigated 6 members (11663 1997 GO24, 32794 1989
UE5, 56968 2000 SA92, 99328 2001 UY123, 105685 2000 SC51 and 120453
1988 RE12) out of the 12 grouped at a relative velocity of 150 m/s. The
family survives with the same members also at a stringent cut-off velocity
of 120 m/s. The spectra, together with magnitude color indices transformed
into linear reflectance, are shown in Fig. 2, while the color indices are reported
in Table 3. All the spectra are featureless with different spectral slope values
covering the 4.6–15.9 %/103Å range (Table 4): 1988 RE12 has the lowest
spectral slope and is classified as P–type, 3 objects (11663, 32794 and 2000
SC51) are in the transition region between P– and D– type, with very similar
spectral behavior, while the two other observed members are D–types. Of
these last, 56968 has the highest spectral slope not only inside the family
(15.86 %/103Å) but also inside the whole L5 sample analyzed in this paper.
All the investigated Misenus members are quite faint and have diameters
of a few tens of kilometers. No clear size-slope relationship has been found
inside this family (Fig. 7).
No other data on the Misenus family members are available in the literature,
so we do not know if the large gap between the spectral slope of 56968 and
those of the other 5 investigated objects is real or it could be filled by other
members not yet observed. If real, 56968 can be an interloper inside the
family.
3.1.3 Panthoos
[FIGURE 3]
The Panthoos family has 59 members for a relative velocity cutoff of 150
m/s. We obtained new spectroscopic and photometric data of 5 members:
4829 Sergestus, 30698 Hippokoon, 31821 1999 RK225, 76804 2000 QE and
111113 2001 VK85 (Fig. 3). Three objects presented by Fornasier et al.
(2004a) as belonging to the Astyanax family (23694 1997 KZ3, 32430 2000
RQ83, 30698 Hippokoon) and one to the background population (24444 2000
OP32) are now included among the members of the Panthoos family. Peri-
odic updates of the proper elements can change the family membership. In
particular the Astyanax group disappeared in the latest revision of dynami-
cal families, and its members are now in the Panthoos family within a cutoff
of 150m/s. The Panthos family survives also a cutoff of 120 m/s, with 7
members, and 90 m/s, with 6 members.
We observed 30698 Hippokoon during two different runs (on 9 Nov. 2002
and on 18 Jan. 2005), and both spectral slopes and colors are in agreement
inside the error bars (see Table 3, Table 4, and Fornasier et al., 2004a). No
other data on the Panthoos family are available in the literature.
The analysis of the 8 members (for 24444 only photometry is available)
show featureless spectra with slopes that seem to slightly increase as the
asteroid size decreases (Table 4 and Fig. 7). However, all the members have
dimensions very similar within the uncertainties, making it difficult for any
slope-size relationship to be studied. The largest member, 4829 Sergestus, is
a PD–type with a slope of about 5 %/103Å, while all the other investigated
members are D–types.
3.1.4 Cloantus
[FIGURE 4]
We observed only 2 out of 8 members of the Cloantus family (5511 Cloan-
thus and 51359 2000 SC17, see Fig. 4) as grouped at a cutoff corresponding
to relative velocities of 150 m/s. This family survives at a stringent cutoff
and 3 members (including the two that we observed) also survive for relative
velocities of 60 m/s. Both of the observed objects are D–types with very
similar, featureless, reddish spectra (Table 4 and Fig. 7). 5511 Cloanthus
was observed also by Bendjoya et al. (2004), who found a slope of 13.0±0.1
%/103Å in the 5000-7500 Å wavelength range, while we measure a value of
10.84±0.15 %/103Å. Our spectrum has a higher S/N ratio than the spectrum
by Bendjoya et al. (2004), and it is perfectly matched by our measured color
indices that confirm the spectral slope. This difference cannot be caused by
the slightly different spectral ranges used to measure the slope, but could
possibly be due to heterogeneous surface composition.
3.1.5 Phereclos
The Phereclos family comprises 15 members at a cutoff of 150 m/s. The
family survives with 8 members also at a cutoff of 120m/s. We obtained
spectroscopic and photometric data of 3 members (9030 1989 UX5, 11488
1988 RM11 and 31820 1999 RT186, see Fig 4), that, together with the 4
spectra (2357 Phereclos, 6998 Tithonus, 9430 1996 HU10, 18940 2000QV49)
already presented by Fornasier et al. (2004a), allow us to investigate about
half of the Phereclos family population defined at a cutoff of 150m/s. The
spectral slope of these objects, all classified as D–type except one PD–type
(11488), varies from 5.3 to 11.3 %/103Å (Table 4). The size of the fam-
ily members ranges from about 20 km in diameter for 31820 to 95 km for
2357, but we do not observe any clear slope-diameter relationship (Fig. 7 and
Table 4).
3.1.6 Sarpedon
We obtained new spectroscopic and photometric data of 2 members of the
Sarpedon family (48252 2001 TL212 and 84709 2002 VW120), whose spectra
and magnitude color indices are reported in Fig. 4 and Table 4. Including
the previous observations (Fornasier et al., 2004a) of 4 other members (2223
Sarpedon, 5130 Ilioneus, 17416 1988 RR10, and 25347 1999 RQ116), we have
measurements of 6 of the 21 members of this family dynamically defined at a
cutoff of 150 m/s. All the 6 aforementioned objects, except 25347, constitute
a robust clustering which survives up to 90 m/s with 9 members. The cluster
which contains (2223) Sarpedon was also recognized as a family by Milani
(1993).
All the 6 investigated members have very similar colors (see Table 3) and
spectral behavior. The spectral slope (Fig. 7) varies over a very restricted
range, from 9.6 to 11.6 %/103Å (Table 4), despite a significant variation of
the estimated size (from the 18 km of 17416 to the 105 km of 2223). Con-
sequently, the surface composition of the Sarpedon family members appears
to be very homogeneous.
3.2 Dynamical families: L4 swarm
3.2.1 Eurybates
[FIGURE 5]
Eurybates family members were observed in May 2004. The selection of
the targets was made on the basis of a very stringent cutoff, corresponding
to relative velocities of 70 m/s, that gives a family population of 28 objects.
We observed 17 of these members (see Table 2) that constitute a very robust
clustering in the space of the proper elements: all the members we studied,
except 2002 CT22, survive at a cutoff of 40 m/s.
The spectral behavior of these objects (Fig. 5) is quite homogeneous with
10 asteroids classified as C–type and 7 as P–type. The spectral slopes (Ta-
ble 5) range from neutral to moderately red (from -0.5 to 4.6 %/103Å). The
slopes of six members are close to zero (3 slightly negative) with solar-like col-
ors. The asteroids 18060, 24380, 24420, and 39285, all classified as C–types,
clearly show a drop off of reflectance for wavelength shorter than 5000–5200
Å. The presence of the same feature in the spectra of 2 other members (1996
RD29 and 28958) is less certain due to the lower S/N ratio. This absorp-
tion is commonly seen on main belt C–type asteroids (Vilas 1994; Fornasier
et al. 1999), where is due to the intervalence charge transfer transitions
(IVCT) in oxidized iron, and is often coupled with other visible absorption
features related to the presence of aqueous alteration products (e.g. phyl-
losilicates, oxides, etc). These IVCTs comprise multiple absorptions that are
not uniquely indicative of phyllosilicates, but are present in the spectrum
of any object containing Fe2+ and Fe3+ in its surface material (Vilas 1994).
Since no other phyllosilicate absorption features are present in the C-type
spectra of the Eurybates family, there is no evidence that aqueous alteration
processes occurred on the surface of these bodies.
In Fig. 8 we show the spectral slopes versus the estimated diameters for
the Eurybates family members. All the observed objects, except the largest
member (3548) that has a diameter of about 70 km and exhibit a neutral (∼
solar-like) spectral slope, are smaller than ∼ 40 km and present both neutral
and moderately red colors. The spectral slopes are strongly clustered around
S = 2%/103Å, with higher S values restricted to smaller objects (D< 25
3.2.2 1986 WD
[FIGURE 6]
We investigated 6 out of 17 members of the 4035 1986 WD family that is
dynamically defined at a cutoff of 130 m/s (Fig. 6 and Table 2). Three of our
targets (4035, 6545 and 11351) were already observed by Dotto et al. (2006):
for 6545 and 11351 there is a good consistency between our spectra and
those already published. 4035 was observed also by Bendjoya et al. (2004):
all the spectra are featureless, but Bendjoya et al. (2004) obtain a slope of
8.8 %/103Å, comparable to the one here presented, while Dotto et al. (2006)
found a higher value (see Table 5). This could be interpreted as due to the
different rotational phases seen in the three observations, and could indicate
some inhomogeneities on the surface of 4035.
The observed family members show heterogeneous behaviors (Fig. 8),
with spectral slopes ranging from neutral values for the smaller members
(24341 and 14707) to reddish ones for the 3 members with size bigger than
50 km (4035, 6545, and 11351). For this family, it seems that a size-slope
relationship exists, with smaller members having solar colors and spectral
slopes increasing with the object’ sizes.
3.2.3 1986 TS6
The 1986 TS6 family includes 20 objects at a cut-off of 100 m/s. We present
new spectroscopy and photometry of a single member, 12921 1998 WZ5
(Fig. 6). The spectrum we present here is flat and featureless, with a spectral
slope of 4.6±0.8%/103Å. Dotto et al. (2006) presented a spectrum obtained
a month after our data (in May 2003) that has a very similar spectral slope
3.7± 0.8%/103Å. Previously, 12917 1998 TG16, 13463 Antiphos, 12921 1998
WZ5, 15535 2000 AT177, 20738 1999 XG191, and 24390 2000 AD177 were
included in the Makhoan family. Refined proper elements now place all of
these bodies in the 1986 TS6 family.
In Fig. 8 we report the spectral slopes vs. estimated diameters of the
6 observed members. The family shows different spectral slopes with the
presence of both P–type (12921 and 13463) and D–type asteroids (12917,
15535, 20738, and 24390). Due to the very similar diameters, a slope-size
relationship is not found.
[FIGURE 7 AND 8]
4 Discussion
The spectra of Jupiter Trojan members of dynamical families show a range
of spectral variation from C– to D–type asteroids. With the exception of the
L4 Eurybates family, all the observed objects have featureless spectra, and
we cannot find any spectral bands which could help in the identification of
minerals present on their surfaces. The lack of detection of any mineralogy
diagnostic feature might indicate the formation of a thick mantle on the Tro-
jan surfaces. Such a mantle could be formed by a phase of cometary activity
and/or by space weathering processes as demonstrated by laboratory exper-
iments on originally icy surfaces (Moore et al., 1983; Thompson et al., 1987;
Strazzulla et al., 1998; Hudson & Moore, 1999).
A peculiar case is constituted by the Eurybates family, which shows a pre-
ponderance of C–type objects and a total absence of D–types. Moreover,
this is the only family in which some members exhibit spectral features at
wavelengths shorter than 5000–5200 Å, most likely due to the intervalence
charge transitions in materials containing oxidized iron (Vilas 1994).
4.1 Size vs spectral slope distribution:
Individual families
The plots of spectral slopes vs. diameters are shown in Fig. 7 and 8. A
relationship between spectral slopes and diameters seems to exist for only
three of the nine families we studied. In the Anchises and Panthoos families,
smaller objects have redder spectra, while for the 1986 WD family larger
objects have the redder spectra.
Moroz et al. (2004) have shown that ion irradiation on natural complex
hydrocarbons gradually neutralizes the spectral slopes of these red organic
solids. If the process studied by Moroz et al. (2004) occurred on the surface of
Jupiter Trojans, the objects having redder spectra have to be younger than
those characterized by bluish-neutral spectra. In this scenario the largest
and spectrally reddest objects of the 1986 WD family could come from the
interior of the parent body and expose fresh material. In the case of the
Anchises and Panthoos families the spectrally reddest members, being the
smallest, could come from the interior of the parent body, or alternatively
could be produced by more recent secondary fragmentations. In particular,
small family members may be more easily resurfaced, as significant collisions
(an impactor having a size greater than a few percent of the target), as well
as seismic shaking and recoating by fresh dust, may occur frequently at small
sizes.
[FIGURE 9]
4.2 Size vs slope distribution:
The Trojan population as a whole
[TABLE 6]
As compared to the data available in literature, our sample strongly con-
tributed to the analysis of fainter and smaller Trojans, with estimated di-
ameters smaller than 50 km. Jewitt & Luu (1990), analyzing a sample of
32 Trojans, found that the smaller objects were redder than the bigger ones.
However, our data play against the existence of a possible color-dimension
trend. In fact, the spectral slope’s range of the objects smaller than 50 km
is similar to that of the larger Trojans, as shown in Fig. 9.
The Eurybates family strongly contributes to the population of small
spectrally neutral objects, filling the region of bodies with mean diameter
D<40 km and with spectral slopes smaller than 3 %/103Å.
In order to carry out a complete analysis of the spectroscopic and pho-
tometric characteristics of the whole available data set on Jupiter Trojans,
we considered all the visible spectra published in the literature: Jewitt &
Luu (1990, 32 objects), Fitzimmons et al. (1994, 3 objects), Bendjoya et
al. (2004, 34 objects), Fornasier et al. (2004a, 26 L5 objects), and Dotto et
al. (2006, 24 L4 Trojans). We also add several Trojans spectra (11 L4 and
3 L5 Trojans) from the files available on line (Planetary data System archive,
pdssbn.astro.umd.edu, and www.daf.on.br/∼lazzaro/S3OS2-Pub/s3os2.htm)
from the SMASS I, SMASS II and S3OS2 surveys (Xu et al., 1995; Bus &
Binzel, 2003; Lazzaro et al., 2004). Including all these data, we compile a
sample of 142 different Trojans, 68 belonging to the L5 cloud and 74 belong-
ing to the L4. We performed the taxonomic classification of this enlarged
sample, on the basis of the Dahlgren and Lagerkvist (1995) scheme, by ana-
lyzing spectral slopes computed in the range 5500-8000 Å. Different authors,
of course, considered different spectral ranges for their own slope gradient
evaluations: Jewitt & Luu (1990) and Fitzimmons et al. (1994) use the
4000-7400 Å and Bendjoya et al. (2004, Table 2) used a slightly different
ranges around 5200-7500 Å. Since all the cited papers show spectra with lin-
ear featureless trends, the different wavelength ranges used for the spectral
gradient computation by Bendjoya et al. (2004) and Jewitt & Luu (1990)
are not expected to influence the obtained slopes.
In order to search for a dependency of the spectral slope distribution with
the size of the objects, all observations (from this paper as well as from the
literature) were combined. The objects were isolated in 5 size bins (smaller
than 25 km, 25–50 km, 50–75 km, 75-100 km and larger than 100 km). Each
bin contains between 20 and 50 objects. These subsamples are large enough
to be compared using classical statistical tests: the t-test, which estimates
if the mean values are compatible, the f-test, which checks if the widths of
the distributions are compatible (even if they have different means), and
the KS test, which compares directly the full distributions. A probability is
computed for each test; a small probability indicates that the tested distri-
butions are not compatible, i.e. the objects are not randomly extracted from
the same population, while a large probability value has no meaning (i.e. it
is not possible to assure that both samples come from the same population,
we can just say in that case that they are not incompatible). In order to
quantify the probability levels that we consider as significant, the same tests
were run on randomized distributions (see Hainaut & Delsanti 2002 for the
method). Since probability lower than 0.04-0.05 does not appear in these
randomized distributions, we consider that values smaller than 0.05 indicate
a significant incompatibility.
Each sub-sample was compared with the four others – the results are sum-
marized in Table 6. The average slope of the 5 bins are all compatible among
each other. The only marginally significant result is that the width of the
slope distribution among the larger objects (diam. > 100 km) is narrower
than that of all the smaller objects.
This narrower color distribution could be due to the aging processes affect-
ing the surface of bigger objects, which are supposed to be older. The wider
color distribution of small members is possibly related to the different ages
of their surfaces: some of them could be quite old, while some other could
have been recently refreshed.
4.3 Spectral slopes and L4/L5 Clouds
[HERE FIGURES 10 AND 11]
Considering only the Trojan observations reported in this paper, the aver-
age slope is 8.84±3.03%/103Å for the L5 population, and 4.57±4.01%/103Å
for the L4.
Considering now all the spectra available in the literature, the 68 L5
Trojans have an average slope of 9.15±4.19%/103Å, and the 78 L4 objects,
6.10±4.48%/103Å. Performing the same statistical tests as above, it appears
that these two populations are significantly different. In particular, the av-
erage slopes are incompatible at the 10−5 level.
Nevertheless, as described in Section 3.2.1, the Eurybates family members
have quite different spectral characteristics than the other objects and con-
stitute a large subset of the whole sample. Indeed, comparing their distribu-
tion with the whole populations, they are found significantly different at the
10−10 level. In other words, the Eurybates family members do not constitute
a random subset of the other Trojans.
Once excluded the Eurybates family, the remaining 61 Trojans from the
L4 swarm have an average slope of 7.33±4.24%/103Å. The very slight dif-
ference of average slope between the L5 and remaining L4 objects is very
marginally significant (probability of 1.6%), and the shape and width of the
slope distributions are compatible with each other.
The taxonomic classification we have performed shows that the majority
(73.5%) of the observed L5 Trojans (Fig. 10) are D–type (slope > 7 %/103
Å) with featureless reddish spectra, 11.8% are DP/PD –type (slope between
5 and 7 %/103 Å), 10.3% are P–type, and only 3 objects are classified as
C–type (4.4%).
In the L4 swarm (Fig. 11), even though the D–type still dominate the
population (48.6%), the spectral types are more heterogeneous as compared
to the L5 cloud, with a higher percentage of neutral-bluish objects: 20.3%
are P–type, 8.1% are DP/PD-type, 12.2% are C–type, and 10.8% of the
bodies have negative spectral slope. The higher percentage of C– and P–
type as compared to the L5 swarm is strongly associated with the presence
of the very peculiar Eurybates family. Among 17 observed members 10 are
classified as C–types (among which 3 have negative spectral slopes) and 7 are
P–types. Considering the 57 asteroids that compose the L4 cloud without
the Eurybates family, we find percentages of P, and PD/DP –types very
similar to those of the L5 cloud (14.0% and 10.5% respectively), a smaller
percentage of D–types (63.2%) and of the C–types (3.5%), and the presence
of a 8.8% Trojans with negative spectral slopes.
The visible spectra of the Eurybates members are very similar to those of
C–type main belt asteroids, Chiron-like Centaurs, and cometary nuclei. This
similarity is compatible with three different scenarios: the family could have
been produced by the fragmentation of a parent body very different from
all the other Jupiter Trojans (in which case the origin of such a peculiar
parent must still be assessed); this could be a very old family where space
weathering processes have covered any differences in composition among the
family members and flattened all the spectra; this could be a young family
where space weathering processes occurred within time scales smaller than
the age of the family. In the last two cases the Eurybates family would give
the first observational evidence of spectra flattened owing to space weathering
processes. This would then imply, according to the results of Moroz et al.
(2004), that its primordial composition was rich in complex hydrocarbons.
The knowledge of the age of the Eurybates family is therefore a fundamental
step to investigate the nature and the origin of the parent body, and to assess
the effect of space weathering processes on the surfaces of its members.
The present sample of Jupiter Trojans suggests a more heterogeneous
composition of the L4 swarm as compared to the L5 one. As previously
noted by Bendjoya et al. (2004), the L4 swarm contains a higher percentage
of C– and P–type objects. This result is enhanced by members of the Eu-
rybates family, but remains even when these family members are excluded.
Moreover, the dynamical families belonging to the L4 cloud are more robust
than those of the L5 one, surviving with densely populated clustering even
at low relative velocity cut-off. We therefore could argue that the L4 cloud
is more collisionally active than the L5 swarm. Nevertheless, we still cannot
intepret this in terms of the composition of the two populations, since we
cannot exclude that as yet unobserved C– and P–type families are present
in the L5 cloud.
4.4 Orbital Elements
[HERE FIGURE 12 and TABLES 7 and 8 ]
We analyzed the spectral slope as a function of the Trojans’ orbital el-
ements. As an illustration, Fig. 12 shows the B − R color distribution as
a function of the orbital elements. In order to investigate variations with
orbital parameters, the Trojan population is divided in 2 sub samples: those
with the considered orbital element lower than the median value, and those
with the orbital element higher than the median (by construction, the two
subsamples have the same size). Taking a as an example, half the Trojans
have a < 5.21AU, and half have a larger than this value.
The mean color, the color dispersion, and the color distribution of the
2 subsamples are compared using the three statistical tests mentioned in
Section 4.2. The method is discussed in details in Hainaut & Delsanti (2002).
The tests are repeated for all color and spectral slope distributions. The
results are the following.
• q, perihelion distance: the color distribution of the Trojans with small
q is marginally broader than that of Trojans with larger q. This result
is not very strong (5%), and is dominated by the red-end of the visible
wavelength. Removing the Eurybates from the sample maintains the
result, at the same weak level.
• e, eccentricity: the distribution shows a similar result, also at the weak
5% significance. The objects with larger e have broader color distribu-
tion than those with lower e. This result is entirely dominated by the
Eurybates’ contribution.
• i, inclination: objects with smaller inclination are significantly bluer
than those with larger i. This result is observed at all wavelengths. It
is worth noting that this is contrary to what is usually observed on
other Minor Bodies in the Outer Solar System survey (MBOSSes),
where objects with high i, or more generally, high excitation E =√
e2 + sin2 i, are bluer (Hainaut & Delsanti, 2002; Doressoundiram et
al., 2005). This can also be visually appreciated in Fig. 12. This result
is also completely dominated by the Eurybates’ contribution. The non-
Eurybates Trojans do not display this trend.
• E =
e2 + sin2 i, orbital excitation: the objects with small E are also
significantly bluer than those with high E. This result is also com-
pletely dominated by the Eurybates’ contribution. The non-Eurybates
Trojans do not display this trend.
In summary this analysis shows that the Eurybates sub-sample of the
Trojans is well separated in orbital elements and in colors.
For the other Minor Bodies in the outer Solar System, the relation be-
tween color and inclination–orbital excitation (objects with a higher orbital
excitation tend to be bluer) is interpreted as a relation between excitation
and surface aging/rejuvenating processes (Doressoudiram et al., 2005). The
Eurybates family has low excitation and neutral-blue colors, suggesting that
the aging/rejuvenating processes affecting them are different from the other
objects. This could be due to different surface compositions, different irradi-
ation processes, or different collisional properties – which would be natural
for a collisional family.
5 Comparison with other outer Solar System
minor bodied
5.1 Introduction and methods
[HERE FIGURES 13 AND 14]
The statistical tests set described in section 4.2 has also been applied
to compare the colors and the spectral slopes distribution of the Trojans
with those of the other minor bodies in the outer Solar System taken from
the updated, on-line version of the Hainaut & Delsanti (2002) database.
Figure 13, as an example, displays the (R-I) vs (V-R) diagrams, while Fig. 14
shows the (B-V) and (V-R) color distributions, as well as the spectral slope
distribution of the different classes of objects. The tests were performed
on all the color indices derived from filters in the visible (UBVRI) and near
infrared range (JHK) but in Table 7 and 8 we summarize the most significant
results.
In order to “calibrate” the significant probabilities, additional artificial
classes are also compared: first, the objects which have an even internal
number in the database with the odd ones. As this internal number is purely
arbitrary, both classes are statistically indistinguishable. The other tested
pair is the objects with a “1999” designation versus the others. Again, this
selection criterion is arbitrary, so the pseudo-classes it generates are sub-
sample of the total population, and should be indistinguishable. However, as
many more objects have been discovered in all the other years than during
that specific year, the size of these sub-samples are very different. This
permits us to estimate the sensitivity of the tests on sample of very different
sizes. Some of the tests found the arbitrary populations incompatible at the
5% level, so we use 0.5% as a conservative threshold for statistical significance
of the distribution incompatibility
5.2 Results
Table 7 and Fig. 14 clearly show that the Trojans’ colors distribution is
different as compare to that of Centaurs, TNOs and comets. Trojans are at
the same time bluer, and their distribution is narrower than all the other
populations. Using the statistical tests (see Table 8), we can confirm the
significance of these results.
• The average colors of the Trojans are significantly different from those
of all the other classes of objects (t-test), with the notable exception of
the short period comet nuclei. Refining the test to the Eurybates/non-
Eurybates, it appears that the Eurybates have marginally different
mean colors, while the non-Eurybates average colors are indistinguish-
able from those of the comets.
• Considering the full shape of the distribution (KS test), we obtain the
same results: the Trojans colors distributions are significantly differ-
ent from those of all the other classes, with the exception of the SP
comets, which are compatible. Again, this result becomes stronger
separating the Eurybates: their distributions are different from those
of the comets, while the non-Eurybates ones are indistinguishable.
• The results when considering the widths of the color distributions (f-
test) are slightly different. Classes of objects with different mean colors
could still have the same distribution width. This could suggest that a
similar process (causing the width of the distribution) is in action, but
reached a different equilibrium point (resulting in different mean val-
ues). This time, all classes are incompatible with the Trojans, including
the comets, with strong statistical significance.
In order to further explore possible similarities between Trojans and other
classes, the comparisons were also performed with the neutral Centaurs.
These were selected with S < 20%/103Å); this cut-off line falls in the gap
between the ”neutral” and ”red” Centaurs (Peixinho et al., 2003, Fornasier
et al., 2004b).
The t-Test (mean color) only reveals a very moderate incompatibility be-
tween the Trojans and neutral Centaurs, at the 5% level, i.e. only marginally
significant. On the other hand, the f-Test gives some strong incompatibilities
in various colors (moderate in B− V and H −K, very strong in R− I), but
the two populations are compatible for most of the other colors. Similarly,
only the R − I KS-test reveals a strong incompatibility. It should also be
noted that only 18 neutral Centaurs are known in the database. In summary,
while the Trojans and neutral Centaurs have fairly similar mean colors, their
color distributions are also different.
6 Conclusions
From 2002, we carried out a spectroscopic and photometric survey of Jupiter
Trojans, with the aim of investigating the members of dynamical families.
In this paper we present new data on 47 objects belonging to several dy-
namical families: Anchises (5 members), Cloanthus (2 members), Misenus (6
members), Phereclos (3 members), Sarpedon (2 members) and Panthoos (5
members) from the L5 swarm; Eurybates (17 members), 1986 WD (6 mem-
bers), and Menelaus (1 member) for the L4 swarm. Together with the data
already published by Fornasier et al. (2004a) and Dotto et al. (2006), taken
within the same observing program, we have a total sample of 80 Trojans,
the largest homogeneous data set available to date on these primitive aster-
oids. The main results coming from the observations presented here, and
from the analysis including previously published visible spectra of Trojans,
are the following:
• Trojans’ visible spectra are mostly featureless. However, some mem-
bers of the Eurybates family show a UV drop-off in reflectivity for
wavelength shorter than 5000–5200 Å that is possibly due to interva-
lence charge transfer transitions (IVCT) in oxidized iron.
• The L4 Eurybates family strongly differs from all the other families
in that it is dominated by C– and P–type asteroids. Also its spectral
slope distribution is significantly different when compared to that of
the other Trojans (at the 10−10 level).
This family is very peculiar and is dynamically very strong, as it sur-
vives also at a very stringent cutoff (40 m/s). Further observations in
the near-infrared region are strongly encouraged to look for possible
absorption features due to water ice or to material that experienced
aqueous alteration.
• The average spectral slope for the L5 Trojans is 9.15±4.19%/103Å, and
6.10±4.48%/103Å for the L4 objects. Excluding the Eurybates, the L4
average slope values becomes 7.33±4.24%/103Å. The slope distribu-
tions of the L5 and of the non-Eurybates L4 are indistinguishable.
• Both L4 and L5 clouds are dominated by D–type asteroids, but the L4
swarm has an higher presence of C– and P–type asteroids, even when
the Eurybates family is excluded, and appears more heterogeneous in
composition as compared to the L5 one.
• We do not find any size versus spectral slope relationship inside the
whole Trojans population.
• The Trojans with higher orbital inclination are significantly redder than
those with lower i. While this trend is the opposite of that observed
for other distant minor bodies, this effect is entirely dominated by the
Eurybates family.
• Comparing the Trojans colors with those of other distant minor bod-
ies, they are the bluest of all classes, and their colors distribution is the
narrowest. This difference is mostly due to the Eurybates family. In
fact, if we consider only the Trojan population without the Eurybates
members, their average colors and overall distributions are not distin-
guishable from that of the short period comets. However, the widths
of their color distributions are not compatible. The similarity in the
overall color distributions might be caused by the small size of the short
period comet sample rather than by a physical analogy. The Trojans
average colors are also fairly similar to those of the neutral Centaurs,
but the overall distributions are not compatible.
After this study, we have to conclude that Trojans have peculiar charac-
teristics very different from those of all the other populations of the outer
Solar System.
Unfortunately, we still cannot assess if this is due to differences in the physi-
cal nature, or in the aging/rejuvenating processes which modified the surface
materials in different way at different solar distances. Further observations,
mainly in V+NIR spectroscopy and polarimetry, are absolutely needed to
better investigate the nature of Jupiter Trojans and to definitively assess if a
genetical link might exist with Trans-Neptunian Objects, Centaurs and short
period comets.
Acknowledgments
We thank Beaugé and Roig for kindly providing us with updated Trojan
family list, and R.P. Binzel and J.P. Emery for their useful comments in the
reviewing process.
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Tables
Table 1: Observing conditions of the investigated L5 asteroids. For each
object we report the observational date and universal time, total exposure
time, number of acquisitions with exposure time of each acquisition, airmass,
and the observed solar analogs with their airmass.
Obj Date UT Texp (s) nexp air. Solar An. (air.)
Anchises
1173 17 Jan 05 06:06 60 1×60s 1.42 HD76151 (1.48)
23549 17 Jan 05 07:20 480 2×240s 1.60 HD76151 (1.48)
24452 17 Jan 05 07:54 960 4×240s 1.44 HD76151 (1.48)
47967 17 Jan 05 05:34 800 2×400s 1.38 HD76151 (1.48)
2001 SB173 17 Jan 05 06:28 1200 2×600s 1.35 HD76151 (1.48)
Cloanthus
5511 19 Jan 05 06:04 960 4×240s 1.26 HD76151 (1.12)
51359 19 Jan 05 04:13 660 1×660s 1.36 HD76151 (1.12)
Misenus
11663 17 Jan 05 05:13 400 1×400s 1.21 HD44594 (1.12)
32794 18 Jan 05 03:13 1800 2×900s 1.39 HD28099 (1.44)
56968 17 Jan 05 04:31 400 2×400s 1.21 HD44594 (1.12)
1988 RE12 18 Jan 05 04:12 2000 2×1000s 1.31 HD28099 (1.44)
2000 SC51 18 Jan 05 06:09 1320 2×660s 1.16 HD44594 (1.17)
2001 UY123 18 Jan 05 06:46 1320 2×660s 1.32 HD44594 (1.17)
Phereclos
9030 18 Jan 05 08:19 1000 1×1000s 1.37 HD44594 (1.17)
11488 19 Jan 05 03:31 1320 2×660s 1.99 HD76151 (1.12)
31820 19 Jan 05 07:02 1320 2×660s 1.35 HD76151 (1.11)
Sarpedon
48252 18 Jan 05 02:32 1320 2×660s 1.30 HD28099 (1.44)
84709 19 Jan 05 05:35 1320 2×660s 1.34 HD76151 (1.12)
Panthoos
4829 17 Jan 05 08:37 720 3×240s 1.45 HD76151 (1.48)
30698 18 Jan 05 01:54 1320 2×660s 1.73 HD28099 (1.44)
31821 18 Jan 05 05:27 1320 2×660s 1.35 HD28099 (1.44)
76804 17 Jan 05 03:35 1800 3×600s 1.38 HD44594 (1.12)
2001 VK85 18 Jan 05 07:31 2000 2×1000s 1.23 HD44594 (1.17)
Table 2: Observing conditions of the investigated L4 asteroids. For each
object we report the observational date and universal time, total exposure
time, number of acquisitions with exposure time of each acquisition, airmass,
and the observed solar analogs with their airmass.
Obj Date UT Texp (s) nexp air. Solar An. (air.)
Eurybates
3548 25 May 04 05:14 600 2×300s 1.02 SA107-684 (1.19)
9818 26 May 04 00:13 780 1×780s 1.19 SA102-1081(1.15)
13862 25 May 04 03:35 1200 2×600s 1.09 SA107-998 (1.15)
18060 25 May 04 02:47 1500 2×750s 1.07 SA107-998 (1.15)
24380 25 May 04 06:53 780 1×780s 1.18 SA107-684 (1.19)
24420 25 May 04 08:49 900 1×900s 1.59 SA112-1333 (1.17)
24426 26 May 04 00:13 1440 2×720s 1.13 SA107-684 (1.17)
28958 26 May 04 07:14 1800 2×900s 1.35 SA107-684 (1.17)
39285 25 May 04 05:40 2700 3×900s 1.09 SA107-684 (1.19)
43212 25 May 04 07:39 2340 3×780s 1.39 SA110-361 (1.15)
53469 25 May 04 02:05 1800 2×900s 1.04 SA107-998 (1.15)
65150 26 May 04 01:59 3600 4×900s 1.07 SA102-1081 (1.20)
65225 26 May 04 03:40 3600 4×900s 1.04 SA107-684 (1.17)
1996RD29 26 May 04 05:12 2700 3×900s 1.10 SA107-684 (1.17)
2000AT44 25 May 04 04:14 1800 2×900s 1.04 SA107-684 (1.19)
2002CT22 26 May 04 00:49 2400 4×600s 1.08 SA102-1081 (1.15)
2002EN68 26 May 04 08:10 1800 2×900s 1.62 SA107-684 (1.17)
1986 WD
4035 10 Apr 03 03:28 600 1×600s 1.09 SA107-684 (1.15)
6545 10 Apr 03 02:39 900 1×900s 1.16 SA107-684 (1.15)
11351 10 Apr 03 09:21 900 1×900s 1.28 SA107-684 (1.15)
14707 11 Apr 03 08:11 1200 1×1200s 1.15 SA107-684 (1.15)
24233 11 Apr 03 02:29 1200 1×1200s 1.39 SA107-684 (1.37)
24341 11 Apr 03 05:47 900 1×900s 1.16 SA107-684 (1.17)
1986 TS6
12921 10 Apr 03 07:33 900 1×900s 1.39 SA107-684 (1.40)
Table 3: Visible photometric observations of L4 and L5 Trojans (ESO-NTT
EMMI): for each object, date, computed V magnitude, B-V, V-R and V-
I colors are reported. The given UT is for the V filter acquisition. The
observing photometric sequence (V-R-B-I) took a few minutes.
Object date UT V B-V V-R V-I
1986 WD
4035 10 Apr 03 03:11 16.892±0.031 0.752±0.040 0.473±0.042 0.926±0.055
4035 10 Apr 03 04:22 16.981±0.031 0.752±0.040 0.495±0.042 0.945±0.055
6545 10 Apr 03 02:22 17.558±0.031 0.734±0.041 0.499±0.042 0.935±0.055
11351 10 Apr 03 09:03 18.407±0.032 0.739±0.044 0.498±0.044 0.900±0.057
14707 11 Apr 03 06:46 18.666±0.031 0.751±0.041 0.401±0.033 0.804±0.055
14707 11 Apr 03 08:37 18.873±0.031 0.754±0.041 0.424±0.033 0.790±0.056
24233 11 Apr 03 01:33 18.894±0.034 0.704±0.051 0.481±0.037 0.899±0.058
24341 11 Apr 03 05:05 19.376±0.032 0.713±0.043 0.369±0.035 0.759±0.057
1986 TS6
12921 10 Apr 03 07:12 18.393±0.031 0.673±0.040 0.421±0.042 0.786±0.055
L5 cut off 150m/s
Anchises
1173 17 Jan 05 05:54 16.595±0.024 0.811±0.034 0.402±0.035 0.805±0.038
23549 17 Jan 05 07:09 18.969±0.050 0.800±0.071 0.485±0.068 0.872±0.075
24452 17 Jan 05 07:48 18.757±0.043 0.872±0.056 0.441±0.056 0.847±0.066
47967 17 Jan 05 05:27 19.382±0.044 0.899±0.058 0.489±0.069 0.965±0.075
2001 SB173 17 Jan 05 06:20 19.882±0.043 0.992±0.060 0.503±0.064 0.927±0.078
Cloanthus
5511 19 Jan 05 05:52 17.968±0.020 0.906±0.027 0.442±0.027 0.968±0.032
51359 19 Jan 05 03:54 19.631±0.102 0.864±0.201 0.447±0.131 0.885±0.164
Misenus
11663 17 Jan 05 05:05 18.473±0.022 0.837±0.030 0.409±0.030 0.872±0.039
32794 18 Jan 05 03:07 19.685±0.038 0.923±0.065 0.393±0.056 0.879±0.057
56968 17 Jan 05 04:18 18.596±0.026 0.986±0.040 0.494±0.033 1.003±0.036
1988 RE12 18 Jan 05 04:00 20.892±0.081 0.826±0.132 0.388±0.108 0.871±0.106
2000 SC51 18 Jan 05 06:03 19.876±0.038 1.016±0.055 0.444±0.059 0.896±0.056
2001 UY123 18 Jan 05 06:41 19.869±0.047 0.890±0.058 0.537±0.056 0.971±0.063
Phereclos
9030 18 Jan 05 08:14 18.397±0.020 0.887±0.024 0.493±0.027 0.973±0.028
11488 19 Jan 05 02:57 18.931±0.066 0.868±0.101 0.430±0.079 0.848±0.084
31820 19 Jan 05 06:39 20.041±0.077 0.889±0.093 0.520±0.091 0.916±0.123
Sarpedon
48252 18 Jan 05 02:25 19.878±0.060 0.949±0.100 0.467±0.093 0.903±0.090
84709 19 Jan 05 05:10 19.862±0.068 0.855±0.087 0.462±0.090 1.010±0.094
Panthoos
4829 17 Jan 05 08:18 18.430±0.029 0.851±0.050 0.420±0.039 0.792±0.052
30698 18 Jan 05 01:45 19.353±0.036 – 0.472±0.042 0.865±0.047
31821 18 Jan 05 05:21 19.328±0.076 0.980±0.111 0.440±0.097 0.901±0.108
76804 17 Jan 05 03:21 19.471±0.065 0.803±0.082 0.446±0.070 0.889± 0.080
2001 VK85 18 Jan 05 07:23 20.179±0.038 0.822±0.063 0.462±0.048 1.020±0.050
Table 4: L5 families. We report for each target the absolute magnitude H
and the estimated diameter (diameters marked by ∗ are taken from IRAS
data), the spectral slope S computed between 5500 and 8000 Å and the
taxonomic class (T) derived following Dahlgren & Lagerkvist (1995) classi-
fication scheme. The asteroids marked with a were observed by Fornasier
et al. (2004a), and their spectral slope values have been recomputed in the
5500-8000 Å wavelength range; asteroids 23694, 30698 and 32430, previously
Astyanax members, have been reassigned to the Panthoos family due to re-
fined proper elements.
Obj H D (km) S (%/103Å) T
Anchises
1173 8.99 ∗126+11
3.87±0.70 P
23549 12.04 26+4
8.49±0.88 D
24452 11.85 29+5
7.42±0.70 D
47967 12.15 25+4
9.21±0.78 D
2001 SB173 12.77 19+3
14.78±0.99 D
Cloanthus
5511 10.43 55+8
10.84±0.65 D
51359 12.25 24+6
12.63±1.30 D
Misenus
11663 10.95 44+7
6.91±0.70 DP
32794 12.77 19+3
6.59±0.88 DP
56968 11.72 30+5
15.86±0.71 D
1988 RE12 13.20 16+2
4.68±1.20 P
2000 SC51 12.69 20+3
6.54±0.98 DP
2001 UY123 12.75 19+3
8.28±0.88 D
Phereclos
a2357 8.86 ∗95+4
9.91±0.68 D
a6998 11.43 34+5
11.30±0.75 D
9030 11.14 40+6
10.35±0.76 D
a9430 11.47 35+5
10.02±0.90 D
11488 11.82 29+5
5.37±0.92 PD
a18940 11.81 29+4
7.13±0.75 D
31820 12.63 20+3
7.53±0.80 D
Sarpedon
a2223 9.25 ∗95+4
10.20±0.65 D
a5130 9.85 71+11
10.45±0.65 D
a17416 12.83 18+3
10.80±0.90 D
a25347 11.59 33+5
10.11±0.83 D
48252 12.84 18+3
9.62±0.82 D
84709 12.70 19+3
11.64±0.84 D
Panthoos
4829 11.16 39+6
5.03±0.70 PD
a23694 11.61 32+5
8.20±0.72 D
30698 12.14 25+4
8.23±1.00 D
a30698 12.27 25+4
9.08±0.82 D
a32430 12.23 25+4
8.12±1.00 D
31821 11.99 27+4
10.58±0.82 D
76804 12.16 25+4
7.29±0.71 D
2001 VK85 12.79 19+3
14.39±0.81 D
Table 5: L4 Families. We report for each target the absolute magnitude H
and the estimated diameter (diameters marked by ∗ are taken from IRAS
data, while absolute magnitudes marked by ∗∗ are taken from the astorb.dat
file of the Lowell Observatory), the spectral slope S computed between 5500
and 8000 Å, and the taxonomic class (T) derived following Dahlgren &
Lagerkvist (1995) classification scheme. The asteroids marked with a were
observed by Dotto et al. (2006), and their spectral slope values have been
recomputed in the 5500-8000 Å wavelength range.
Obj H D (km) S (%/103Å) T
Eurybates
3548 9.50∗∗ ∗72+4
-0.18±0.57 C
9818 11.00∗∗ 42+6
2.12±0.72 P
13862 11.10∗∗ 40+6
1.59±0.70 C
18060 11.10∗∗ 40+6
2.86±0.60 P
24380 11.20∗∗ 38+6
0.34±0.65 C
24420 11.50∗∗ 33+5
1.65±0.70 C
24426 12.50∗∗ 21+3
4.64±0.80 P
28958 12.10∗∗ 25+4
-0.04±0.80 C
39285 12.90∗∗ 17+3
0.25±0.69 C
43212 12.30∗∗ 23+4
1.19±0.78 C
53469 11.80∗∗ 29+4
0.17±0.80 C
65150 12.90∗∗ 17+3
4.14±0.70 P
65225 12.80∗∗ 18+3
0.97±0.85 C
1996RD29 13.06∗∗ 16+3
2.76±0.89 P
2000AT44 12.16∗∗ 24+3
-0.53±0.83 C
2002CT22 12.04∗∗ 26+4
2.76±0.73 P
2002EN68 12.30∗∗ 23+3
3.60±0.98 P
1986 WD
4035 9.72 ∗68+5
9.78±0.61 D
a4035 9.30∗∗ ∗68+5
15.19±0.61 D
6545 10.42 55+8
11.32±0.63 D
a6545 10.00∗∗ 66+10
9.88±0.56 D
11351 10.88 44+7
10.26±0.67 D
a11351 10.50∗∗ 53+8
10.44±0.61 D
14707 11.25 38+6
−9.4 -1.06±1.00 C
24233 11.58 33+5
−8.0 6.37±0.67 DP
24341 11.99 27+4
-0.26±0.71 C
1986 TS6
12917 11.61 32+5
10.98±0.68 D
12921 11.12 40+6
4.63±0.75 P
a12921 10.70∗∗ 48+7
3.74±1.00 P
13463 11.27 37+6
4.37±0.65 P
15535 10.70 48+7
10.67±0.65 D
20738 11.67 31+5
8.84±0.70 D
24390 11.80 29+5
9.53±0.62 D
Table 6: Results of the statistical analysis on the spectral slope distribution as a function of the diameters. For each
test bin, the average slope and the dispersion are listed; the size of the sample is reported in parenthesis. For each pair
of subsamples, the probability that both are randomly extracted from the same global sample is listed, as estimated
by the t-, f- and ks-test, respectively. Low probability indicates significant differences between the subsamples.
Diameter range 0–25 km 25–50 km 50–75 km 75–100 km >100 km
S average±σ 7.17±4.79 (22) 6.92±4.69 (48) 8.91±4.68 (26) 6.74±5.85 (21) 7.87±2.88 (21)
(%/103Å)
0–25 0.842 0.876 0.579 0.213 0.903 0.575 0.792 0.370 0.775 0.551 0.017 0.494
25–50 0.088 0.985 0.150 0.897 0.216 0.519 0.286 0.011 0.275
50–75 0.176 0.289 0.469 0.344 0.019 0.440
75–100 0.442 0.001 0.469
Table 7: Mean color indices and spectral slope of different classes of minor bodies of the outer Solar System. For each
class the number of objects considered is also listed.
Color Plutinos Cubewanos Centaurs Scattered Comets Trojans
B-V 36 87 29 33 2 74
0.895± 0.190 0.973± 0.174 0.886± 0.213 0.875± 0.159 0.795± 0.035 0.777± 0.091
V-R 38 96 30 34 19 80
0.568± 0.106 0.622± 0.126 0.573± 0.127 0.553± 0.132 0.441± 0.122 0.445± 0.048
V-I 34 64 25 25 7 80
1.095± 0.201 1.181± 0.237 1.104± 0.245 1.070± 0.220 0.935± 0.141 0.861± 0.090
V-J 10 14 11 8 1 12
2.151± 0.302 1.750± 0.456 1.904± 0.480 2.041± 0.391 1.630± 0.000 1.551± 0.120
V-H 3 7 11 4 1 12
2.698± 0.083 2.173± 0.796 2.388± 0.439 2.605± 0.335 1.990± 0.000 1.986± 0.177
V-K 2 5 9 2 1 12
2.763± 0.000 2.204± 1.020 2.412± 0.396 2.730± 0.099 2.130± 0.000 2.125± 0.206
R-I 34 64 25 26 8 80
0.536± 0.135 0.586± 0.148 0.548± 0.150 0.517± 0.102 0.451± 0.059 0.416± 0.057
J-H 11 17 21 11 1 12
0.403± 0.292 0.370± 0.297 0.396± 0.112 0.348± 0.127 0.360± 0.000 0.434± 0.064
H-K 10 16 20 10 1 12
-0.034± 0.171 0.084± 0.231 0.090± 0.142 0.091± 0.136 0.140± 0.000 0.139± 0.041
Slope 38 91 30 34 8 80
(%/103Å) 19.852± 10.944 25.603± 13.234 20.601± 13.323 18.365± 12.141 10.722± 6.634 7.241± 3.909
Table 8: Statistical tests performed to compare the color and slope distributions of different classes of minor bodies
(Plt= Plutinos, Resonant TNOs; QB1= Cubiwanos, Classical TNOs; Cent= Centaurs; Scat= scattered TNOs; Com=
Short Period Comet nuclei) with those of Trojans. The first five columns consider all the Trojans, the second five only
the Eurybates family, the third five only the non-Eurybates family Trojans. For each color, the first line shows the
number of objects used for the comparison (2nd is the number of Trojans), and the second line reports the probability
resulting from the test. A very low value indicates that the two compared distributions are not statistically compatible.
Probabilities are in boldface when the size of the samples is large enough for the value to be meaningful.
f-test
Color All Trojans Only Eurybates Only NON-Eurybates
Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com
B-V 36 74 83 74 29 74 33 74 2 74 36 14 83 14 29 14 33 14 2 14 36 60 83 60 29 60 33 60 2 60
0.000 0.000 0.000 0.000 0.600 0.001 0.001 0.000 0.005 0.722 0.000 0.000 0.000 0.000 0.598
V-R 38 80 92 80 30 80 34 80 19 80 38 17 92 17 30 17 34 17 19 17 38 63 92 63 30 63 34 63 19 63
0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000
R-I 34 80 62 80 25 80 26 80 8 80 34 17 62 17 25 17 26 17 8 17 34 63 62 63 25 63 26 63 8 63
0.000 0.000 0.000 0.000 0.773 0.000 0.000 0.000 0.001 0.335 0.000 0.000 0.000 0.000 0.185
Slope 38 80 87 80 30 80 34 80 8 80 38 17 87 17 30 17 34 17 8 17 38 63 87 63 30 63 34 63 8 63
0.000 0.000 0.000 0.000 0.020 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000
t-test
Color All Trojans Only Eurybates Only NON-Eurybates
Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com
B-V 36 74 83 74 29 74 33 74 2 74 36 14 83 14 29 14 33 14 2 14 36 60 83 60 29 60 33 60 2 60
0.001 0.000 0.012 0.002 0.608 0.000 0.000 0.001 0.000 0.139 0.003 0.000 0.025 0.006 0.858
V-R 38 80 92 80 30 80 34 80 19 80 38 17 92 17 30 17 34 17 19 17 38 63 92 63 30 63 34 63 19 63
0.000 0.000 0.000 0.000 0.916 0.000 0.000 0.000 0.000 0.083 0.000 0.000 0.000 0.000 0.532
R-I 34 80 62 80 25 80 26 80 8 80 34 17 62 17 25 17 26 17 8 17 34 63 62 63 25 63 26 63 8 63
0.000 0.000 0.000 0.000 0.154 0.000 0.000 0.000 0.000 0.001 0.000 0.000 0.001 0.000 0.502
Slope 38 80 87 80 30 80 34 80 8 80 38 17 87 17 30 17 34 17 8 17 38 63 87 63 30 63 34 63 8 63
0.000 0.000 0.000 0.000 0.185 0.000 0.000 0.000 0.000 0.008 0.000 0.000 0.000 0.000 0.404
KS-test
Color All Trojans Only Eurybates Only NON-Eurybates
Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com Plt QB1 Cent Scat Com
B-V 36 74 83 74 29 74 33 74 2 74 36 14 83 14 29 14 33 14 2 14 36 60 83 60 29 60 33 60 2 60
0.001 0.000 0.001 0.004 0.330 0.002 0.000 0.035 0.000 0.065 0.003 0.000 0.002 0.047 0.468
V-R 38 80 92 80 30 80 34 80 19 80 38 17 92 17 30 17 34 17 19 17 38 63 92 63 30 63 34 63 19 63
0.000 0.000 0.000 0.000 0.040 0.000 0.000 0.000 0.000 0.008 0.000 0.000 0.000 0.000 0.056
R-I 34 80 62 80 25 80 26 80 8 80 34 17 62 17 25 17 26 17 8 17 34 63 62 63 25 63 26 63 8 63
0.000 0.000 0.000 0.000 0.201 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.000 0.587
Slope 38 80 87 80 30 80 34 80 8 80 38 17 87 17 30 17 34 17 8 17 38 63 87 63 30 63 34 63 8 63
0.000 0.000 0.000 0.000 0.088 0.000 0.000 0.000 0.000 0.002 0.000 0.000 0.000 0.000 0.211
Figure captions
Fig. 1 - Reflectance spectra of 5 Anchises family members (L5 swarm).
The photometric color indices are also converted to relative reflectance and
overplotted on each spectrum. Spectra and photometry are shifted by 0.5 in
reflectance for clarity.
Fig. 2 - Reflectance spectra of 6 Misenus family members (L5 swarm).
The photometric color indices are also converted to relative reflectance and
overplotted on each spectrum. Spectra and photometry are shifted by 0.5 in
reflectance for clarity.
Fig. 3 - Reflectance spectra of 5 Panthoos family members (L5 swarm).
The photometric color indices are also converted to relative reflectance and
overplotted on each spectrum. Spectra and photometry are shifted by 0.5 in
reflectance for clarity. For asteroid 30698, the B-V color is missing as a B
filter measurement was not available.
Fig. 4 - Reflectance spectra of 2 Cloantus, 3 Phereclos and 2 Sarpedon
family members (L5 swarm). The photometric color indices are also con-
verted to relative reflectance and overplotted on each spectrum. Spectra and
photometry are shifted by 1.0 in reflectance for clarity.
Fig. 5 - Reflectance spectra of the 17 Eurybates family members (L4
swarm). Spectra are shifted by 0.5 in reflectance for clarity.
Fig. 6 - Reflectance spectra of the 6 1986 WD family members and 12921,
which is a member of the 1986 TS6 family (all belonging to the L4 swarm).
Spectra are shifted by 1.0 in reflectance for clarity.
Fig. 7 - Plot of the spectral slope versus the estimated diameter for the
families observed in the L5 swarm.
Fig. 8 - Plot of the spectral slope versus the estimated diameter for the
families observed in the L4 swarm.
Fig. 9 - Plot of the observed spectral slopes versus the estimated diameter
for the whole population of Jupiter Trojans investigated by us and available
from the literature. The errors on slopes and diameters are not plotted to
avoid confusion.
Fig. 10 - Histogram of L5 Trojans taxonomic classes.
Fig. 11 - Histogram of L4 Trojans taxonomic classes (Neg indicates ob-
jects with negative spectral slope).
Fig. 12 - Color distributions as functions of the absolute magnitude
M(1, 1), the inclination i [degrees], the orbital semi-major axis a [AU], the
perihelion distance q [AU], the eccentricity e, and the orbital energy E (see
text for definition). We include all the available colors for distant minor bod-
ies (TNOs, Centaurs, and cometary nuclei, see Hainaut & Delsanti 2002).
The Plutinos (resonant TNOs) are red filled triangles, Cubiwanos (classical
TNOs) are pink filled circles, Centaurs are green open triangles, Scattered
TNOs are blue open circles, and Trojans are cyan filled triangles.
Fig. 13 - V −R versus R−I color-color diagram for the observed Trojans
and all distant minor bodies available in the updated Hainaut & Delsanti
(2002) database. The solid symbols are for the Trojans (square for Eurby-
bates, triangles for others). The open symbols are used as following: tri-
angles for Plutinos, circles for Cubiwanos, squares for Centaurs, pentagons
for Scattered, and starry square for Comets. The continuous line represents
the ”reddening line”, that is the locus of objects with a linear reflectivity
spectrum. The star symbol represents the Sun.
Fig. 14 - Cumulative function and histograms of the B − V and V − R
color distributions and of the spectral slope for all the considered classes of
objects. The dotted line marks the solar colors.
4000 5000 6000 7000 8000 9000
Figure 1:
4000 5000 6000 7000 8000 9000
Figure 2:
4000 5000 6000 7000 8000 9000
Figure 3:
Figure 4:
Figure 5:
Figure 6:
Figure 7:
Figure 8:
Figure 9:
Figure 10:
Figure 11:
M(1,1) E
a [AU] q [AU]
Figure 12:
Figure 13:
Figure 14:
Introduction
Observations and data reduction
Results
Dynamical families: L5 swarm
Anchises
Misenus
Panthoos
Cloantus
Phereclos
Sarpedon
Dynamical families: L4 swarm
Eurybates
1986 WD
1986 TS6
Discussion
Size vs spectral slope distribution:Individual families
Size vs slope distribution: The Trojan population as a whole
Spectral slopes and L4/L5 Clouds
Orbital Elements
Comparison with other outer Solar System minor bodied
Introduction and methods
Results
Conclusions
|
0704.0351 | FIRST-based survey of Compact Steep Spectrum sources, V.
Milliarcsecond-scale morphology of CSS objects | Astronomy & Astrophysics manuscript no. 6364 c© ESO 2021
June 8, 2021
FIRST-based survey of compact steep spectrum sources
V. Milliarcsecond-scale morphology of CSS objects
M. Kunert-Bajraszewska1 and A. Marecki1
Toruń Centre for Astronomy, N. Copernicus University, 87-100 Toruń, Poland
Received 8 September 2006; Accepted 7 March 2007
ABSTRACT
Aims. Multifrequency VLBA observations of the final group of ten objects in a sample of FIRST-based compact steep spectrum
(CSS) sources are presented. The sample was selected to investigate whether objects of this kind could be relics of radio−loud AGNs
switched off at very early stages of their evolution or possibly to indicate intermittent activity.
Methods. Initial observations were made using MERLIN at 5 GHz. The sources have now been observed with the VLBA at 1.7, 5 and
8.4 GHz in a snapshot mode with phase-referencing. The resulting maps are presented along with unpublished 8.4-GHz VLA images
of five sources.
Results. Some of the sources discussed here show a complex radio morphology and therefore a complicated past that, in some cases,
might indicate intermittent activity. One of the sources studied – 1045+352 – is known as a powerful radio and infrared-luminous
broad absorption line (BAL) quasar. It is a young CSS object whose asymmetric two-sided morphology on a scale of several hundred
parsecs, extending in two different directions, may suggest intermittent activity. The young age and compact structure of 1045+352
is consistent with the evolution scenario of BAL quasars. It has also been confirmed that the submillimetre flux of 1045+352 can be
seriously contaminated by synchrotron emission.
Key words. galaxies: active, galaxies: evolution, quasars: absorption lines
1. Introduction
Following early hypotheses (Phillips & Mutel, 1982; Carvalho,
1985) suggesting that the gigahertz-peaked spectrum (GPS)
and compact steep spectrum (CSS) could be young objects,
Readhead et al. (1996) proposed an evolutionary scheme uni-
fying three classes of radio-loud AGNs (RLAGNs): symmet-
ric GPS objects – CSOs (compact symmetric objects); sym-
metric CSS objects – MSOs (medium-sized symmetric objects)
and large symmetric objects (LSOs). In this scheme GPS/CSO
sources with linear sizes less than 1 kpc1 would evolve into
CSS/MSOs with subgalactic sizes (<20 kpc) and these in turn
would eventually become LSOs during their lifetimes. Two
pieces of evidence definitely point towards GPS/CSS sources
being young objects: lobe proper motions (up to 0.3c) giving
kinematic ages as low as ∼103 years for CSOs (Owsianik et al.,
1998; Giroletti et al., 2003; Polatidis & Conway, 2003) and ra-
diative ages typically ∼105 years for MSOs (Murgia et al.,
1999). Although these AGNs are small-scale objects, in some
cases CSO/GPS sources are associated with much larger ra-
dio structures that extend out to many kiloparsecs. In these
cases, it has been suggested that the CSO/GPS stage rep-
resents a period of renewed activity in the life cycle of
the AGN (Stanghellini et al., 2005, and references therein).
Reynolds & Begelman (1997) have also proposed a model in
which extragalactic radio sources are intermittent on timescales
Send offprint requests to: M. Kunert-Bajraszewska
e-mail: [email protected]
1 For consistency with earlier papers in this field, the following
cosmological parameters have been adopted throughout this paper:
H0=100 km s
−1 Mpc−1 and q0=0.5. Throughout this paper, the spectral
index is defined such that S ∝ να.
of ∼104–105 years. Following the above scenarios and also an
earlier suggestion by Readhead et al. (1994) and O’Dea & Baum
(1997) that there exists a large population of compact, short-
lived objects, Marecki et al. (2003, 2006) concluded that the evo-
lutionary track proposed by Readhead et al. (1996) is only one
of many possible tracks. A lack of stable fuelling from the black
hole can inhibit the growth of a radio source, and consequently
it will never reach the LSO stage, at least in a given phase of its
activity.
Observational support for the above ideas has been pro-
vided by Gugliucci et al. (2005). They calculated the kinematic
ages for a sample of CSOs with well-identified hotspots. It ap-
pears that the kinematic age distribution drops sharply above
∼500 years, suggesting that in many CSOs activity may cease
early. It is, therefore, possible that only some of them evolve
any further. Our observations have shown that young, fading
compact sources do indeed exist (Kunert-Bajraszewska et al.,
2005; Marecki et al., 2006; Kunert-Bajraszewska et al., 2006,
hereafter Papers II, III, and IV, respectively). A double source,
0809+404, described in Paper IV is our best example of a very
compact – i.e. very young – fader. The VLBA multifrequency
observations have shown it to have a diffuse, amorphous struc-
ture, devoid of a dominant core and hotspots. Giroletti et al.
(2005) have analysed the properties of a sample of small-size
sources and found a very good example of a kiloparsec-scale
fader (1855+37). It is to be noted that re-ignition of activity in
compact radio sources is not ruled out. In this paper – the fifth
and the last of the series – VLBA observations of 10 CSS and
CSO sources that are potential candidates for compact faders
or objects with intermittent activity are presented. One of these
sources, 1045+352, is of particular interest not only because it
http://arxiv.org/abs/0704.0351v2
2 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
has a puzzling radio structure, but it also appears to be a broad
absorption line (BAL) quasar.
As their name somewhat suggests, BAL quasars have very
broad, blue-shifted absorption lines arising from high-ionization
transitions such as C IV, Si IV, N V, etc. (e.g C IV 1549Å). They
constitute ∼10% of the optically selected radio-quiet quasars
with the absorption arising from gas outflow at velocities up to
∼0.2 c (Hewett & Foltz, 2003). In fact, BAL quasars have been
divided into two categories, as 10% of them also show absorp-
tion troughs in low-ionization lines such as Mg II 2800Å. This
group has been designated as LoBAL quasars and the others as
HiBAL ones. The high ionization level and continuous absorp-
tion over a wide velocity range is hard to reconcile with absorp-
tion by individual clouds. Rather, they indicate that BAL regions
exist in both BAL and non-BAL quasars and evidence, accumu-
lated from optically selected BAL quasars, indicates an orienta-
tion hypothesis to explain their nature. It would appear that BAL
quasars are normal quasars seen along a particular line of sight,
e.g. a line of sight skimming the edge of the accretion disk or
torus (Weymann et al., 1991; Elvis, 2000). Murray et al. (1995)
have proposed a model in which the line of sight to a BAL quasar
intersects an outflow or wind that is not entirely radial, e.g. an
outflow that initially emerges perpendicular to the accretion disk
and is then accelerated radially.
For quite a long time it was believed that BAL quasars were
never radio-loud. This view was challenged by Becker et al.
(1997), who discovered the first radio-loud BAL quasar when
using the VLA FIRST survey to select quasar candidates.
Five radio-loud BAL quasars were then identified in NVSS
by Brotherton et al. (1998). Since then, the number of radio-
loud BAL QSOs has increased considerably (Becker et al., 2000;
Menou et al., 2001), following identification of new quasar can-
didates selected from the FIRST survey. Most of the BAL
quasars in the Becker et al. (2000) sample tended to be com-
pact at radio frequencies with either a flat or steep radio spec-
trum. Those with steep spectra could be related to GPS and CSS
sources. A variety of their spectral indices also suggested a wide
range of orientations, contrary to the interpretation favoured
from optically selected quasars. Moreover, Becker et al. (2000)
indicated that the frequency of BAL quasars in their sample was
significantly greater (factor ∼2) than inferred from optically se-
lected samples and that the frequency of BAL quasars appeared
to show a complex dependence on radio loudness.
The radio morphology of BAL quasars is important because
it can indicate inclination in BALs, and therefore yields a di-
rect test of the orientation model. However, information about
the radio structure of BAL quasars is still very limited. Prior
to 2006, only three BAL quasars, FIRST J101614.3+520916
(Gregg et al., 2000), PKS 1004+13 (Wills et al., 1999), and
LBQS 1138−0126 (Brotherton et al., 2002) were known to have
a double-lobed FR II radio morphology on kiloparsec scales,
although this interpretation was doubtful for PKS 1004+13
(Gopal-Krishna & Wiita, 2000). Recently, the population of
FR II-BAL quasars has increased to ten objects (excluding PKS
1004+13) following the discoveries of Gregg et al. (2006) and
Zhou et al. (2006), although some of these still require confir-
mation. Their symmetric structures indicate an “edge-on” ori-
entation, which in turn supports an alternative hypothesis de-
scribed as “unification by time”, with BAL quasars charac-
terised as young or recently refuelled quasars (Becker et al.,
2000; Gregg et al., 2000). There has been only one attempt (at
1.6 GHz with the EVN) to image radio structures of the smallest
(and possibly the youngest) BAL quasars (Jiang & Wang, 2003)
from the Becker et al. (2000) sample. This paper presents high
frequency VLBA images of another very compact BAL quasar
— 1045+352, which makes it the BAL quasar with the best
known radio structure to date.
2. The observations and data reduction
The five papers of this series are concerned with a sample
of 60 candidates selected from the VLA FIRST catalogue
(White et al., 1997)2 which could be weak CSS sources. The
sample selection criteria have been given in Kunert et al. (2002)
(hereafter Paper I). All the sources were initially observed with
MERLIN at 5 GHz and the results of these observations led to
the selection of several groups of objects for further study with
MERLIN and the VLA (Paper II), as well as the VLBA and
the EVN (Papers III and IV). The last of those groups contains
10 sources that, because of their structures (very faint “haloes”
or possible core-jet structures), were not included in the other
groups as they were less likely to be candidates for faders.
However, to complete the investigation of the primary sample,
1.7, 5, and 8.4-GHz VLBA observations of 10 sources listed in
Table 1 together with their basic properties, were carried out on
13 November 2004 in a snapshot mode with phase-referencing.3
Each target source scan was interleaved with a scan on a phase
reference source and the total cycle time (target and phase-
reference) was ∼9 minutes including telescope drive times, with
∼7 minutes actually on the target source per cycle. The cycles
for a given target-calibrator pair were grouped and rotated round
the three frequencies, although the source 1059+351 was only
observed at 1.7 GHz with the VLBA because of its very low flux
density as measured at 5 GHz by MERLIN (13 mJy).
The whole data reduction process was carried out using
standard AIPS procedures but, in addition to this, corrections
for Earth orientation parameter (EOP) errors introduced by the
VLBA correlator also had to be made. For each target source
and at each frequency, the corresponding phase-reference source
was mapped, and the phase errors so determined were applied to
the target sources, which were then mapped using a few cycles
of phase self-calibration and imaging. For some of the sources
a final amplitude self-calibration was also applied. IMAGR was
used to produce the final “naturally weighted”, total intensity im-
ages shown in Figs. 1 to 10. Three of the ten sources (1056+316,
1302+356, 1627+289) were not detected in the 8.4-GHz VLBA
observations, and 1425+287 has not been detected in any VLBA
observations. Flux densities of the principal components of the
sources were measured using the AIPS task JMFIT and are listed
in Table 3.
In addition to the observations described above, unpub-
lished 8.4-GHz VLA observations of five sources – 1056+316,
1126+293, 1425+287, 1627+289, 1302+356 – made in A-conf.
by Glen Langston (first four objects) and Patnaik et al. (1992)
have been included (Figs. 3, 5, 9, 10, and 7, respectively).
It was realised that because of poor u-v coverage at the higher
frequencies, some flux density could be missing and the resul-
tant spectral index maps were not considered to be reliable. Any
calculation of spectral indices from the flux densities quoted in
Table 3 should also be treated only as coarse approximations.
For 1045+352, 30-GHz continuum observations using the
Toruń 32-m radio telescope and a prototype (two-element
2 Official website: http://sundog.stsci.edu
3 Including this paper, the results of the observations of 46 sources
out of 60 candidates from the primary sample have been published. The
observations of 14 objects failed for different reasons.
http://sundog.stsci.edu
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 3
Table 1. Basic parameters of target sources
Source RA Dec ID mR z S 1.4 GHz logP1.4GHz S 4.85 GHz α
4.85GHz
1.4GHz LAS LLS
Name h m s ◦ ′ ′′ mJy W Hz−1 mJy ′′ h−1 kpc
(1) (2) (3) (4) (5) (6) (7) (8) (9) (10) (11) (12)
1045+352 10 48 34.247 34 57 24.99 Q 20.86 1.604 1051 27.65 439 −0.70 ∼0.50 2.1
1049+384 10 52 11.797 38 11 43.83 G 20.76 1.018 712 27.04 205 −1.00 0.14 0.6
1056+316 10 59 43.236 31 24 20.59 G 21.10 0.307∗ 459 25.72 209 −0.63 0.50 1.4
1059+351 11 02 08.686 34 55 10.74 G 19.50 0.594∗ 702 26.52 252 −0.82 3.03 11.5
1126+293 11 29 21.738 29 05 06.40 EF — — 729 — 213 −0.99 0.79 —
1132+374 11 35 05.927 37 08 40.80 G — 2.880 638 28.00 218 −0.86 ∼0.30 1.1
1302+356 13 04 34.477 35 23 33.93 EF — — 483 — 185 −0.77 ∼0.20 —
1407+369 14 09 09.528 36 42 08.06 q 21.51 0.996∗ 538 26.89 216 −0.73 ∼0.25 1.1
1425+287 14 27 40.281 28 33 25.78 EF — — 859 — 198 −1.18 0.75 —
1627+289 16 29 12.290 28 51 34.25 EF — — 526 — 162 −0.95 ∼0.65 —
Description of the columns: (1) source name in the IAU format; (2) source right ascension (J2000) extracted from FIRST; (3) source declination
(J2000) extracted from FIRST; (4) optical identification: G - galaxy, Q - quasar, EF - empty field, q - star-like object, i.e. unconfirmed QSO; (5)
red magnitude extracted from SDSS/DR5; (6) redshift; (7) total flux density at 1.4 GHz extracted from FIRST; (8) log of the radio luminosity
at 1.4 GHz; (9) total flux density at 4.85 GHz extracted from GB6; (10) spectral index between 1.4 and 4.85 GHz calculated using flux densities
in columns (7) and (9); (11) largest angular size (LAS) measured in the 5-GHz MERLIN image – in most cases, as a separation between the
outermost component peaks, otherwise “∼” means measured in the image contour plot; (12) largest linear size (LLS).
∗ photometric redshift extracted from SDSS/DR5
receiver) of the One-Centimeter Receiver Array (OCRA-p,
Lowe et al., 2005) have also been made. The recorded output
from the receiver was the difference between the signals from
two closely-spaced horns effectively separated in azimuth so
that atmospheric variations were mostly cancelled out. The ob-
serving technique was such that the respective two beams were
pointed at the source alternately with a switching cycle of ∼50
seconds for a period of ∼6 minutes, thus measuring the source
flux density relative to the sky background on either side of the
source. The telescope pointing was determined from azimuth
and elevation scans across the point source Mrk 421. The pri-
mary flux density calibrator that was used was the planetary neb-
ula NGC 7027, which has an effective radio angular size of ∼8
arcseconds (Bryce et al., 1997) and for which a correction of the
flux density scale had to be made. However, as NGC 7027 was at
some distance from the target source, the point source 1144+402
was used as a secondary flux density calibrator. Corrections for
the effects of the atmosphere were determined from system tem-
perature measurements at zenith distances of 0◦ and 60◦.
3. Comments on individual sources
1045+352. The MERLIN and VLBA maps (Fig. 1) show this
source to be extended in both the NE/SW and NW/SE directions.
The central compact feature visible in all the maps is probably a
radio core with a steep spectrum. The VLBA image at 1.7 GHz
shows two symmetric protrusions – possibly jets – straddling the
core in a NE/SW direction, the SW emission being weaker than
in the NE. This structure is aligned with the NE/SW emission
visible in the 5-GHz MERLIN image, but the more extended dif-
fuse emission has been resolved out in the VLBA images. The
5-GHz VLBA image shows a core and a one-sided jet pointing
to the East. Some compact features in a NE direction are also
visible. The radio structure in the 8.4-GHz VLBA image is sim-
ilar to that at 5 GHz: an extended radio core and a jet pointing in
an easterly direction.
The observed radio morphology of 1045+352 could indicate
a restart of activity with the NE/SW radio emission being the
first phase of activity, now fading away, and the extension in the
NW/SE direction being a signature of the current active phase.
However, the above is only one of a number of possible interpre-
tations of the structure of 1045+352 – see further discussion in
Sect. 4.
According to Sloan Digital Sky Survey/Data Release 5
(SDSS/DR5), 1045+352 is a galaxy at RA= 10h48m34.s242,
Dec=+34◦57′24.′′95, which is marked with a cross in the
MERLIN map but the spectral observations carried out by
Willott et al. (2002) have shown 1045+352 to be a quasar with
a redshift of z = 1.604. It has been also classified as a HiBAL
object based upon the observed very broad C IV absorption, and
it is a very luminous submillimetre object with detections at both
850µm and 450µm (Willott et al., 2002).
The total flux of 1045+352 at 30 GHz measured by us using
OCRA-p is S 30GHz=69 mJy±7 mJy, which gives a steep spectral
index α = −1.01 between 4.85 GHz and 30 GHz.
1049+384. The 5-GHz MERLIN image (Fig. 2) shows it as
a triple core-jet structure with the brightest component re-
solved into a double structure extended in a NW/SE direc-
tion in the high resolution VLBA observations. The 1.7-GHz
VLBA image shows four radio components (in agreement with
Dallacasa et al., 2002), whereas the 5-GHz and 8.4-GHz VLBA
maps show only three components. However, the 5-GHz VLBA
image published by Orienti et al. (2004) shows all four compo-
nents, and they suggest that the two western components and the
two eastern ones are two independent radio sources. As pointed
by Orienti et al. (2004), it is difficult to classify the object, al-
though the idea that 1049+384 consists of two separate com-
pact, double sources is not very plausible because of the very
small separation, ∼ 0.09′′ (0.4 kpc), between these two poten-
tial objects. Although the spectral index calculations are very
uncertain, it is suggested that one of the eastern components
at RA= 10h52m11.s797, Dec=+38◦11′44.′′027 is a radio core (in
agreement with Orienti et al., 2004) from which jets emerge al-
ternately in opposite directions.
1049+384 is a galaxy with a redshift z = 1.018
(Riley & Warner, 1994), but according to Allington-Smith et al.
(1988) the optical spectrum of 1049+384 shows interme-
diate properties between a galaxy and a quasar. The opti-
4 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
1045+352 4994.000 MHz
peak flux density=230.21 mJy/beam, beam size=56 x 41 mas
first contour level=0.12 mJy/beam
RIGHT ASCENSION (J2000)
10 48 34.30 34.28 34.26 34.24 34.22 34.20
34 57 25.8
1045+352 1667.474 MHz
peak flux density=118.71 mJy/beam, beam size=13.1 x 8.2 mas
first contour level=0.80 mJy/beam
RIGHT ASCENSION (J2000)
10 48 34.256 34.254 34.252 34.250 34.248 34.246 34.244 34.242 34.240
34 57 25.14
25.12
25.10
25.08
25.06
25.04
25.02
25.00
24.98
24.96
24.94
1045+352 4987.474 MHz
peak flux density=13.64 mJy/beam, beam size=4.7 x 2.4 mas
first contour level=0.14 mJy/beam
RIGHT ASCENSION (J2000)
10 48 34.254 34.252 34.250 34.248 34.246 34.244
34 57 25.12
25.10
25.08
25.06
25.04
25.02
25.00
24.98
1045+352 8421.474 MHz
peak flux density=4.03 mJy/beam, beam size=2.7 x 1.5 mas
first contour level=0.14 mJy/beam
RIGHT ASCENSION (J2000)
10 48 34.251 34.250 34.249 34.248 34.247 34.246 34.245
34 57 25.08
25.07
25.06
25.05
25.04
25.03
25.02
25.01
Fig. 1. The MERLIN 5-GHz (upper left) and VLBA 1.7, 5, and 8.4-GHz maps of 1045+352. Contours increase by a factor 2, and
the first contour level corresponds to ≈ 3σ. A cross indicates the position of an optical object found using the SDSS/DR5.
cal object was included in SDSS/DR5 (RA= 10h52m11.s802,
Dec=+38◦11′44.′′00) and is marked in all maps with a cross.
1056+316. The 8.4-GHz VLA image (Fig. 3) shows this source
to have a double structure that, in the 5-GHz MERLIN image,
has been resolved into a radio core and probably a hotspot in
a NW radio lobe. Both components are visible in the 1.7-GHz
VLBA image, but neither has been detected in the higher fre-
quency VLBA images. The two weak features on either side of
the NW component in the 1.7-GHz VLBA image may be the
remains of extended emission that has been resolved out.
The optical counterpart of 1056+316 was included
in SDSS/DR5 (RA= 10h59m43.s145, Dec=+31◦24′23.′′31), to-
gether with a photometric redshift (Table 1). Its position is
marked with a cross in 8.4-GHz VLA map.
1059+351. The 5-GHz MERLIN map (Fig. 4) shows a bright
component that is probably a radio core, on almost opposite
sides of which is emission from compact features (hotspots)
within the two radio lobes. This structure agrees with the 1.4-
GHz VLA observations presented by Gregorini et al. (1988) and
Machalski & Condon (1983). Their images clearly show an S-
shaped morphology of 1059+351 with two very diffuse compo-
nents, the brighter one resolved into a double structure in 5-GHz
VLA observations (Machalski, 1998). One of these two com-
ponents is the NW hotspot visible in the 5-GHz MERLIN map,
and the second is probably a radio core visible in both the 5-GHz
MERLIN and 1.7-GHz VLBA images.
The optical counterpart of 1059+351 was included
in SDSS/DR5 (RA= 11h02m08.s727, Dec=+34◦55′08.′′79), to-
gether with a photometric redshift (Table 1). The position of the
optical object is marked with a cross in all maps and is well
correlated with the position of the radio core. Machalski (1998)
also measured a photometric redshift for 1059+351, which is
z = 0.37 and which differs from that in SDSS/DR5.
1126+293. The VLA 8.4-GHz and MERLIN 5-GHz maps
(Fig. 5) show three radio components, the brighter one proba-
bly being the core that was resolved into a core-jet structure in
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 5
1049+384 4994.000 MHz
peak flux density=116.27 mJy/beam, beam size=62 x 38 mas
first contour level=0.40 mJy/beam
RIGHT ASCENSION (J2000)
10 52 11.84 11.83 11.82 11.81 11.80 11.79 11.78 11.77 11.76
38 11 44.6
1049+384 1667.474 MHz
peak flux density=181.64 mJy/beam, beam size=11.6 x 8.2 mas
first contour level=0.09 mJy/beam
RIGHT ASCENSION (J2000)
10 52 11.810 11.805 11.800 11.795 11.790 11.785
38 11 44.14
44.12
44.10
44.08
44.06
44.04
44.02
44.00
43.98
43.96
43.94
1049+384 4987.474 MHz
peak flux density=21.07 mJy/beam, beam size=4.2 x 2.3 mas
first contour level=0.09 mJy/beam
RIGHT ASCENSION (J2000)
10 52 11.805 11.800 11.795 11.790
38 11 44.08
44.07
44.06
44.05
44.04
44.03
44.02
44.01
44.00
43.99
43.98
1049+384 8421.474 MHz
peak flux density=68.39 mJy/beam, beam size=2.6 x 1.2 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
10 52 11.800 11.798 11.796 11.794 11.792 11.790 11.788 11.786
38 11 44.06
44.05
44.04
44.03
44.02
44.01
44.00
43.99
Fig. 2. The MERLIN 5-GHz (upper left) map and VLBA 1.7, 5, and 8.4-GHz maps of 1049+384. Contours increase by a factor 2,
and the first contour level corresponds to ≈ 3σ. Crosses indicate the position of an optical object found using the SDSS/DR5 .
the 1.7-GHz VLBA image. The source was not detected in the 5
and 8.4-GHz VLBA observations.
1132+374. The 5-GHz MERLIN image shows (Fig. 6) a core-jet
structure that was resolved into a triple CSO object in the 1.7-
GHz VLBA image. The 5 and 8.4-GHz VLBA images show only
two components: a hotspot in the NE lobe and a radio core. This
source is identified with a very high redshift (z = 2.88) galaxy
(Eales & Rawlings, 1996).
1302+356. This source was observed with the VLA at 8.4 GHz
as a part of the JVAS survey (Patnaik et al., 1992). The result-
ing map shows a slightly extended EW object (Fig. 7). The 5-
GHz MERLIN image shows this to be a double source, and the
weak (∼10 mJy) eastern component could be part of a jet. The
bright component was resolved into a diffuse structure in the 1.7-
GHz VLBA image. The 5-GHz VLBA image shows only a sin-
gle component at the position of the maximum emission in the
1.7-GHz VLBA image, which is probably a radio core (Fig. 7).
There is no trace of this source in the 8.4-GHz VLBA image.
1407+369. The 5-GHz MERLIN image shows a core-jet struc-
ture in a NW direction that is resolved into a core and jet in
6 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
1056+316 8439.900 MHz
peak flux density=118.41 mJy/beam, beam size=270 x 257 mas
first contour level=0.06 mJy/beam
RIGHT ASCENSION (J2000)
10 59 43.45 43.35 43.25 43.15 43.05
31 24 23
1056+316 4994.000 MHz
peak flux density=80.83 mJy/beam, beam size=60 x 43 mas
first contour level=0.16 mJy/beam
RIGHT ASCENSION (J2000)
10 59 43.28 43.26 43.24 43.22 43.20
31 24 21.2
1056+316 1667.474 MHz
peak flux density=9.10 mJy/beam, beam size=13.9 x 5.5 mas
first contour level=0.30 mJy/beam
RIGHT ASCENSION (J2000)
10 59 43.265 43.255 43.245 43.235 43.225
31 24 20.6
Fig. 3. The VLA 8.4-GHz map, MERLIN 5-GHz map, and VLBA 1.7-GHz map of 1056+316. Contours increase by a factor 2, and
the first contour level corresponds to ≈ 3σ. A cross on the VLA map indicates the position of an optical object found using the
SDSS/DR5.
1059+351 4994.000 MHz
peak flux density=10.03 mJy/beam, beam size=89 x 69 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
11 02 08.85 08.80 08.75 08.70 08.65 08.60
34 55 10.0
1059+351 1667.474 MHz
peak flux density=8.07 mJy/beam, beam size=10.9 x 7.9 mas
first contour level=0.08 mJy/beam
RIGHT ASCENSION (J2000)
11 02 08.735 08.730 08.725 08.720
34 55 08.80
08.75
08.70
08.65
08.60
Fig. 4. The MERLIN 5-GHz map and VLBA 1.7-GHz map of 1059+351. Contours increase by a factor 2, and the first contour level
corresponds to ≈ 3σ. Crosses indicate the position of an optical object found using the SDSS/DR5.
all the VLBA maps (Fig. 8). The optical object was included
in SDSS/DR5 (RA= 14h09m09.s509, Dec=+36◦42′08.′′15) and is
marked with a cross in all maps. The redshift quoted in Table 1
is photometric.
1425+287. Both the VLA 8.4-GHz and MERLIN 5-GHz images
(Fig. 9) show a double structure for this source. The brighter
component seems to be a radio core, although this cannot be
confirmed because the source was not detected in the VLBA ob-
servations (Fig. 9).
1627+289. Both the VLA 8.4-GHz and MERLIN 5-GHz images
(Fig. 10) show this source to have a core-jet structure. The 1.7-
GHz VLBA image shows only the central extended feature that
was resolved into a core-jet structure in the 5-GHz VLBA image.
The source was not detected in the 8.4-GHz VLBA image.
4. Discussion
4.1. 1045+352 — a BAL quasar
1045+352 is a HiBAL quasar with a very reddened spectrum
showing a C IV broad absorption system (Willott et al., 2002).
Its projected linear size is only 2.1 kpc, which is consistent with
the observation of Becker et al. (2000) that, amongst radio loud
quasars, broad absorption lines are more commonly observed in
the smallest radio sources.
It is a very luminous submillimetre object, which together
with the template dust spectrum adopted by Willott et al. (2002),
indicates this source to be a hyperluminous infrared quasar, with
large amounts of dust in its host galaxy. Although 1045+352
is quite luminous at 151 MHz (2.88 Jy, Waldram et al., 1996),
which suggests the presence of some extended emission and
which, indeed, appears to be present in our MERLIN 5-GHz
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 7
1126+293 8439.900 MHz
peak flux density=65.97 mJy/beam, beam size=368 x 310 mas
first contour level=0.08 mJy/beam
RIGHT ASCENSION (J2000)
11 29 23.9 23.8 23.7 23.6 23.5 23.4
29 05 01
04 59
1126+293 4994.000 MHz
peak flux density=60.53 mJy/beam, beam size=62 x43 mas
first contour level=0.14 mJy/beam
RIGHT ASCENSION (J2000)
11 29 21.80 21.75 21.70 21.65
29 05 07.5
1126+293 1667.474 MHz
peak flux density=6.03 mJy/beam, beam size=14.0 x 4.2 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
11 29 21.762 21.760 21.758 21.756 21.754 21.752 21.750 21.748 21.746
29 05 06.50
06.48
06.46
06.44
06.42
06.40
06.38
06.36
06.34
06.32
06.30
Fig. 5. The VLA 8.4-GHz map, MERLIN 5-GHz map and VLBA 1.7-GHz map of 1126+293. Contours increase by a factor 2, and
the first contour level corresponds to ≈ 3σ.
image, the VLBA maps show the radio structure to be domi-
nated by jets and a core. The 30-GHz flux density of 1045+352
is also high, as would be expected from the VLBA structure.
Consequently, there could be synchrotron contamination of the
submillimetre flux. As shown by Blundell et al. (1999), either
the first-order or second-order polynomials can accurately pre-
dict the shape of the radio spectrum. Both models have been ap-
plied to the radio data of 1045+352 taken from the literature and
from this paper (Fig. 11), and show that a non-thermal compo-
nent could constitute at least ∼40% of the entire 850µm flux (the
parabolic fit). The linear fit agrees with calculations based upon
the 1.25 mm flux measured by Haas et al. (2006), who derived a
value of 94% for the non-thermal component part of the detected
850µm flux. It has to be noted here that the linear fit should be
treated as an upper limit for the synchrotron emission at submil-
limetre wavelengths, since the spectrum may steepen in the inter-
val between 30 GHz and the SCUBA wavebands. However, the
above can indicate values of infrared emission and dust mass of
1045+352 lower than estimated (Willott et al., 2002). This also
appears be consistent with the findings of Willott et al. (2003),
who have shown that there is no difference between the submil-
limetre luminosities of BAL and non-BAL quasars, which sug-
gest that a large dust mass is not required for quasars to show
BALs.
The radio luminosity at 1.4 GHz is high (Table 1), making
this source one of the most radio-luminous BAL quasars, with
a value similar to that of the first known radio-loud BAL QSO
with an FR II structure, FIRST J101614.3+520916 (Gregg et al.,
2000). Following Stocke et al. (1992), a radio-loudness param-
eter, R∗, defined as the K-corrected ratio of the 5-GHz radio
flux to 2500Å optical flux (Table 2) was calculated. For this,
a global radio spectral index, αradio = −0.8 and an optical spec-
tral index, αopt = −1.0, were assumed, and the SDSS g
′ mag-
nitude defined by Fukugita et al. (1996) was converted to the
Johnson-Morgan-Cousins B magnitude using the formula given
by Smith et al. (2002). Corrections were also made for intrin-
8 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
1132+374 4994.000 MHz
peak flux density=122.40 mJy/beam, beam size=58 x 44 mas
first contour level=0.18 mJy/beam
RIGHT ASCENSION (J2000)
11 35 05.98 05.96 05.94 05.92 05.90 05.88
37 08 41.6
1132+374 1667.474 MHz
peak flux density=38.64 mJy/beam, beam size=9.5 x 4.0 mas
first contour level=0.40 mJy/beam
RIGHT ASCENSION (J2000)
11 35 05.940 05.938 05.936 05.934 05.932 05.930 05.928 05.926
37 08 40.86
40.84
40.82
40.80
40.78
40.76
40.74
40.72
40.70
40.68
40.66
1132+374 4987.474 MHz
peak flux density=12.57 mJy/beam, beam size=3.1 x 1.2 mas
first contour level=0.14 mJy/beam
RIGHT ASCENSION (J2000)
11 35 05.936 05.934 05.932 05.930 05.928
37 08 40.82
40.80
40.78
40.76
40.74
40.72
40.70
1132+374 8421.474 MHz
peak flux density=8.74 mJy/beam, beam size=2.2 x 1.5 mas
first contour level=0.10 mJy/beam
RIGHT ASCENSION (J2000)
11 35 05.936 05.935 05.934 05.933 05.932 05.931 05.930
37 08 40.83
40.82
40.81
40.80
40.79
40.78
40.77
40.76
40.75
40.74
Fig. 6. The MERLIN 5-GHz (upper left) map and VLBA 1.7, 5, and 8.4-GHz maps of 1132+374. Contours increase by a factor 2,
and the first contour level corresponds to ≈ 3σ.
sic extinction (local to the quasar) calculated by Willott et al.
(2002), who assumed a Milky-Way extinction curve. Even af-
ter correction, log(R∗) > 1, which means that 1045+352 is still
radio-loud object. The angle between the jet axis and the line
of sight can be estimated using the core radio-to-optical lumi-
nosity ratio defined by Wills & Brotherton (1995) as log(RV ) =
log(Lcore) + 0.4MV − 13.69, where Lcore is a radio luminosity of
the core at 5-GHz rest frequency (the core flux density at 5 GHz
were taken from the VLBA image; see also Table 3), and MV
is the K-corrected absolute magnitude calculated using transfor-
mation equation V = g′−0.55(g′−r′)−0.03 (Smith et al., 2002).
From this, a value of ∼3.2 has been obtained for 1045+352, im-
plying an angle in the range θ ∼ 10◦ − 30◦ for the jet in the
observed asymmetric MERLIN 5-GHz radio morphology, and
can explain the high value of the radio-loudness parameter. An
assumption of θ = 20◦ yields the deprojected linear size of the
source of ∼ 6 kpc. As shown by White et al. (2006), BAL QSOs
are systematically brighter than non-BAL objects, which indi-
cates we are looking closer to the jet axis in quasars with BALs.
Based upon the small inclination angles of their BAL quasars,
Zhou et al. (2006) suggest that BAL features can be caused by
polar disk winds. Also, Saikia et al. (2001) and Jeyakumar et al.
(2005) found that the radio properties of CSS sources are con-
sistent with the unified scheme in which the axes of the quasars
are observed close to the line of sight. On the other hand, it
has been shown (Saikia et al., 2001; Jeyakumar et al., 2005) that
many CSS objects interact with an asymmetric medium in the
central regions of their host galaxies, and this can cause the ob-
served asymmetries. It is then likely that, also in the case of the
CSS quasar 1045+352, the environmental asymmetries might
play an important role. The jet power can be estimated from
the relationship between the radio luminosity and the jet power
given by Willott et al. (1999, Eq.(12)). However, because some
of the flux density of the 1045+352 can be beamed, the calcu-
lations have to be treated as an approximation. Assuming the
151-MHz flux density, which accounts for the extended emis-
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 9
1302+356 8452.400 MHz
peak flux density=109.96 mJy/beam, beam size=252 x 230 mas
first contour level=0.09 mJy/beam
RIGHT ASCENSION (J2000)
13 04 34.75 34.70 34.65 34.60 34.55 34.50 34.45 34.40 34.35 34.30
35 23 36
1302+356 4994.500 MHz
peak flux density=129.54 mJy/beam, beam size=62 x 39 mas
first contour level=0.18 mJy/beam
RIGHT ASCENSION (B1950)
13 02 13.86 13.84 13.82 13.80 13.78 13.76 13.74 13.72 13.70 13.68
35 39 38.5
1302+356 1667.474 MHz
peak flux density=19.62 mJy/beam, beam size=10.0 x 4.0 mas
first contour level=0.14 mJy/beam
RIGHT ASCENSION (J2000)
13 04 34.502 34.500 34.498 34.496 34.494 34.492 34.490 34.488 34.486
35 23 33.64
33.62
33.60
33.58
33.56
33.54
33.52
33.50
33.48
33.46
33.44
1302+356 4987.474 MHz
peak flux density=4.23 mJy/beam, beam size=3.8 x 1.5 mas
first contour level=0.07 mJy/beam
RIGHT ASCENSION (J2000)
13 04 34.498 34.497 34.496 34.495 34.494 34.493 34.492
35 23 33.57
33.56
33.55
33.54
33.53
33.52
33.51
33.50
33.49
Fig. 7. The VLA 8.4-GHz map (upper left), MERLIN 5-GHz map (upper right) and VLBA 1.7 and 5-GHz maps of 1302+356.
Contours increase by a factor 2, and the first contour level corresponds to ≈ 3σ.
sion and the radio emission from the jets, the jet kinetic power is
Q jet ∼ 10
44erg sec−1.
The projected linear size D of a radio quasar or radio
galaxy can be approximately related to the time, from the trig-
gering of activity, as the relationship between these variables
is only weakly dependent upon the radio luminosity. Using
the model of radio source evolution from Willott et al. (1999),
the age of 1045+352 was estimated to be ∼ 105 years (see
also Willott et al., 2002; Rawlings et al., 2004). For the calcu-
lations we assumed: θ = 20◦, β = 1.5, c1 = 2.3, n100 =
3000 e− m−3, a0 = 100 kpc (see Willott et al., 1999, for defini-
tions). Both the MERLIN and VLBA high frequency images
have revealed that two cycles of activity may have occurred dur-
ing these ∼ 105 years. The extended NE/SW emission is prob-
ably the remnant of the first phase of activity, which has been
very recently replaced by a new phase of activity pointing in a
NW/SE direction. It has been shown by Stanghellini et al. (2005)
that the extended emission observed for small-scale objects can
be the remnants of an earlier period of activity in these sources.
In the case of 1045+352, renewal of activity has been accompa-
nied by a reorientation of the jet axis.
Several processes can be used to explain a jet reorientation
in AGNs. There are strong observational and theoretical grounds
for believing that accretion disks around black holes may be
twisted or warped, and this can be caused by a number of pos-
sible physical processes. In particular, if there is a misalignment
between the axis of rotating black hole and the axis of its rotating
accretion disk, then the Lense-Thirring precession produces a
warp in the disk. This process is called the Bardeen-Peterson ef-
fect (Bardeen & Petterson, 1975). According to Pringle (1997),
disk warping can also be induced by internal instabilities in
the accretion disk caused by radiation pressure from the central
source.
A reorientation of the jet axis may also result from a merger
with another black hole. Merritt & Ekers (2002) have shown that
a rapid change in jet orientation can be caused by even a mi-
10 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
1407+369 4994.500 MHz
peak flux density=109.87 mJy/beam, beam size=52 x 46 mas
first contour level=0.12 mJy/beam
RIGHT ASCENSION (J2000)
14 09 09.56 09.54 09.52 09.50 09.48 09.46
36 42 08.8
1407+369 1667.474 MHz
peak flux density=147.37 mJy/beam, beam size=10.2 x 5.1 mas
first contour level=0.18 mJy/beam
RIGHT ASCENSION (J2000)
14 09 09.516 09.512 09.508 09.504 09.500
36 42 08.30
08.25
08.20
08.15
08.10
08.05
1407+369 4987.474 MHz
peak flux density=60.90 mJy/beam, beam size=3.6 x 2.0 mas
first contour level=0.16 mJy/beam
RIGHT ASCENSION (J2000)
14 09 09.512 09.510 09.508 09.506 09.504
36 42 08.22
08.20
08.18
08.16
08.14
08.12
08.10
1407+369 8421.474 MHz
peak flux density=24.34 mJy/beam, beam size=2.1 x 1.0 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
14 09 09.511 09.510 09.509 09.508 09.507 09.506
36 42 08.18
08.17
08.16
08.15
08.14
08.13
08.12
Fig. 8. The MERLIN 5-GHz map (upper left) and VLBA 1.7, 5, and 8.4-GHz maps of 1407+369. Contours increase by a factor 2,
and the first contour level corresponds to ≈ 3σ. Crosses indicate the position of an optical object found using the SDSS/DR5.
nor merger because of a spin-flip of the central active black
hole arising from the coalescence of inclined binary black holes.
According to Liu (2004), the Bardeen-Peterson effect can also
cause a realignment of a rotating SMBH and a misaligned ac-
cretion disk, where the timescale of such a realignment t < 105
years. If it is assumed that the typical speed of advance of ra-
dio lobes of young AGNs is υ ∼0.3c (Owsianik et al., 1998;
Giroletti et al., 2003; Polatidis & Conway, 2003), then distorted
jets of length, tυ <10 kpc for some CSS and GPS sources should
be observed, although the character of these disturbances is not
known. Liu (2004) shows that the interaction/realignment of
a binary and its accretion disk leads to the development of X-
shaped sources. 1045+352 is not a typical X-shaped source like
3C 223.1 or 3C 403 (Dennett-Thorpe et al., 2002; Capetti et al.,
2002). However, according to Cohen et al. (2005) the realign-
ment of a rotating SMBH followed by a repositioning of the ac-
cretion disk and jets is a plausible interpretation for misaligned
radio structures, even if they are not conspicuously X-shaped.
It is likely that in young sources such as 1045+352, the
gas has not yet settled into a regular disk following a merger
event and that separate clouds of gas and dust reaching the very
central regions of the source at different times disturb the sta-
bility of the accretion disk and affect the jet formation. Later,
these clouds could cause a renewal of activity. Numerical simu-
lations of colliding galaxies show that these usually merge com-
pletely after a few encounters in timescales up to ∼ 108 years
(Barnes & Hernquist, 1996). According to Schoenmakers et al.
(2000), multiple encounters between interacting galaxies can
cause interruptions of activity and lead to the many types of
sources that are observed in a restarted phase, such as double-
double radio galaxies. Nevertheless, it is unclear whether such
encounters can cause jet reorientation. On the other hand, the
dense medium of a host galaxy can frustrate the jets, and their
collisions with the dense surrounding medium can cause rapid
bends through large angles. In the case of 1045+352, the VLBA
images at the higher frequencies seem to show a jet emerging in
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 11
1425+287 8439.900 MHz
peak flux density=64.89 mJy/beam, beam size=305 x 267 mas
first contour level=0.08 mJy/beam
RIGHT ASCENSION (J2000)
14 27 38.7 38.6 38.5 38.4 38.3 38.2
28 33 17
1425+287 4994.500 MHz
peak flux density=62.37 mJy/beam, beam size=74 x 40 mas
first contour level=0.80 mJy/beam
RIGHT ASCENSION (J2000)
14 27 40.40 40.35 40.30 40.25 40.20
28 33 27.5
Fig. 9. The VLA 8.4-GHz map and MERLIN 5-GHz map of 1425+287. Contours increase by a factor 2, and the first contour level
corresponds to ≈ 3σ.
Frequency (Hz)
1045+352
Fig. 11. Spectral Energy Distribution (SED) of 1045+352 from
radio to submillimetre wavelengths. The errors are smaller than
the size of the symbols; 1.25 mm point (Haas et al., 2006) is
shown as a triangle, 850µm and 450 µm points (Willott et al.,
2002) are shown as filled circles, radio observations are shown
as asterisks. The solid curve is the parabolic fit f (x) = ax2+bx+c
to all radio data (yi), with a = −0.14, b = 1.91, c = −5.68, and
reduced χ2 = 12. The dashed curve is the linear fit f (x) = ax+ b
to radio data with ν > 1GHz, with a = −0.86, b = 7.91, and
reduced χ2 = 0.5.
a S/SE direction, but being bent through ∼ 60◦ to a NE direction
in the lower resolution 1.7-GHz image. The MERLIN lower res-
olution 5-GHz image might indicate that the jet has been bent
again and now emerges from the core in a NW direction.
It is difficult to find a convincing argument in favour of one
of the above-mentioned alternatives or to rule any of them out
based upon the extensive multifrequency data on 1045+352 pre-
sented here. However, if it is assumed that a merger is the most
probable cause of the ignition and restart of activity in radio
galaxies, this could mean that 1045+352 has undergone two
merger events in a very short period of time (∼ 105), which is un-
Table 2. 1045+352 properties
Parameter Value
u′ 22.12
g′ 21.38
r′ 20.81
i′ 20.14
z′ 20.08
AB 2.0
MB -22.05 (-24.05)
AV 1.5
MV -22.83 (-24.33)
log(R∗)(total) 4.9 (4.1)
log(R∗)(core) 3.8 (3.0)
Notes: Optical photometry from SDSS, corrected for Galactic extinc-
tion. AV taken from Willott et al. (2002). Quantities in parentheses are
corrected for intrinsic extinction.
likely. More probable is that the ignition of activity in 1045+352
has occurred during a merger event that is, as yet, incomplete and
that disturbed, misaligned radio jets result from the realignment
of a rotating SMBH or intermittent gas injection that interrupts
jet formation.
4.2. Other nine sources
Three sources from our sample (1126+293, 1407+369,
1627+289) show one- or two-sided core-jet structures, indicat-
ing that they are in an active phase of their evolution, although
the core-jet structure of 1126+293 is controversial. Our images
indicate that the western components are parts of the jet, which is
possibly precessing or being bent by interactions with the inter-
stellar medium. They could, however, also be hotspots of a radio
lobe. Unfortunately, our high frequency VLBA observations are
not sensitive enough to settle this problem. Three other sources
(1056+316, 1132+374, 1425+287) have visible radio cores and
parts of lobes or hotspots, indicating activity. 1132+374 is a CSO
object. In the case of one source, 1059+351, the VLBA obser-
12 M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources
1627+289 8439.900 MHz
peak flux density=75.42 mJy/beam, beam size=271 x 263 mas
first contour level=0.07 mJy/beam
RIGHT ASCENSION (J2000)
16 29 12.50 12.45 12.40 12.35 12.30 12.25 12.20 12.15 12.10 12.05
28 51 37
1627+289 4994.500 MHz
peak flux density=77.77 mJy/beam, beam size=70 x 39 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
16 29 12.36 12.34 12.32 12.30 12.28 12.26 12.24 12.22 12.20
28 51 35.5
1627+289 1667.474 MHz
peak flux density=40.35 mJy/beam, beam size=10.6 x 5.1 mas
first contour level=0.30 mJy/beam
RIGHT ASCENSION (J2000)
16 29 12.270 12.268 12.266 12.264 12.262 12.260 12.258
28 51 34.16
34.14
34.12
34.10
34.08
34.06
34.04
34.02
34.00
33.98
33.96
1627+289 4987.474 MHz
peak flux density=5.57 mJy/beam, beam size=4.2 x 1.9 mas
first contour level=0.15 mJy/beam
RIGHT ASCENSION (J2000)
16 29 12.267 12.266 12.265 12.264 12.263 12.262 12.261 12.260
28 51 34.11
34.10
34.09
34.08
34.07
34.06
34.05
34.04
34.03
34.02
Fig. 10. The VLA 8.4-GHz map (upper left), MERLIN 5-GHz map (upper right), and VLBA 1.7 and 5-GHz maps of 1627+289.
Contours increase by a factor 2, and the first contour level corresponds to ≈ 3σ.
vations show only a radio core, although the 5-GHz MERLIN
image of 1059+351 also shows remnants of the two radio
lobes of its “S” shaped structure visible at the VLA resolutions
(Machalski & Condon, 1983; Machalski, 1998). According to
Taylor et al. (1996) and Readhead et al. (1996), “S” symmetry
is observed in many compact sources and can be explained
by precession of the central engine. 1059+351 is the largest
source in our sample with a linear size of 45 kpc based upon
its largest angular size measured from 1.46-GHz VLA image
(Machalski & Condon, 1983).
The compact 1049+384 and 1302+356 steep spectrum
sources appeared to be low-frequency variables (LFV) at
151 MHz with very high (≥0.99) probabilities that their variabil-
ity is real (Minns & Riley, 2000). According to them, LFV ob-
jects are generally more compact than other CSS sources and
tend to exhibit steeper spectra than typical CSS sources. This
may be because of rapid spectral ageing, which might be ex-
pected for frustrated sources, or it might simply be because the
sources are at very high redshifts.
5. Conclusions
VLBA, VLA, and MERLIN images of ten compact steep
spectrum sources have been presented. One of these sources,
1045+352, is a very radio-luminous BAL quasar, whose com-
plex structure suggests restarted activity. This may have resulted
either from a merger event or from the infall of a cloud of gas,
that had cooled in the halo of the galaxy into the core region of
the source. The asymmetric radio jets of 1045+352 and the es-
timated angle suggest that some of the emission can be boosted,
although the intrinsic asymmetries cannot be ruled out. It has
also been confirmed that the 850µm flux of 1045+352 can be
severely contaminated by synchrotron emission, which may sug-
gest less than previously estimated values of infrared emission
and dust mass. Most of the radio-loud BAL quasars detected to
M. Kunert-Bajraszewska and A. Marecki: FIRST-based survey of compact steep spectrum sources 13
Table 3. Flux densities of sources principal components from the VLBA observations
Source RA DEC S1.7 GHz S5 GHz S8.4GHz θ1 θ2 PA
Name h m s ◦ ′ ′′ mJy mJy mJy mas mas ◦
(1) (2) (3) (4) (5) (6) (7) (8) (9)
1045+352 10 48 34.248 34 57 25.044 303.2 − − 15.0 11.0 60
10 48 34.249 34 57 25.061 − 3.5 − 2.0 1.0 76
10 48 34.248 34 57 25.041 − 21.8 7.1 7.0 1.0 101
10 48 34.248 34 57 25.043 − 32.7 12.3 4.0 3.0 95
1049+384 10 52 11.803 38 11 44.018 13.6 − − 3.0 1.0 14
10 52 11.797 38 11 44.027 11.4 3.9 6.9 2.0 2.0 121
10 52 11.789 38 11 44.031 182.1 33.6 12.9 8.0 1.0 119
10 52 11.787 38 11 44.048 218.5 23.9 2.3 5.0 3.0 177
1056+316 10 59 43.254 31 24 20.106 8.8 − − 0.9 0.1 7
10 59 43.235 31 24 20.538 43.6 − − 33.0 8.0 6
1059+351 11 02 08.726 34 55 08.709 8.1 − − 0.7 0.3 124
1126+293 11 29 21.755 29 05 06.402 7.3 − − 3.0 1.0 84
11 29 21.753 29 05 06.401 10.4 − − 13.0 4.0 53
1132+374 11 35 05.934 37 08 40.810 124.1 6.6 1.6 18.0 2.0 57
11 35 05.932 37 08 40.775 36.3 13.8 9.4 2.0 0.4 8
11 35 05.931 37 08 40.715 14.5 − − 5.0 0.8 105
1302+356 13 04 34.495 35 23 33.534 46.8 5.9 − 11.0 6.0 97
13 04 34.494 35 23 33.538 60.5 − − 15.0 7.0 147
1407+369 14 09 09.504 36 42 08.195 81.0 1.9 − 17.0 3.0 138
14 09 09.508 36 42 08.164 192.8 76.7 42.0 8.0 1.5 141
14 09 09.508 36 42 08.152 − 9.5 4.8 0.7 0.2 140
1627+289 16 29 12.264 28 51 34.062 111.5 8.1 − 10.0 6.0 58
Description of the columns: (1) source name in the IAU format; (2) component right ascension (J2000) as measured at 1.7 GHz; (3) component
declination (J2000) as measured at 1.7 GHz; (4) VLBA flux density in mJy at 1.7 GHz from the present paper; (5) VLBA flux density in mJy at
5 GHz from the present paper; (6) VLBA flux density in mJy at 8.4 GHz from the present paper; (7) deconvolved component major axis angular
size at 1.7 GHz obtained using JMFIT; (8) deconvolved component minor axis angular size at 1.7 GHz obtained using JMFIT; (9) deconvolved
major axis position angle at 1.7 GHz obtained using JMFIT. In the case the component is not visible in 1.7 GHz map the values for the last three
columns are taken from the 5-GHz image.
date have very compact radio structures similar to GPS and CSS
sources which are thought to be young. Therefore, the compact
structure and young age of 1045+352 fit well to the evolutionary
interpretation of radio-loud BAL QSOs.
According to the evolutionary model recently proposed by
Lipari & Terlevich (2006), BAL quasars are young systems with
composite outflows, and they are accompanied by absorption
clouds. The radio-loud systems may be associated with the later
stages of evolution, when jets have removed the clouds respon-
sible for the generation of BALs. The effect of orientation could
play a secondary role here. The above could explain the rarity of
extended radio structures showing BAL features (Gregg et al.,
2006).
Acknowledgements.
The VLBA is operated by the National Radio Astronomy Observatory (NRAO),
a facility of the National Science Foundation (NSF) operated under cooperative
agreement by Associated Universities, Inc. (AUI).
This research has made use of the NASA/IPAC Extragalactic Database (NED),
which is operated by the Jet Propulsion Laboratory, California Institute
of Technology, under contract with the National Aeronautics and Space
Administration.
Use has been made of the Sloan Digital Sky Survey (SDSS) Archive. The
SDSS is managed by the Astrophysical Research Consortium (ARC) for the
Participating Institutions: The University of Chicago, Fermilab, the Institute for
Advanced Study, the Japan Participation Group, The Johns Hopkins University,
Los Alamos National Laboratory, the Max-Planck-Institute for Astronomy
(MPIA), the Max-Planck-Institute for Astrophysics (MPA), New Mexico State
University, University of Pittsburgh, Princeton University, the United States
Naval Observatory, and the University of Washington.
We thank M. Gawroński for his help with the OCRA-p observations. The OCRA
project was supported by the Polish Ministry of Science and Higher Education
under grant 5 P03D 024 21 and the Royal Society Paul Instrument Fund.
We thank P.J. Wiita for a discussion and P. Thomasson for reading of the paper
and a number of suggestions.
This work was supported by the Polish Ministry of Science and Higher
Education under grant 1 P03D 008 30.
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List of Objects
‘1045+352’ on page 3
‘1049+384’ on page 3
‘1056+316’ on page 3
‘1059+351’ on page 3
‘1126+293’ on page 4
‘1132+374’ on page 4
‘1302+356’ on page 4
‘1407+369’ on page 5
‘1425+287’ on page 5
‘1627+289’ on page 5
Introduction
The observations and data reduction
Comments on individual sources
Discussion
1045+352 — a BAL quasar
Other nine sources
Conclusions
|
0704.0352 | Investigation of relaxation phenomena in high-temperature
superconductors HoBa2Cu3O7-d at the action of pulsed magnetic fields | Microsoft Word - article.doc
Investigation of of of of relaxation phenomena in high-temperature superconductors HoBa2Cu3O7-δ at
the action of pulsed magnetic fields
J.G. Chigvinadze*, J.V. Acrivos**, S.M. Ashimov*, A.A. Iashvili*, T. V. Machaidze*, Th. Wolf***
* E. Andronikashvili Institute of Physics, 0177 Tbilisi, Georgia
** San Jose’ State University, San Jose’ CA 95192-0101, USA
*** Forschungszentrum Karlsruhe, Institut für Festkörperphysik, 76021 Karlsruhe, Germany
Summary
It is used the mechanical method of Abrikosov vortex
stimulated dynamics investigation in superconductors.
With its help it was studied relaxation phenomena in
vortex matter of high-temperature superconductors. It
established that pulsed magnetic fields change the course
of relaxation processes taking place in vortex matter. The
study of the influence of magnetic pulses differing by
their durations and amplitudes on vortex system of
isotropic high-temperature superconductors system
HoBa2Cu3O7-δ showed the presence of threshold
phenomena. The small duration pulses doesn’t change
the course of relaxation processes taking place in vortex
matter. When the duration of pulses exceeds some
critical value (threshold), then their influence change the
course of relaxation process which is revealed by
stepwise change of relaxing mechanical moment
relτ .
These investigations showed that the time for
formatting of Abrikosov vortex lattice in HoBa2Cu3O7-δ
is of the order of 20 �s which on the order of value
exceeds the time necessary for formation of a single
vortex observed in type II superconductors.
1. Introduction
The present communication is devoted to the
experimental investigation of relaxation phenomena in
high-temperature superconductors of HoBa2Cu3O7-δ
system.
High-temperature superconductors are characterized
by such high critical transition temperatures Тс in the
superconducting state, they remain superconductors at
temperatures when their thermal fluctuations energy
becomes compared with the elastic energy, and also with
the pinning energy [1]. It creates prerequisites for phase
transitions. Due to the layered crystal structure and
anisotropy, which is a characteristic high-temperature
superconductors, they reveal conditions for the
appearance of different phases on B-T diagram.( B is
magnetic induction, T-is temperature)[2-13]. As
example, Abrikosov vortex lattice begin melting near the
critical Тс temperature what is followed by the essential
change of vortex continuum flow dynamics along with
sharp change of character (dynamics) of relaxation
phenomena. In high-temperature superconductors it is
observed such relaxation processes as a slow logarithmic
decrease of captured flux with time at temperatures much
below their superconductive critical transition
temperature Тс [14-16]. The logarithmic character of
relaxation is explained by the Anderson [17]. Near Тс, in
the range of Abrikosov vortex lattice melting, the
logarithmic character of relaxation is changed by the
power one with 2/3 exponent [18].
Consequently, the study of relaxation processes in
high-temperature superconductors is an important
problem.
2. Experimental
For Investigation it was used currentless mechanical
method of Abrikosov vortex stimulated dynamics study
by magnetic pulses revealing relaxation phenomena in
vortex matter described in work [19]. This method is a
development of currentless mechanical method of
pinning investigations [20,21] and is based on pinning
forces countermoments measurements and viscous
friction, acting on a axially symmetrical superconducting
sample in an outer (transverse) magnetic field.
Countermoments of pinning forces and of viscous
friction, acting on a superconductive sample from
quantized vortex lines side (Abrikosov vortices) are
defined the way as it was described in work [22,23]. The
sensitivity of the method accordingly works [24], is
equivalent to 10-8 V×cm-1 in the method of V-A
characteristics.
The high-temperature superconducting samples of
HoBa2Cu3O7-δ system were prepared by the standard
solid state reaction method. Samples were made
cylindrical with height L=13mm and diameter d=6mm.
Their critical temperature was Tc=92 K. The investigated
samples were isotropic what was established by
mechanical moment τ measurements appearing
H > 1cH with the penetration of Abrikosov vortices
into a freely suspended on a thin elastic thread
superconducting sample. The appearance of such
moment ατ sinMH= , characteristic for anisotropic
superconductors, is related with penetrating Abrikosov
vortices and the mean magnetic moment M
of a sample
which could deviate on angle α from the direction of
outer magnetic field H
. In superconducting anisotropic
samples it is presented energetically favorable directions
for the arrangement of emerging (penetrating) vortex
lines which in their turn are fastened by pinning centers
creating aforementioned moment τ . The lack of τ
moment is characteristic for isotropic and investigated by
us samples, no matter magnetic field value and its
previous orientation in respect to H
in the axial
symmetry plane. Pulsed magnetic fields were created by
Helmholtz coils. The value of pulsed magnetic fields was
changed in Oeh 2002 ÷=∆ limits.
In experiments it was used both single and continuous
pulsed with repetition frequency ν from 2.5 s-1 to 500s-1
. The duration
x of pulses was changed from 0,5 до
500 �s. Magnetic pulse could be directed both parallel
h||H) and perpendicularly (
h⊥H) to applied steady
magnetic field H
, creating mixed state of
superconducting sample. The standard pulsed generator
and amplifier were used to feed Helmholtz coils. The
current strength in coils reached up to 40÷50 A.
Samples were high-temperature superconductors of
HoBa2Cu3O7-δ system placed in the center between
Helmholtz coils.
The principal set-up of experiment is shown in fig.1
[19,20]. In experiments it is measured the rotation angle
2ϕ of sample depending on the angle of rotation of a
torsion head 1ϕ , transmitting the rotation to a sample by
means of suspension having the torsion stiffness
K ≈4·10-1 [dyn•cm], which can be replaced when
necessary by a less stiff or stiffer one.
The measurements were carried out at a constant speed
of rotation of the torsion head, making ω1=1,8·10
-2 rad/s .
Angles of rotation поворота 2ϕ and 1ϕ were
determined with an accuracy of ±4,6·10-3 and ±2,3·10-3
rad, respectively. The uniformity of the magnetic field’s
strength along a sample was below H
∆ = 10-3.
Fig. 1. The schematic diagram and the geometry of the experiment. 1-sample, 2-upper elastic filament, 3-lower filament, 4 - leading head,
5 - glass road. φ is angle between Mr and Hr
To avoid effects, connected with the frozen magnetic
fluxes, the lower part of the cryostat with the sample was
put into a special cylindrical Permalloy screen, reducing
the Earth magnetic field by the factor of 1200. After a
sample was cooled by liquid nitrogen to the
superconducting state, the screen was removed, a
magnetic field of necessary intensity H was applied and
the 2 1( )ϕ ϕ dependences were measured. To carry out
measurements at different values of H , the sample was
brought to the normal state by heating it to до T > cT at
H =0, and only after returning sample and torsion head
to the initial state 1 2 0ϕ ϕ= = , the experiment was
repeated.
3. Results and discussions
During rotation of the sample both of normal and
superconducting states in the
absence external magnetic field ( H =0) the 2ϕ
dependence versus 1ϕ is linear and the condition is
satisfied.
tωϕϕ == 21
The character of the 2 1( )ϕ ϕ dependence is changed
significantly, when the sample is in magnetic fields
H > 1cH at T < cT . Typical 2 1( )ϕ ϕ dependences at
T=77K and various magnetic fields for HoBa2Cu3O7-δ
sample ( length of a cylindrical sample L=13mm and
diameter d=6mm ) is shown in Fig.2.
Fig.2. Dependence of the rotation angle of the sample HoBa2Cu3O7-δ
2ϕ on the rotation angle of the leading head 1ϕ in magnetic field
H=1000 Oe at T=77K.
Three distinct regions are observed in Fig.2. In the
first (initial) region, the sample does not respond to the
increase in 1ϕ , i.e. to the applied and increase with
time torsion torque as 1ϕ ~ )( 21 ϕϕτ −= K or
responds weakly. Such behavior of the sample can be
explained by fact that Abrikosov vortices are not
detached from pinning centers at small values of 1ϕ ~τ ,
but if the sample is still turned slightly, this can be
caused by elastic deformation of magnetic force lines
beyond it or, possibly, by separation of the most weakly
fixed vortices. As it is seen from fig.2 , as soon as a
certain critical value φсmin depending on H is reached ,
the first region under goes a transition to the second
region in which the velocity of the sample increases
gradually with 1ϕ increasing resulting from the
progressive process of detachment of vortices from their
corresponding pinning centers. One should expect that
just in this region, in the rotating sample “the vortices
fan” begins to unfold, in with the vortices are distributed
according to the instantaneous angles of orientations with
respect to the fixed external magnetic field. In this case
the of orientation angles of separate vortex filaments are
limited from frϕ to pinfr ϕϕ + , where frϕ is the angle
on which the vortex filament can be turned with respect
to H
by forces of viscous friction with the matrix of
superconductor, and pinϕ is the angle on with the
vortex filament can be turned by the most strong pinning
center, studied for the first time in [25].
The gradual transition (at high 1ϕ values) to the third
region where the linear 2 1( )ϕ ϕ dependence was
observed, allows one to define the countermoments of
pinning forces pτ and frτ , independently. Just in this
region, when 21 ωω = the torque τ , appeared to the
uniformly rotating sample, is balanced by the
countermoment pτ and frτ . In particular, in the case
of continuously rotating sample with frequency
21 ωω = one could find similarly to [26,27] the
expression for the total braking torque τ [19] .
Indeed, if we consider in this case a vortex element
moving with velocity ⊥υ
perpendicular to sd
then the average force acting on this elements is
dsFdsfd l
⊥ += υ
and the associated braking torque, exerted on the rotating
specimen becomes:
υτ fdrd
where r
is the vector pointing from the rotational axis
to the vortex elements, lF is the pinning force per flux
thread per unit length, and η is the viscosity coefficient.
For a cylindrical specimen of radius R and height L
integrating over the individual contribution of all vortex
gives a total braking torque τ
ωτττ 0+= p (1)
with
=τ ,
and
B ηπτ
= ,
Where B is the inductivity averaged over the sample,
0Φ is the flux quantum , L is the height and R is the
radius of the sample.
As it is shown in Fig.2, starting with the point
(a), where 21 ωω = , to the superconducting sample
uniformly rotating in the homogeneous stationary
magnetic field H=1000 Oe, is applied stationary
dynamic torsion moment fr
p τττ += .
If in this region the torsion head is stopped, then at the
expense of relaxation processes connected with the
presence of viscous forces acting on vortex filaments, the
sample will continue the rotation in the same direction
(with decreasing velocity) until it reaches a certain
equilibrium position, depending on the H value. The
Fig.3 shows curves of
2ϕ∆ time dependences at the
stopped leading head for HoBa2Cu3O7-δ sample at
T=77K and H=1000 Oe.
Fig.3. Dependence of momentum
relτ on time t after the stopping of
rotating head for HoBa2Cu3O7-δ sample at T=77K and H=1000 Oe.
If during the relaxation after rotation of sample one
applies the pulsed magnetic field in parallel to the outer
magnetic field H
, then additional vortices, created as
result of magnetic pulse, influence the structure already
existing in the sample as “the vortex fan” what could
result in the decrease of the angle of its unfolding or to
its folding. The letter in its turn, would cause the
additional change in the relaxation process taking place
in the sample, and, correspondingly, results in the
stepwise decrease of moment related with viscous
forces frτ .
But the change of relaxation process character and,
correspondently, the stepwise decrease of moment could
happen if the duration of magnetic pulse is larger as
compared with the time necessary for creation of a new
vortex structure, which will influence the
superconducting sample relaxing in magnetic field. If it
is the case, then at the small durations of magnetic pulses
the relaxation curve, presented in Fig.3, doesn’t change,
but when this duration becomes the order of a time for
penetration of vortices into the sample and the creation
of vortex structure, then the aforementioned change of
relaxation processes could principally appear. Namely,
this situation when the duration �x of magnetic pulses
was larger then the time for Abrikosov vortex lattice
creation �xс, have been described by us our previous
work [19], when it was shown that the influence of one
magnetic pulse �h≈400 Oe (�h||H) with duration
30�сек>�xс was stepwisely decreased the relτ moment
and the relaxation process continued with the reduction
relτ on a level as far as a new magnetic pulse similar
the first one is not applied.
In the presented work it was studied the influence of
different duration and amplitude pulses on relaxation
processes in vortex matter. The results shown in Fig.4
on action of single pulses of different durations on
relaxation processes in vortex matter and, consequently,
on mechanical moment
relτ revealed that at small pulses
durations up to 15 �s the relτ doesn’t change, but at
duration of applied pulse >15 �s it is observed the
stepwise change
relτ , what speaks on the existence of
the �xс threshold.
Fig.4. Dependence of momentum � relτ on the duration �x of
magnetic field single pulse � h=172 Oe applied in parallel to the
main magnetic field H=1000 Oe at T=77K for HoBa2Cu3O7-δ
sample.
This way one could say that the Abrikosov vortex
lattice creation time in high-temperature isotropic
superconductor of HoBa2Cu3O7-δ makes value on the
order of 20�s. This value approximately on the order of
value higher then time for the single-vortex creation for
the first time measured by G. Boato, G.Gallinaro and C.
Rizzuto [28], who showed that this time is less than 10-5
sec.
In work [19] it was also shown that continuous action
of aforementioned pulses with the train frequency equal
to 2,5 s-1 more sharply reveals their influence on
relaxation processes in vortex matter and in these
conditions the processes of penetration of vortices into
superconductors bulk are made more sharply expressed.
In fig.5 it is presented the clear picture of magnetic
pulses continuous action with �h=172 Oe (�h||H) , and
the duration 20�sec, what is larger than the �xс with the
train frequency ν = 2,5 s-1. As it is seen from picture the
pulses of 5, 10 and 15 �sec durations doesn’t change
)(tfrel =τ which is observed at absence of magnetic
pulses.
The results presented in Fig.5 show that at durations of
pulses in 20�sec, 30�sec and 40�sec the Abrikosov
vortices penetrate into the superconductor. This way the
threshold value on the magnetic pulses duration observed
at the action of single pulses (Fig.4) coinside with the
threshold observed when their repetition frequency is ν
= 2.5 s-1.
Fig.5. Dependence of momentum
relτ on time t after the stopping of
rotating head with the influence since t=5 min on the relaxation
process of HoBa2Cu3O7-δ sample of the continuous magnetic field
pulses h=172 Oe with ν=2,5 s-1 frequency and different
durations x=5; 10; 15; 20; 30; and 40 �s. Pulsed magnetic field
was parallel to the main magnetic field H=1000 Oe at T=77K.
In Fig.6 it is presented the curve of relτ =f(t)
dependence on time at the influence of magnetic pulses �h=172 (�h||H) the duration of which is below the time
of Abrikosov vortex system creation �x=5�s<�xс
(�xс≥15�s for the investigated HoBa2Cu3O7-δ). As it is
seen from the picture when �x<�xс, the relaxation curve
doesn’t change in spite the increase of the repetition
frequency of magnetic pulses ν from 2.5 up to 500 s-1. As
soon as the duration of pulses exceeds the critical value
and becomes �x=30�s, the relaxation curve undergoes
the essential (stepwise) change. For example in Fig.6 it
is presented measurement for ν=5s-1 и ν= 500s-1.
Fig.6. Dependence of momentum
relτ on the time t after the
stopping of rotating head with the influence since t=10 min on the
HoBa2Cu3O7-δ sample relaxation process of the continuous
pulses magnetic field with frequency ν=2,5 ÷500s-1 at �x=5�s< �xс ,
and also at �x= 30�s > �xс. The pulsed magnetic field �h=172 Oe was
parallel to the main magnetic field H=1000 Oe at T=77K.
And finally, we have observed the threshold on the
value of applied pulses. In Fig.7 it is shown that in spite
the fact that we applied magnetic pulses of the large
duration 300�s>>�xс, much longer as compared with the
time of Abrikosov vortex creation at small amplitudes of
pulsed field �h ~7, 11, 14 Oe relτ =f(t) doesn’t change.
The stepwise change of the relaxing moment
relτ is
revealed only at �h ~18 Oe and higher.
Fig.6. Dependence of momentum
relτ on the time t after the
stopping of rotating head with the application after 5 minutes on the
HoBa2Cu3O7-δ sample relaxation process of the single magnetic
field pulses �h=(7÷36) Oe with duration �x= 300�s >>�xс.
The pulsed magnetic field was parallel to the main magnetic field
H=400 Oe at T=77K.
The further investigations of relaxation phenomena are
anticipated for anisotropic high-temperature
superconductors among them in strongly anisotropic
high-temperature superconductors of Bi-Pb-Sr-Ca-Cu-O
system.
4. Conclusion
The simple mechanical method of Abrikosov vortex
stimulated dynamics investigations it was applied for the
study of pulsed magnetic fields influence on relaxation
phenomena in vortex matter of high- temperature
superconductors. It was observed the change of
relaxation processes in vortex matter as a result of pulsed
magnetic field influence on it.
The study of influence of different duration and
amplitude pulsed magnetic fields influence was revealed
the existence of threshold phenomena. A small duration
pulse doesn’t change the course of relaxation processes
in vortex matter of isotropic high- temperature
superconductor HoBa2Cu3O7-δ. When the duration of
pulses exceeds some critical value (threshold), then their
influence change the course of relaxation processes. The
latter is revealed in a stepwise decrease of relaxing
mechanical momentum
relτ , apparently, related with a
sharp change of pinning and the rearrange of vortex
system of superconducting sample as a result of
penetration into its bulk of a new portion of vortices at
application of pulsed field on the outer magnetic field
creating the main vortex structure in the investigated
HoBa2Cu3O7-δ sample. A new portion of vortices
“shakes” the vortex lattice existing in a sample causing
the detachment of vortices from a weak pinning centers
what, apparently, is the reason for the stepwise decrease
of mechanical momentum
relτ .
All these made it possible to define the Abrikosov
vortex lattice creation time in HoBa2Cu3O7-δ which
turned out to be on the order of value higher as compared
with the time of single- vortex creation observed in type
II superconductors.
Acknowledgements
The work was supported by the grants of International
Science and Technology Center (ISTC) G-389 and G-
593.
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|
0704.0353 | Spin and pseudospin symmetries and the equivalent spectra of
relativistic spin-1/2 and spin-0 particles | Spin and pseudospin symmetries and the equivalent spectra of relativistic spin-1/2
and spin-0 particles
P. Alberto
Physics Department and Center for Computational Physics,
University of Coimbra, P-3004-516 Coimbra, Portugal
A. S. de Castro
Departamento de F́ısica e Qúımica, Universidade Estadual Paulista, 12516-410 Guaratinguetá, SP, Brazil
M. Malheiro
Departamento de F́ısica, Instituto Tecnológico de Aeronáutica,
CTA, 12228-900, São José dos Campos, SP, Brazil
and Instituto de F́ısica, Universidade Federal Fluminense, 24210-340 Niterói, Brazil
(Dated: November 4, 2018)
We show that the conditions which originate the spin and pseudospin symmetries in the Dirac
equation are the same that produce equivalent energy spectra of relativistic spin-1/2 and spin-0
particles in the presence of vector and scalar potentials. The conclusions do not depend on the
particular shapes of the potentials and can be important in different fields of physics. When both
scalar and vector potentials are spherical, these conditions for isospectrality imply that the spin-
orbit and Darwin terms of either the upper component or the lower component of the Dirac spinor
vanish, making it equivalent, as far as energy is concerned, to a spin-0 state. In this case, besides
energy, a scalar particle will also have the same orbital angular momentum as the (conserved) orbital
angular momentum of either the upper or lower component of the corresponding spin-1/2 particle.
We point out a few possible applications of this result.
PACS numbers: 11.30.-j,03.65.Pm
When describing some strong interacting systems it is often useful, because of simplicity, to approximate the
behavior of relativistic spin-1/2 particles by scalar spin-0 particles obeying the Klein-Gordon equation. An example
is the case of relativistic quark models used for studying quark-hadron duality because of the added complexity of
structure functions of Dirac particles as compared to scalar ones. It turns out that some results (e.g., the onset of
scaling in some structure functions) almost do not depend on the spin structure of the particle [1]. In this work we
will give another example of an observable, the energy, whose value may not depend on the spinor structure of the
particle, i.e., whether one has a spin-1/2 or a spin-0 particle. We will show that when a Dirac particle is subjected
to scalar and vector potentials of equal magnitude, it will have exactly the same energy spectrum as a scalar particle
of the same mass under the same potentials. As we will see, this happens because the spin-orbit and Darwin terms
in the second-order equation for either the upper or lower spinor component vanish when the scalar and vector
potentials have equal magnitude. It is not uncommon to find physical systems in which strong interacting relativistic
particles are subject to Lorentz scalar potentials (or position-dependent effective masses) that are of the same order
of magnitude of potentials which couple to the energy (time components of Lorentz four-vectors). For instance, the
scalar and vector (hereafter meaning time-component of a four-vector potential) nuclear mean-field potentials have
opposite signs but similar magnitudes, whereas relativistic models of mesons with a heavy and a light quark, like D-
or B-mesons, explain the observed small spin-orbit splitting by having vector and scalar potentials with the same sign
and similar strengths [2].
It is well-known that all the components of the free Dirac spinor, i.e., the solution of the free Dirac equation, satisfy
the free Klein-Gordon equation. Indeed, from the free Dirac equation
(i~γµ∂µ −mc)Ψ = 0 (1)
one gets
(−i~γν∂ν −mc)(i~γ
µ∂µ −mc)Ψ = (~
2∂µ∂µ +m
2c2)Ψ = 0 , (2)
where use has been made of the relation γµγν∂µ∂ν = ∂µ∂
µ. In a similar way, for the time-independent free Dirac
equation we would have
(cα · p+ βmc2)ψ = (−i~cα · ∇+ βmc2)ψ = Eψ , (3)
http://arxiv.org/abs/0704.0353v1
where, as usual, ψ(r) = Ψ(r, t) exp (i E t/~), α = γ0γ and β = γ0. Then, by left multiplying Eq. (3) by cα ·p+βmc2,
one gets the time-independent free Klein-Gordon equation
(c2p2 +m2c4)ψ = (−~2c2∇2 +m2c2)ψ = E2ψ , (4)
where the relation {β,α} = 0 was used. This all means that the free four-component Dirac spinor, and of course all
of its components, satisfy the Klein-Gordon equation. This is not surprising, because, after all, both free spin-1/2 and
spin-0 particles obey the same relativistic dispersion relation, E2 = p2c2 +m2c4, in spite of having different spinor
structures and thus different wave functions. Since there is no spin-dependent interaction, one expects both to have
the same energy spectrum.
We consider now the case of a spin-1/2 particle subject to a Lorentz scalar potential Vs plus a vector potential Vv.
The time-independent Dirac equation is given by
[cα · p+ β(mc2 + Vs)]ψ = (E − Vv)ψ (5)
It is convenient to define the four-spinors ψ± = P±ψ = [(I ± β)/2]ψ such that
, (6)
where φ and χ are respectively the upper and lower two-component spinors. Using the properties and anti-
commutation relations of the matrices β and α we can apply the projectors P± to the Dirac equation (5) and
decompose it into two coupled equations for ψ+ and ψ−:
cα · pψ− + (mc
2 + Vs)ψ+ = (E − Vv)ψ+ (7)
cα · pψ+ − (mc
2 + Vs)ψ− = (E − Vv)ψ− . (8)
Applying the operator cα · p on the left of these equations and using them to write ψ+ and ψ− in terms of α · pψ−
and α · pψ+ respectively, we finally get second-order equations for ψ+ and ψ−:
c2p2 ψ+ + c
[α · p∆]α · pψ+
E −∆+mc2
= (E −∆+mc2)(E − Σ−mc2)ψ+ (9)
c2p2 ψ− + c
[α · pΣ]α · pψ−
E − Σ−mc2
= (E −∆+mc2)(E − Σ−mc2)ψ− (10)
where the square brackets [ ] mean that the operator α · p only acts on the potential in front of it and we defined
Σ = Vv + Vs and ∆ = Vv − Vs. The second term in these equations can be further elaborated noting that the Dirac
αi matrices satisfy the relation αiαj = δij +
iǫijkSk where Sk, k = 1, 2, 3, are the spin operator components. The
second-order equations read now
c2 p2 ψ+ + c
[p∆] · pψ+ +
[p∆]× p · S ψ+
E −∆+mc2
= (E −∆+mc2)(E − Σ−mc2)ψ+ (11)
c2 p2 ψ− + c
[pΣ] · pψ− +
[pΣ]× p · S ψ−
E − Σ−mc2
= (E −∆+mc2)(E − Σ−mc2)ψ−. (12)
Now, if p∆ = 0, meaning that ∆ is constant or zero (if ∆ goes to zero at infinity, the two conditions are equivalent),
then the second term in eq. (11) disappears and we have
c2 p2ψ+ = (E −∆+mc
2)(E − Σ−mc2)ψ+ = [(E − Vv)
2 − (mc2 + Vs)
2]ψ+ , (13)
which is precisely the time-independent Klein-Gordon equation for a scalar potential Vs plus a vector potential Vv[14].
Since the second-order equation determines the eigenvalues for the spin-1/2 particle, this means that when p∆ = 0,
a spin-1/2 and a spin-0 particle with the same mass and subject to the same potentials Vs and Vv will have the same
energy spectrum, including both bound and scattering states. This last sufficient condition for isospectrality can be
relaxed to demand that just the combination mc2+Vs be the same for both particles, allowing them to have different
masses. This is so because this weaker condition does not change the gradient of ∆ and Σ and therefore the condition
p∆ = 0 will still hold. On the other hand, if the scalar and vector potentials are such that pΣ = 0, we would obtain a
Klein-Gordon equation for ψ−, and again the spectrum for spin-0 and spin-1/2 particles would be the same, provided
they are subjected to the same vector potential and mc2 + Vs is the same for both particles. If both Vs and Vv are
central potentials, i.e., only depend on the radial coordinate, then the numerators of the second terms in equations
(11) and (12) read
[p∆] · pψ+ +
[p∆]× p · S ψ+ = −~
∆′L · S ψ+ (14)
[pΣ] · pψ− +
[pΣ]× p · S ψ− = −~
Σ′L · S ψ− , (15)
where ∆′ and Σ′ are the derivatives with respect to r of the radial potentials ∆(r) and Σ(r), and L = r × p is the
orbital angular momentum operator. From these equations ones sees that these terms, which set apart the Dirac
second-order equations for the upper and lower components of the Dirac spinor from the Klein-Gordon equation and
thus are the origin of the different spectra for spin-1/2 and spin-0 particles, are composed of a derivative term, related
to the Darwin term which appears in the Foldy-Wouthuysen expansion, and a L · S spin-orbit term. If ∆′ = 0
(Σ′ = 0), then there is no spin-orbit term for the upper (lower) component of the Dirac spinor. In turn, since
the second-order equation determines the energy eigenvalues, this means that the orbital angular momentum of the
respective component is a good quantum number of the Dirac spinor. This can be a bit surprising, since one knows
that in general the orbital quantum number is not a good quantum number for a Dirac particle, since L2 does not
commute with a Dirac Hamiltonian with radial potentials. The reason why this does not happen in these cases was
reported in Refs. [3, 4], and we now review it in a slight different fashion. Let us consider in more detail the case of
spherical potentials such that ∆′ = 0. One knows that a spinor that is a solution of a Dirac equation with spherically
symmetric potentials can be generally written as
ψjm(r) =
gj l(r)
Yj lm(r̂)
j l̃ m
. (16)
where Yj lm are the spinor spherical harmonics. These result from the coupling of spherical harmonics and two-
dimensional Pauli spinors χms , Yj lm =
〈 l ml ; 1/2ms | j m 〉Ylmlχms , where 〈 l ml ; 1/2ms | j m 〉 is a
Clebsch-Gordan coefficient and l̃ = l ± 1, the plus and minus signs being related to whether one has aligned or
anti-aligned spin, i.e., j = l ± 1/2. The spinor spherical harmonics for the lower component satisfy the relation
j l̃m
= −σ · r̂Yj lm. The fact that the upper and lower components have different orbital angular momenta is related
to the fact, mentioned before, that L2 does not commute with the Dirac Hamiltonian
H = cα · p+ β(Vs +mc
2) + Vv = cα · p+ βmc
2 +ΣP+ +∆P− , (17)
where P± are the projectors defined above. However, when ∆
′ = 0, there is an extra SU(2) symmetry of H (so-called
“spin symmetry”) as first shown by Bell and Ruegg [5]. When we have spherical potentials, Ginocchio showed that
there is an additional SU(2) symmetry (for a recent review see [4]). The generators of this last symmetry are
L = LP+ +
α · pLα · pP− =
0 Up LUp
, (18)
where Up = σ · p/(
p2) is the helicity operator. One can check that L commutes with the Dirac Hamiltonian,
[H,L] = [cα · p,LP+ +
α · pLα · pP−] + [∆,
α · pLα · p] + [Σ,L]
= [∆,
α · pLα · p ] = 0 , (19)
where the last equality comes from the fact that ∆′ = 0. The Casimir L2 operator is given by L2 = L2P+ +
α · pP−. Applying this operator to the spinor ψjm (16), we get
2ψjm = L
α · pL2 α · pψ−
= ~2l(l+ 1)ψ+
α · p cL2 ψ+jm
E −∆+mc2
= ~2l(l + 1)ψ+
+ ~2l(l + 1)ψ−
= ~2l(l + 1)ψjm , (20)
where ψ±jm = P±ψjm and we used the relation, valid when ∆
′ = 0, ψ+jm = (E −∆ +mc
α · p
ψ−jm. From (20) we
see that ψjm is indeed an eigenstate of L
2. Thus the orbital quantum number of the upper component l is a good
quantum number of the system when the spherical potentials Vs(r) and Vv(r) are such that Vv(r) = Vs(r)+C∆, where
C∆ is an arbitrary constant. Also, according to we have said before, there is a state of a spin-0 particle subjected to
these same spherical potentials (or, at least, with a scalar potential such that the sum Vs +mc
2 is the same) that
has the same energy and the same orbital angular momentum as ψjm. In addition, the wave function of this scalar
particle would be proportional to the spatial part of the wave function of the upper component.
Note that the generator of the “spin symmetry” S is given by a similar expression as (18) just replacing L by ~/2σ
[4, 5], meaning that S2 ≡ S2 = 3/4 ~2I so that spin is also a good quantum number, as would be expected. Actually,
one can show that the total angular momentum operator J can be written as L + S, so that l, ml (eigenvalue of
Lz), s = 1/2, ms (eigenvalue of Sz) are good quantum numbers. Then, of course, j and m = ml +ms are also good
quantum numbers, but only in a trivial way, because there is no longer spin-orbit coupling. Therefore, in the spinor
(16) one could just replace the spinor spherical harmonic Yj lm by Yl mlχms and Yj l̃m by −σ · r̂ Yl mlχms . Note that if
∆ is a nonrelativistic potential, ∆ ≪ mc2 and ∆′ ≪ m2c4/(~c), i.e., it is slowly varying over a Compton wavelength.
In this case, the spin-orbit term will also get suppressed. In fact, the derivative of the ∆ potential is the origin of
the well-known relativistic spin-orbit effect which appears as a relativistic correction term in atomic physics or in the
v/c Foldy-Wouthuysen expansion (only the derivative of Vv appears because usually no Lorentz scalar potential Vs is
considered, and therefore ∆ = Vv).
When Σ′ = 0, or Vv(r) = −Vs(r) + CΣ, with CΣ an arbitrary constant, there is again a SU(2) symmetry, usually
called pseudospin symmetry ([5, 6]) which is relevant for describing the single-particle level structure of several nuclei.
This symmetry has a dynamical character and cannot be fully realized in nuclei because in Relativistic Mean-field
Theories the Σ potential is the only binding potential for nucleons [7, 8]. For harmonic oscillator potentials this is
no longer the case, since ∆, acting as an effective mass going to infinity, can bind Dirac particles [9, 10], even when
Σ = 0. As before, in the special case of spherical potentials, there is another SU(2) symmetry whose generators are
α · pLα · pP+ +LP− =
Up LUp 0
. (21)
In the same way as before, applying L̃
to ψjm, we would find that L̃
ψjm = ~
2 l̃(l̃ + 1)ψjm, that is, this time it
is the orbital quantum number of the lower component l̃ which is a good quantum number of the system and can
be used to classify energy levels. Again, provided the vector and scalar potentials are adequately related, there
would be a corresponding state of a spin-0 particle with the same energy and same orbital angular momentum l̃,
and, furthermore, its wave function would be proportional to the spatial part of the wave function of the lower
component. As before, the pseudospin symmetry generator S̃ can be obtained from L̃ by replacing L by ~/2σ. The
good quantum numbers of the system would be, besides l̃, m
, s̃ ≡ s = 1/2 and ms̃. Again, J = L̃ + S̃. It is
interesting that, as has been noted by Ginocchio [9], the generators of spin and pseudospin symmetries are related
through a γ5 transformation since S̃ = γ5Sγ5 and L̃ = γ5Lγ5. This property was used in a recent work to relate
spin symmetric and pseudospin symmetric spectra of harmonic oscillator potentials [11]. There it was shown that
for massless particles (or ultrarelativistic particles) the spin- and pseudo-spin spectra of Dirac particles are the same.
In addition, this means that spin-symmetric massless eigenstates of γ5 would be also pseudo-spin symmetric and
vice-versa. Since in this case ∆ = Σ = 0, or Vv = Vs = 0, this is, of course, just another way of stating the well-known
fact that free massless Dirac particles have good chirality.
Naturally, for free spin-1/2 particles described by spherical waves, both l and l̃ are good quantum numbers, which
just reflects the fact that one can have free spherical waves with any orbital angular momentum for the upper or
lower component and still have the same energy, as long as their linear momentum magnitude is the same, or, put in
another way, the energy of a free spin-1/2 particle cannot depend on its direction of motion.
In summary, we showed that when a relativistic spin-1/2 particle is subject to vector and scalar potentials such
that Vv = ±Vs + C±, where C± are constants, its energy spectrum does not depend on their spinorial structure,
being identical to the spectrum of a spin-0 particle which has no spinorial structure. This amounts to say that if
the potentials have these configurations there is no spin-orbit coupling and Darwin term. If the scalar and vector
potentials are spherical, one can classify the energy levels according to the orbital angular momentum quantum
number of either the upper or the lower component of the Dirac spinor. This would then correspond to having a
spin-0 particle with orbital angular momentum l or l̃, respectively. This spectral identity can of course happen only
with potentials which do not involve the spinorial structure of the Dirac equation in an intrinsic way. For instance, a
tensor potential of the form iβσµν (∂µAν − ∂νAµ) does not have an analog in the Klein-Gordon equation, so that one
could not have a spin-0 particle with the same spectrum as a spin-1/2 particle with such a potential. This is the case
of the so-called Dirac oscillator [12] (see [10] for a complete reference list), in which the Dirac equation contains a
potential of the form iβσ0imωri = imωβα · r. Another important potential, the electromagnetic vector potential A,
which is the spatial part of the electromagnetic four-vector potential, can be added via the minimal coupling scheme
to both the Dirac and the Klein-Gordon equations. Since α · (p− eA)α · (p− eA) = (p− eA)2 + 2e~∇×A · S, the
spectra of spin-0 and spin-1/2 particles cannot be identical as long as there is a magnetic field present, even though
the condition Vv = ±Vs +C± is fulfilled. It is important also to remark that, since for an electromagnetic interaction
Vv is the time-component of the electromagnetic four-vector potential, this last condition is gauge invariant in the
present case, in which we are dealing with stationary states, i.e, time-independent potentials. So, in the absence of a
external magnetic field (allowing, for instance, an electromagnetic vector potential A which is constant or a gradient
of a scalar function), a spin-0 and spin-1/2 particle subject to the same electromagnetic potential Vv and a Lorentz
scalar potential fulfilling the above relation would have the same spectrum.
The remark made above about the similarity of spin-0 and spin-1/2 wave functions can be relevant for calculations
in which the observables do not depend on the spin structure of the particle, like some structure functions. One such
calculation was made by Paris [13] in a massless confined Dirac particle, in which Vv = Vs. It would be interesting
to see how a Klein-Gordon particle would behave under the same potentials. More generally, this spectral identity
can also have experimental implications in different fields of physics, since, should such an identity be found, it would
signal the presence of a Lorentz scalar field having a similar magnitude as that of a time-component of a Lorentz
vector field, or at least differing just by a constant.
Acknowledgments
We acknowledge financial support from CNPQ, FAPESP and FCT (POCTI) scientific program.
[1] S. Jeschonnek and J. W. Van Orden, Phys. Rev. D 69, 054006 (2004).
[2] P. R. Page, T. Goldman, and J. N. Ginocchio, Phys. Rev. Lett. 86, 204.
[3] J. N. Ginocchio and A. Leviatan, Phys. Lett. B425, 1 (1998).
[4] J. N. Ginocchio, Phys. Rep. 414 165 (2005).
[5] J. S. Bell and H. Ruegg, Nucl. Phys. B98, 151 (1975).
[6] J. N. Ginocchio, Phys. Rev. Lett. 78, 436 (1997).
[7] P. Alberto, M. Fiolhais, M. Malheiro, A. Delfino, and M. Chiapparini, Phys. Rev. Lett. 86, 5015 (2001).
[8] P. Alberto, M. Fiolhais, M. Malheiro, A. Delfino, and M. Chiapparini, Phys. Rev. C 65, 034307 (2002).
[9] J. N. Ginocchio, Phys. Rev. Lett. 95, 252501 (2005).
[10] R. Lisboa, M. Malheiro, A. S. de Castro, P. Alberto, and M. Fiolhais, Phys. Rev. C 69, 024319 (2004).
[11] A. S. de Castro, P. Alberto, R. Lisboa, and M. Malheiro, Phys. Rev. C 73, 054309 (2006).
[12] D. Itô, K. Mori, and E. Carriere, Nuovo Cimento A 51, 1119 (1967); M. Moshinsky and A. Szczepaniak, J. Phys. A 22,
L817 (1989).
[13] M. W. Paris, Phys. Rev. C 68, 025201 (2003).
[14] There are some authors who introduce a scalar potential Vs in the Klein-Gordon equation by making the replacement
m2c4 → m2c4 +V2
. Here we introduce it, as most authors do, as an effective mass m∗ 2 = (m+Vs/c
2)2, since it is the way
that it is introduced in the Dirac equation. The two potentials are related by V2
= (mc2 + Vs)
−m2c4.
Acknowledgments
References
|
0704.0354 | General asymptotic solutions of the Einstein equations and phase
transitions in quantum gravity | General asymptoti
solutions
of the Einstein equations and
phase transitions in quantum gravity
Dmitry Podolsky
Helsinki Institute of Physi
s, University of Helsinki,
Gustaf Hällströmin katu 2, FIN00014, Helsinki, Finland
Email: dmitry.podolsky�helsinki.�
Abstra
t
We dis
uss generi
properties of
lassi
al and quantum theories of grav-
ity with a s
alar �eld whi
h are revealed at the vi
inity of the
osmolog-
i
al singularity. When the potential of the s
alar �eld is exponential and
unbounded from below, the general solution of the Einstein equations
has quasi-isotropi
asymptoti
s near the singularity instead of the usual
anisotropi
Belinskii - Khalatnikov - Lifshitz (BKL) asymptoti
s. De-
pending on the strength of s
alar �eld potential, there exist two phases
of quantum gravity with s
alar �eld: one with essentially anisotropi
be-
havior of �eld
orrelation fun
tions near the
osmologi
al singularity, and
another with quasi-isotropi
behavior. The �phase transition� between the
two phases is interpreted as the
ondensation of gravitons.
On leave from Landau Institute for Theoreti
al Physi
s, 119940, Mos
ow, Russia.
http://arxiv.org/abs/0704.0354v2
One pessimisti
quotation from the golden era of �nding exa
t solutions of
the Einstein equations whi
h re�e
ted the relations between parti
le theorists
and experts in GR belongs to Ri
hard Feynman. Taking part in the Interna-
tional Conferen
e on Relativisti
Theories of Gravitation at Warsaw, he was
writing to his wife [1℄: �I am not getting anything out of the meeting. I am
learning nothing. ... I get into arguments outside the formal sessions (say,
at laun
h) whenever anyone asks me a question or starts to tell me about his
�work�. The �work� is always: (1)
ompletely un-understandable, (2) vague and
inde�nite, (3) something
orre
t that is obvious and self-evident but worked out
by a long and di�
ult analysis, and presented as an important dis
overy, or (4)
a
laim based on the stupidity of the author that some obvious and
orre
t fa
t,
a
epted and
he
ked for years, is in fa
t false ... (5) an attempt to do something
probably impossible but
ertainly of no utility whi
h, it is �nally revealed in the
end, fails ... or (6) just plan wrong ... Remind me not to
ome to any more
gravity
onferen
es!� Certainly, I am well aware of that the work presented in
this essay
ould belong to the
lass (3) or (5) in the Feynman's
lassi�
ation
(hopefully, not to the
lass (6)!), but I will follow Feynman's own words [1℄:
�We all do it for the fun of it� trying to �nd my fun in identifying some links
whi
h
onne
t the part of the
ommon lore on general relativity named �Exa
t
solutions of the Einstein equations� to the problem of the GR quantization.
Of
ourse, Feynman's interest was in the quantization of GR by applying the
path integral approa
h working so well in QED. Solutions of the Einstein equa-
tions de�ne saddle points of the a
tion
2 S = Sgravity + Smatter of the quantum
gravity with matter. However, the
ontributions of these saddle points into the
partition fun
tion of the theory and �u
tuations near them
Dφmatter exp
(Sgravity + Smatter)
typi
ally have zero measure. In other words, the probability for an almost any
exa
t solution to des
ribe the observable features of the Universe or some parts
of it, to appear somehow from the quantum foam realized near the singularity
is in�nitely small, and the Feynman's anger is absolutely understandable.
Well, almost absolutely... Of
ourse, there are several
lasses of solutions
whi
h will be important for the quantum part of the story, too, and one
an
without mu
h thinking immediately identify some:
1. Attra
tors : among them are Minkowski spa
etime, de Sitter (at least in
the sense of eternal in�ation [2℄) and anti de Sitter spa
etimes (a set of
AdS domains is mostly probably the global attra
tor of GR realized as
low-energy approximation of string theory [3℄); bla
k holes (S
hwarzs
hild,
Kerr, Reissner-Nordström, Kerr-Newman solutions), et
.
From now on, by the quantum theory of gravity we mean e�e
tive QFT of spin 2 �elds
[4℄ (plus matter �elds) � the one whi
h parti
les with energies E ≪ MP test. In this limit,
the e�e
ts of the non-renormalizability may be negle
ted. Although we dis
uss below the
situation whi
h is realized near the
osmologi
al singularity, we limit the dis
ussion to time
s
ales t ≫ tP .
2. General solutions of the Einstein equations. As usual [5℄, a solution of the
Einstein equations is regarded as general if it
ontains su�
ient number
of arbitrary fun
tions of
oordinates. In the
ase of Ri
i-�at spa
etimes,
this number is 4, and is equal to 8 in the presen
e of hydrodynami
matter.
While any non-attra
tor type solution of the Einstein equations de�nes the
saddle point for the path integral (1) whi
h does have a vanishing
ontribution
into the overall partition fun
tion, eventually it well settle down towards an
attra
tor solution due to the e�e
t of
lassi
al perturbations and/or quantum
�u
tuations. The
ontribution of attra
tor type saddle points into the partition
fun
tion (1) is therefore signi�
ant. However, the key word here is �eventually�.
For any non-attra
tor solution it takes a time tcoll before the solution rea
hes
its attra
tor asymptoti
s.
Let us
onstru
t some initial state |Ψ(t = ti)〉 of quantum matter �elds in
a
urved spa
etime and gravitons. The amplitude 〈Ψ(tf )|Ψ(ti)〉 is then de�ned
by the path integral (1)
al
ulated on the
losed S
hwinger-Keldysh
ontour
from t = ti to t = tf and ba
k. Then, if tf ≪ tcoll, the
orresponding attra
tor
saddle point does not give any noti
eable
ontribution into the amplitude.
it is ne
essary to know the evolution of the quantum state |Ψ(t)〉 at time s
ales
t ≪ tcoll, we are for
ed to pay mu
h more attention to the type of saddle
orresponding to general solutions of the Einstein equations.
Certainly, the Einstein equations are hard to solve, and it is possible to �nd
something like their general solution only in physi
ally simpli�ed situations.
As was �rst shown by Belinskii, Khalatnikov and Lifshitz [6℄, asymptoti
ally,
the general solutions of the Einstein equations near the
osmologi
al singularity
have the very same form for an almost arbitrary
hoi
e of the matter
ontent.
This asymptoti
s in the syn
hronous frame
is given by Kasner-like solution
ds2 = dt2 − γαβ(t,x)dxαdxβ , (2)
γαβ(t, x) = t
lαlβ + t
mαmβ + t
nαnβ . (3)
Both Kasner exponents p1, p2, p3 and Kasner axis ve
tors lα, mα and nα are
arbitrary fun
tions of spa
e
oordinates. The Einstein equations provide two
onstraints on the Kasner exponents
p1 + p2 + p3 = 1, (4)
p21 + p
2 + p
3 = 1, (5)
as well as three other
onstraints on arbitrary fun
tions of spa
e
oordinates
present in (3). Taking into a
ount that the
hoi
e of syn
hronous gauge
g00 = 1, g0α = 0 (6)
Of
ourse, the time s
ale tcoll itself is a fun
tional of the initial state |Ψ(t = ti)〉.
Often, it is impossible to
hoose the globally syn
hronous frame of referen
e due to the
limitations set by the
asuality. However, everywhere in the text we dis
uss the physi
s in a
given
asual pat
h.
leaves the freedom to make three-dimensional spa
e
oordinate transformations,
one
an easily see that the total number of arbitrary
oordinate fun
tions in the
Kasner-like solution (2),(3) is equal to 4 as it should be expe
ted for a general
solution of Einstein equations
orresponding to an empty spa
etime.
In the presen
e of the hydrodynami
matter Kasner solution (2),(3) de-
s
ribes asymptoti
behavior of metri
s near the singularity,
sin
e
omponents
of energy-momentum tensor Tik grow slower at t → 0 then the
omponents of
the Ri
i tensor.
Higher order
orre
tions to the Kasner solution (2),(3), i.e.,
higher order terms in the expansion of γαβ(t, x) over powers of t play the role
of perturbations whi
h give rise to the time dependen
e of Kasner exponents pi
as well as Kasner axis ve
tors lα, mα and nα and to well-known BKL
haoti
behavior. Therefore, the BKL solution is simultaneously a universal attra
-
tor for all solutions of the Einstein equations possessing a spa
elike singularity.
It means that no other saddle points
ontribute into the amplitude (1) in the
vi
inity of the
osmologi
al singularity.
In this essay, it will be �rst of all shown that in the presen
e of a s
alar �eld
with potential V (φ) whi
h is exponential and unbounded from below, the general
asymptoti
solution of the Einstein equations is di�erent from the BKL solution
and is quasi-isotropi
[8℄ (while the BKL solution is essentially anisotropi
). In
parti
ular, we will
hoose potential of the form
V (φ) = −|V0|ch (λφ) . (7)
S
alar �eld potentials of this form appear in problems related to gauged super-
gravity models [10℄ and the ekpyroti
s
enario [11℄. The
osmologi
al singularity
realized in su
h theory is of the Anti de Sitter Big Crun
h type. The physi
s
in its vi
inity it is interesting by itself and even more so sin
e this type of
singularity seems to be realized quite often on the string theory lands
ape [3℄.
As in the
ase dis
ussed in [6℄, it is
onvenient to perform all
al
ulations
in the syn
hronous frame of referen
e where g00 = 1, g0α = 0, gαβ = −γαβ ,
α, β = 1 . . . 3, i.e., the spa
etime interval has the form
ds2 = dt2 − γαβ(t, x)dxαdxβ . (8)
Near the hypersurfa
e t = 0 whi
h
orresponds to the singularity, the spatial
metri
omponents behave as
γαβ(t,x) = aαβ(x)t
2q + cαβ(x)t
d + bαβ(x)t
(i,j)
(x)tfij . (9)
With the same pre
ision, one has in the vi
inity of singularity
φ(t,x) = ψ(x) + φ0(x)log(t) + φ1(x)t
f1 + φ2(x)t
f2 + · · · , (10)
Whi
h
orresponds everywhere below to the spa
elike hypersurfa
e t = 0.
If there is a s
alar �eld in the matter
ontent [7℄, BKL solution (3) remains general solution
of the Einstein equations with
hanged Kasner
onstraints (4),(5).
The quasi-isotropi
solution for su
h potentials was �rst found at the ba
kground level in
[9℄, where it was also shown that it is the attra
tor. The goal we pursue in this essay is to
prove that the quasi-isotropi
solution is also general and to understand how its instability
develops with the
hange of the form of the potential.
with dots
orresponding to higher order terms of φ(t,x) expansion in powers
of t. From the Einstein equations one �nds8 that the leading exponents in the
expansions (9) and (10) are de�ned by the expressions
, n = 2, d = 1− q, (11)
ψ(x) = Const, φ0(x) =
, f1 = 1− 3q, f2 = 2− q, (12)
cαα(x) = 2λφ1(x), c
α;β(x) =
1− 2q
1− 3q
φ0φ1,α(x), (13)
P̃ βα (x) + (1 − q)(qbγγ(x)δβα + (1 + q)bβα(x)) =
e−ψλφ2(x), (14)
− (1− q)bαα(x) =
(1− q)φ0φ2(x)−
−ψλφ2(x), (15)
where P̃ βα (x) is the 3-dimensional Ri
i tensor
onstru
ted from
omponents of
the tensor aβα(x) as
omponents of metri
tensor. Higher order terms in the
expansions (9) and (10)
an be self
onsistently
al
ulated by using the Einstein
equations and the orthogonality
ondition
β = δ
α. (16)
One
an immediately �nd from Eq. (16) that the higher order exponents in the
metri
(9) are de�ned by
fij = i+ 2j − (3i+ 2j − 2)q, (17)
where i, j ∈ N. The n term in the metri
expansion
orresponds to i = 0, j = 1
and d term � to i = 1, j = 0. It is easy to see that there is no other exponents
in the expansion (9).
Let us examine the formulae (11)-(15) more
losely and
al
ulate the number
of arbitrary fun
tions present in this solution. First of all, one
an immediately
see that the tensor aβα(x) is not
onstrained by the Einstein equations. It has 6
omponents, and 3 of them
an be made to be equal to 0 by a three-dimensional
oordinate transformation (the remnant gauge freedom of the syn
hronous gauge
(6)). Sin
e this tensor is used for lowering and rising the indi
es and represents
the leading term in the expansion (9), we will identify the term aαβt
ba
kground
ontribution to γαβ(t, x). Furthermore, we see from Eqs. (14),(15)
that bαβ
an be re
onstru
ted from the known tensor aαβ .
The tensor cαβ
ontains three more arbitrary fun
tions of
oordinates. In-
deed, it
an be represented in the form
cβα(x) =
α + Y
;α + Y
Y γγ δ
α + c
(TT)β
α . (18)
Due to the limitations of spa
e we are unable to present the full derivation of the solution
here. It will be given in the forth
oming publi
ation [12℄.
The indi
es of all matri
es are lowered and raised by the tensor aαβ , for example, b
From Eq. (13) one
an see that its tra
e part de�nes the value of φ2(x)
on-
tributing to Eq. (10) and therefore provides one arbitrary fun
tion. Then, three
omponents of the ve
tor
ontribution Yα(x) are �xed, and transverse tra
eless
part c
(TT)β
α (x) provides remaining two arbitrary fun
tions. We also note that
the cαβ term
an be regarded as the leading term perturbation to the ba
kground
ontribution into γαβ . In parti
ular, it
ontains the
ontribution of s
alar per-
turbations (related to the tra
e of the tensor cβα) and tensor perturbations or
gravitons (related to the transverse tra
eless part of the tensor cβα).
The total number of arbitrary fun
tions in the solution (9),(10) is therefore
6, as one may expe
t for the general solution of the Einstein equations with a
s
alar �eld. By analysis similar to [6℄, one may show [9, 12℄ that the
ontri-
butions of other matter �elds into the overall energy-momentum tensor grow
slower at t → 0 than the
ontribution of the s
alar �eld. We
on
lude that
the solution (9), (10) is the general asymptoti
solution of Einstein equations
(with arbitrary matter
ontent) near the
osmologi
al singularity. Similarly to
the BKL solution, the quasi-isotropi
solution is the universal attra
tor for all
solutions of the Einstein equations with s
alar �eld having the potential (7)
and arbitrary additional matter
ontent whi
h possess the time-like singularity.
Again, under
onsidered
onditions, no other saddle points
ontribute into the
amplitude (1) in the vi
inity of the Big Crun
h singularity.
It is instru
tive to understand how exa
tly the transition from the quasi-
isotropi
regime (9),(10) near the singularity to the BKL anisotropi
regime
(3) happens. This transition
an be a
hieved by
hanging the value of λwhile
keeping V0 �xed (or vise versa).
By
onstru
tion, 2q < d = 1− q, i.e., the exponent aαβt2q in the expansion
(9) is leading. With the in
rease of q, the value of d de
reases and when q rea
hes
the
riti
al value qc = 1/3, the
ontributions aαβ(x)t
and cαβ(x)t
into the
expansion of the metri
(9) be
ome of the same order. Similarly, one
an
he
k
that the values of higher order exponents (17) de
rease with the in
rease of q.
In parti
ular, all exponents with di�erent i's and similar j's be
ome of the same
order of magnitude at qc = 1/3. At q > qc = 1/3 the general asymptoti
solution
of the Einstein equations near the singularity is given by Eq. (3) instead of Eq.
In fa
t, what we have just found is relevant for the quantum part of the
story, too, and in a sense is analogous to the spontaneous symmetry breaking
phenomenon in QFTs. Indeed, let us take the theory with a s
alar �eld
(Φ2 − v)2, (19)
set Φ(x, 0) = 0 as an initial
ondition and
ontinuously
hange the value of
the parameter v. At v > 0 the solution Φ(t, x) = 0 of the
lassi
al equations
of motion is perturbatively stable and
orresponds to the true va
uum of the
theory at the quantum level. At ν < 0 the same solution be
omes
lassi
ally
unstable, and Φ(t, x) rea
hes the �true� va
uum value Φ = ±
v during the time
t ∼ 1
log 1
(with the VEV of the operator Φ̂ having similar behavior at the
quantum level). Similar situation is realized in our
ase.
At q < qc = 1/3 the quasi-isotropi
solution (9),(10) is the general solution of
the Einstein equations; it is perturbatively stable by
onstru
tion (without any
limitations on the weakness of the perturbations). At q > qc the quasi-isotropi
solution be
omes perturbatively unstable (perturbations de�ned by cαβ and
higher order terms grow faster than the ba
kground term aαβ at t→ 0).
Vise versa, at q > qc = 1/3 the BKL anisotropi
solution if the Einstein
equations is general in the vi
inity of the
osmologi
al singularity. It is stable
by
onstru
tion with respe
t to arbitrary perturbations and the stability is lost
at q < qc.
This analysis remains valid for the quantum situation
sin
e the
anoni-
al phase spa
e is in one-to-one
orresponden
e with the spa
e of solutions of
lassi
al �eld equations [13℄, and both quasi-isotropi
and BKL solutions are (a)
general and (b) universal attra
tors for other solutions of the Einstein equations
in the vi
inity of the time-like singularity.
The transition from the regime realized at q < 1/3 to the regime q > 1/3
probably
orresponds in the quantum level to the
ondensation of gravitational
perturbations. Indeed, one
an interpret the higher order
ontributions in the
expansion (9) as terms
orresponding to the intera
tion between gravitational
degrees of freedom as well as higher order nonlinearities in the ba
kground. Our
on
lusion is based on the fa
t that at q = qc the spe
trum of the exponents
in the expansion (9) be
omes in�nitely dense. It is also possible to show that
the point of the �phase transition� qc = 1/3
orresponds at the
lassi
al level
to the situation when the
hoi
e of globally syn
hronous frame of referen
e is
impossible near the singularity [12℄.
Let us summarize what have been found in the present essay. We have shown
that in the presen
e of the s
alar �eld with exponential potential unbounded
from below, the general asymptoti
solution of the Einstein equations near
the
osmologi
al singularity has quasi-isotropi
behavior instead of anisotropi
found by [6℄. We have argued that at the quantum level there should exist a
phase transition between the quasi-isotropi
and anisotropi
phases, governed
by the strength of the s
alar �eld potential and interpreted this phase transition
as the
ondensation of gravitational perturbations.
A
knowledgements
I am thankful to A.A. Starobinsky and D. Wesley for the dis
ussions and to
K. Enqvist for making helpful
omments. While
ondu
ting this work, I was
supported by Marie Curie Resear
h training network HPRN-CT-2006-035863.
One important
omment regarding the quantization should be made. The quantum theory
of the s
alar �eld with the potential (7) is ta
hyoni
ally unstable and has neither well-de�ned
asymptoti
|out〉 states, nor 〈out|in〉 S-matrix. However, the S
hwinger-Keldysh 〈in|in〉 S-
matrix is de�ned, and it is possible to make sense of the
orresponding time-dependent theory
[12℄.
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|
0704.0355 | Trigonometric parallaxes of high velocity halo white dwarf candidates | Astronomy & Astrophysics manuscript no. white2v4 c© ESO 2018
November 4, 2018
Trigonometric parallaxes of high velocity halo white dwarf
candidates ?
C. Ducourant1, R. Teixeira2,1, N.C. Hambly3, B. R. Oppenheimer4, M.R.S. Hawkins3, M. Rapaport1, J. Modolo1, and
J.F. Lecampion1
1 Observatoire Aquitain des Sciences de l’Univers, CNRS-UMR 5804, BP 89, 33270 Floirac, France.
2 Instituto de Astronomia, Geofı́sica e Ciências Atmosféricas, Universidade de São Paulo, Rua do Matão, 1226 - Cidade Universitária,
05508-900 São Paulo - SP, Brasil.
3 Scottish Universities Physics Alliance (SUPA), Institute for Astronomy, School of Physics, University of Edinburgh, Royal
Observatory, Blackford Hill, Edinburgh, EH9 3HJ, UK.
4 Department of Astrophysics, American Museum of Natural History, 79th Street at Central Park West, New York, NY 10024-5192,
Received / Accepted
ABSTRACT
Context. The status of 38 halo white dwarf candidates identified by Oppenheimer et al. (2001) has been intensively discussed
by various authors. In analyses undertaken to date, trigonometric parallaxes are crucial missing data. Distance measurements are
mandatory to kinematically segregate halo object from disk objects and hence enable a more reliable estimate of the local density of
halo dark matter residing in such objects.
Aims. We present trigonometric parallax measurements for 15 candidate halo white dwarfs (WDs) selected from the Oppenheimer et
al. (2001) list.
Methods. We observed the stars using the ESO 1.56-m Danish Telescope and ESO 2.2-m telescope from August 2001 to July 2004.
Results. Parallaxes with accuracies of 1–2 mas were determined yielding relative errors on distances of ∼ 5% for 6 objects, ∼ 12%
for 3 objects, and ∼ 20% for two more objects. Four stars appear to be too distant (probably farther than 100 pc) to have measurable
parallaxes in our observations.
Conclusions. Distances, absolute magnitudes and revised space velocities were derived for the 15 halo WDs from the Oppenheimer
et al. (2001) list. Halo membership is confirmed unambiguously for 6 objects while 5 objects may be thick disk members and 4
objects are too distant to draw any conclusion based solely on kinematics. Comparing our trigonometric parallaxes with photometric
parallaxes used in previous work reveals an overestimation of distance as derived from photometric techniques. This new data set can
be used to revise the halo white dwarf space density, and that analysis will be presented in a subsequent publication.
Key words. Astrometry : trigonometric parallax – Dark matter – Galaxy : halo – Star : kinematics – white dwarfs.
1. Introduction
In the last decade interest in the very cool, old white dwarf (WD)
halo population has grown. This interest is motivated by the pos-
sibility that these objects could account for a significant fraction
of the baryonic dark matter of our Galaxy. This idea is in accord
with discussions attempting to explain the microlensing events
in the Large Magellanic Cloud in terms of a halo WD popula-
tion – see, for example, Chabrier et al. 1996 and Hansen 1998.
Alcock et al. 1999 suggested that massive compact halo objects
(MACHOs) make up 20 to 100% of the dark matter in the halo,
with MACHOs having typical mass m ∼ 0.5 M�; more recently,
Calchi Novati et al. 2005 find a similar result from pixel lensing
in the line of sight to M31. Hence, in this scenario the search
for, and direct study of, halo WDs can provide constraints on
the fraction of dark matter in the Milk Way that is attributable to
these objects.
Oppenheimer et al. (2001, hereafter OHDHS) identified 38
high proper motion WDs; from their kinematics, the authors
Send offprint requests to: [email protected]
? Based on observations collected at the European Southern
Observatory, Chile (067.D-0107, 069.D-0054, 070.D-0028, 071.D-
0005, 072.D-0153, 073.D-0028)
concluded that they were members of a halo population. Since
then an intense discussion concerning the status of these objects
has taken place in the literature. A comprehensive review of
this debate is presented in Hansen and Liebert 2003 where the
conclusion is that the OHDHS interpretation is possibly over-
stated, but that complete conclusions are not possible without
further data. Other studies suggest that the disk and “thick disk”
Galactic populations can be used to explain the great majority
of the objects (Reid 2005, Kilic et al. 2005, Spagna et al. 2004,
Crézé et al. 2004, Holopainen & Flynn 2004, Flynn et al. 2003,
Silvestri et al. 2002). The importance of the high velocity WDs
cannot be understated in other contexts (e.g. the star forma-
tion history of the Galaxy, see also Davies, King & Ritter 2002,
Hansen 2003, Montiero et al. 2006). Moreover, several stud-
ies emphasise the importance of obtaining trigonometric par-
allaxes for candidate halo WDs (Bergeron & Leggett 2002,
Torres et al. 2002, Bergeron 2003). This is especially important
for the coolest WDs, whose spectral energy distributions show
remarkable departures from black–body distributions and which
are proving to be difficult to model accurately (Kowalski 2006,
Gates et al. 2004, Saumon & Jacobson 1999, Hansen 1998). In
the presence of such radical changes to the WD spectrum, the as-
sumption of a monotonic photometric parallax relation (e.g. as
2 C. Ducourant et al.: Parallaxes of halo white dwarf candidates
used in OHDHS) could break down and estimates of intrinsic
space velocities could be in error seriously. Furthermore, a recent
paper (Bergeron et al. 2005) concludes that precise distances are
mandatory to derive accurate kinematics and ages for the puta-
tive halo WDs and in order to derive their evolutionary status.
Aiming to clear up this question, in 2001 we started an ob-
serving program with the ESO 1.56-m Danish and ESO 2.2-m
telescopes to measure the trigonometric parallaxes of these stars.
Trigonometric parallax measurements remain the only direct un-
biased distance determination. They are of great importance in
the debate about the status of cool halo white dwarfs because
they are required to derive precise space velocities and ages
which are used for distinguishing between halo and disk mem-
bership. These trigonometric parallaxes lead to the re-calibration
of photometric distances used until now in this debate and allow
analysis of the cool halo white dwarf population with more con-
fidence. Unfortunately, due to limited observing time, only 15
stars on the OHDHS list have been observed to date. However,
this sub-sample provides important insight into the problem.
2. Observations
Astrometric observations of 15 of the OHDHS list of 38 halo
white dwarf candidates were performed at the ESO 2.2-m tele-
scope equipped with the WFI wide–field mosaic camera (with
0.238 ′′/pixel, a field of view FOV = 34′ × 33′, 4 × 2 mosaic of
2k × 4k CCDs), through the ESO 845 I filter. To reduce astro-
metric distortions and other instrumental effects, only data from
chip 51 (with FOV = 8′× 16′) were used in this work; target stars
were centered in the FOV of this chip.
Four epochs of observation were acquired at maximum
parallactic factor in Right Ascension in November 2002, July
2003, November 2003 and July 2004 with a total of 11 nights of
observations. Two parallactic periods (four observations over 1.5
years) are required, at a minimum, for a unique determination
of the parallax and proper motion. Two preliminary observing
runs were performed at the ESO 1.56-m Danish telescope in
July 2001 and July 2002 but the subsequent closure of the
telescope forced the authors to move the program to the ESO
2.2-m telescope. Data acquired at the Danish telescope were not
included in our final analysis to avoid systematic effects due to
the use of two different telescopes.
To minimize differential colour refraction effects (DCR), ob-
servations were performed around the transit of targets with hour
angles of less than 1 hour. Multiple exposures were taken at
each observation epoch to reduce the astrometric errors and to
estimate the precision of measurements. Exposure times varied
from 100 to 600 seconds depending on the magnitude of the tar-
get. Each field was observed from 20 to 35 times.
3. Astrometric Reduction
3.1. Measurements
Frames were measured using the DAOPHOT II package (Stetson
1987), fitting a PSF. The significance level of a luminosity en-
hancement over the local sky brightness which was regarded as
real was set to 7σ. The PSF routine was used to define a stel-
lar point spread function for each frame. Finally we obtained the
(x, y) measured positions, the internal magnitudes and associated
errors of all stars on each frame. There were typically 300 to 600
stars measured on each frame depending on the exposure time.
From these, a selection on the error in magnitude (ERRMAG)
as derived by the DAOPHOT II software was applied. Any ob-
servation with ERRMAG ≥ 0.15m was rejected. Objects fainter
than 1.5m brighter than a given image’s limiting magnitude were
also rejected from the analysis.
3.2. Cross-Identification
For each of the 15 different fields of view, we selected a “master”
or fiducial image from the set of 20 to 35 images. This master
frame for each object had the deepest limiting magnitude and
highest image quality. For each of the other images for a given
target object, the positions of all stars not rejected by the crite-
ria above were then cross–identified to the master image’s star
positions. Objects not detected on three or more frames were
excluded, yielding 100 to 200 stars in common in each field.
Frames containing less than Nmaster/3 stars in common with the
master frame were removed from the solution (where Nmaster is
the number of stars in the master frame). Note that the master
frame is processed in an identical fashion to the other frames
and is not assumed to be free of errors in the parallax solution.
In other words, the fiducial frames are not taken as an error-free
“truth”, but are simply used as a basis for coordinate transfor-
mations and correlation of star positions that comprise the astro-
metric grid used in the solution.
3.3. Differential Colour Refraction
Atmospheric refraction changes the apparent positions of stars in
ground–based observations and depends on the zenith distance
of the observations. For precision astrometry this effect must be
accounted for, because it can be many tens of milliarcsececonds
at even relatively modest zenith distances. In our case, another
effect becomes important as well, because the atmospheric re-
fraction of our target stars will not be identical to that of the
background stars used for our astrometric reference grid. Our
target stars (WDs) and the background stars (typically main–
sequence G or K stars) have different spectral energy distribu-
tions. Therefore, atmospheric refraction will affect them differ-
ently when observed through a given filter bandpass. This is
called a differential colour refraction (DCR) and is known to
cause spurious parallactic motion Monet et al. 1992. DCR can
affect both the Right Ascension (RA) and Declination of the tar-
get as derived with respect to the field stars. Observations in par-
allax programs are planned to maximize the parallactic factor in
RA so the parallax solution for the target will rely heavily on
the RA measured. Therefore the parallax derived is mainly per-
turbed by the DCR effects in RA which are critically dependent
on the zenith distance of a given observation.
We investigated the impact of such effects on the parallax
of white dwarfs through simulations. Using the usual formula
for atmospheric refraction, a blackbody approximation for white
dwarf and background stellar spectra, the Besançon Galaxy
model for background star characteristics (Robin et al. 1994)
and ESO 845 filter limits, we computed the average differen-
tial colour refraction effects between a white dwarf similar to
those of our list with effective temperatures, Teff , in the range
4000 K to 11000 K (Bergeron et al. 2005, Table 2) and a typical
background star (Teff ∼ 5000 K).
We present in Fig. 1 the effects of DCR in RA for white
dwarfs situated at δ = −30◦, covering the range of temperatures
of our targets. Fig. 1 demonstrates that the impact of DCR effects
were always less than 0.5 mas for observations taken with an
hour angle of less than one hour. Therefore, our observations
C. Ducourant et al.: Parallaxes of halo white dwarf candidates 3
were made specifically so that the hour angle never exceeded
one hour, and DCR corrections were not applied in this work.
Fig. 1. DCR effects in RA between a white dwarf of temperature
Teff and a mean background stars (Teff=5000K) at a Declination
of −30◦ (representative of our sample) for various hour angles of
observation. The DCR effects appear to be always lower than 0.5
mas for observations performed at less than 1 hour from merid-
ian which is the case of the present project. The DCR effects are
then negligible compared with other sources of astrometric error
and were not taken into account in this work.
3.4. Impact of Pixel Scale Errors on Parallax
Proper motions (µx, µy) and trigonometric parallax (πxy) of tar-
gets are determined by comparing the (x, y) measurements ex-
pressed in pixels. A scaling factor S f , the image pixel scale, is
applied to πxy to convert pixel measurements into physical units:
π = S fπxy; d(pc) = π−1, with S f expressed in ′′/pixel.
Derivation of the pixel scale can be achieved through a
cross–correlation between the (x, y) positions of stars on a given
master frame to corresponding values of (α, δ) for the subset of
stars that are also in a reference catalogue. Here we used the
2MASS catalogue (Cutri et al. 2003) to determine the orienta-
tion of the master frame on the sky and for the pixel scale deter-
mination. We selected the 2MASS catalogue as a reference cat-
alogue because of its accuracy and density although we note the
absence of proper motion corrections. Nevertheless the epoch
difference between our observations and the 2MASS catalogue
(3 years) would result in negligible corrections to the catalogue
positions with respect to the catalogue errors.
Errors on the scale so determined, resulting from catalogue
random errors, will produce errors in the distance determination
of the target. It is therefore important to quantify the impact of
the catalogue errors onto the distance of the target.
To measure this impact in the present work, we assumed N
reference stars equally spread over a square detector of side A.
The classical equation relating the (x, y) measurements of a stars
on the frame to its standard coordinates X(α, δ),Y(α, δ) in the
tangent plane to the celestial sphere is (with a similar equation
in the Y coordinate)
X = (ax + by + c)1/F, (1)
where (a,b,c) are the unknown “plate” constants and F the focal
length of the telescope (typically the value indicated in the ref-
erence manual). F is expressed in the same units as (x,y) and A
(pixel, mm). It is then easy to show that a fair approximation of
the variance of the estimation of parameter (a) is given by
�cat, (2)
where �cat is the catalogue precision (expressed in radians).
Similar results can be found in Eichhorn & Williams 1963. We
can express the parallax (in radians) as:
πxy, (3)
σ2π = π
σ2a, σπ =
�cat (4)
with F ∼ 13m, A ∼ 0.03m, we evaluate here σπ ∼ 10−4π. The
impact of the error of the catalogue on the parallax of the target
is far below the measurement errors (typically a few milliarcsec-
onds) and are therefore negligible.
3.5. Global Solution: Relative Parallax
The astrometric reduction of the whole set of data of each
field is performed iteratively through a global central overlap
procedure (Hawkins et al. 1998, Eichhorn 1997) in order to
determine simultaneously the position, the proper motion and
the parallax of each object of the field.
The following condition equations are written for each star
on each of the N frames considered (including the master frame).
These equations relate the measured coordinates to the stellar
astrometric parameters:
X0 + ∆X0 + µX(t − t0) + πFX(t) = a1x(t) + a2y(t) + a3 (5)
Y0 + ∆Y0 + µY (t − t0) + πFY (t) = b1x(t) + b2y(t) + b3 (6)
where (X0,Y0) are the known standard coordinate of the star at
the epoch t0 of the master frame, and (x(t), y(t)) its measured
coordinates on the frame (epoch t) to be transformed into the
master frame system. ∆X0, ∆Y0, µX , µY and π are the unknown
stellar astrometric parameters: (∆X0, ∆Y0) yield correction of the
standard coordinates of the star on the master frame, (µX , µY ) are
the projected proper motion in RA ∗cos(δ) and Dec, and π is the
parallax. Coefficients (ai, bi) are the unknown frame parameters
which describe the transformation to the master frame system.
(FX , FY ) are the parallax factors in standard coordinates. The
unknowns of this large over–determined system of equations are
the stellar astrometric parameters of each object, and the trans-
formation coefficients of each of the N frames considered. The
system of equations is singular and therefore the derived solution
is not unique; any solution will depend on the starting point of
the iterations. The usual technique to obtain a particular solution
is to introduce a set of constraints that the solution must satisfy.
In this work we chose to set strictly to zero the mean parallax of
the reference stars.
We used a Gauss–Seidel type iterative method to solve the
set of equations. At the first iteration all stellar parameters are
assumed null, we then computed the plate constants which are
4 C. Ducourant et al.: Parallaxes of halo white dwarf candidates
injected into the system of equations to derive the stellar param-
eters. These results are then used as the starting point of the fol-
lowing iteration. The iterative procedure converges usually at the
second or third iteration. A test of elimination at 3σ is applied to
remove poor observations either in the master frame fit or in the
stellar parameters fit. The stellar parameters fit equations have
been weighted by the mean residual of the master frame fit. This
weighting represents the quality of the measurements. The stars
used for the master frame fit are called here reference stars.
We applied this global treatment to the various observations
of the 15 fields observed and we derived for the targets a proper
motion and parallax with associated variances.
3.6. Conversion from Relative to Absolute Parallax
The parallaxes that we derived for our targets are relative to the
reference stars (for which we used the constraint
π = 0), sup-
posed placed at infinite distance. In fact these reference stars are
at a finite distance from Sun. We must therefore correct the rela-
tive parallax of the target from an estimate of the mean distance
of the reference stars to obtain the absolute parallax of the target.
The choice we made to keep as many reference stars as possible
in our calculation is interesting because statistically faint stars
have smaller parallax and require smaller correction.
There are several ways to estimate the mean distance of
reference stars: statistical methods relying on a model of the
Galaxy; spectroscopic parallax; and photometric parallax. For
the corrections from relative parallax to absolute parallax we
used a statistical method relying on simulations using the
Besançon Galaxy model (Robin et al. 1994) to derive the theo-
retical mean distance of reference stars. A simulation of each ob-
served field was performed, providing catalogues of distance and
apparent magnitude of simulated stars. We computed in these
catalogues mean distances and associated dispersion in magni-
tude bins of 0.2 mag, establishing a table of theoretical distances
with respect to apparent magnitude. Then we considered our ob-
served fields and we computed the weighted mean parallax and
associated dispersion of our reference stars using the theoretical
table. Finally we added this mean parallax of reference stars to
the relative parallax of our target leading to the absolute parallax
of the white dwarfs.
We give in Table 1 the relative to absolute corrections in mil-
liarcseconds as found from the Besançon Galaxy model in each
of the field treated.
4. Results
4.1. Distances of Halo White Dwarf Candidates
We present in Table 2 the proper motions and absolute parallaxes
of the fifteen halo white dwarf candidates as derived from this
work together with their absolute magnitude MV computed using
CCD V magnitudes from Bergeron et al. (2005).
One notices that WD2326–272, LP586–51, LP588–37, and
WD2324–595 are too distant to have a measurable parallax.
Eleven objects are at distances ranging from 19 pc to 90 pc from
the Sun. The parallax errors are about 1–2 mas corresponding
to relative precisions of 5 to 20%. WD2214–390, which is the
closest and brightest object, has a σπ = 2.6 mas. This poor pre-
cision is due to the short exposure time used to avoid saturation
problems and corresponding lower signal–to–noise ratio.
We present in Figs 8 and ?? the positions (empty circles),
their weighted mean (filled circles) and associated error bars at
Table 1. Relative to absolute corrections ∆π and associated
RMS (σ) as found from the Besançon Galaxy model in the
Galactic direction (l,b) together with number of reference stars
(N*) in magnitude interval [Jmin,Jmax].
Target l b ∆π σ N* Jmin Jmax
[◦] [◦] [mas] [mag]
WD2214-390 2.79 -55.37 1.3 0.3 38 13.1 16.2
WD2242-197 40.01 -59.42 1.0 0.3 97 14.0 18.4
WD2259-465 344.30 -60.62 1.1 0.2 83 13.6 18.0
LHS542 72.40 -59.70 1.2 0.3 42 13.4 17.0
WD2324-595 321.83 -54.34 1.1 0.2 62 13.3 17.0
WD2326-272 27.66 -71.06 1.3 0.4 80 14.2 18.7
LHS4033 90.24 -61.96 1.3 0.2 39 14.2 16.5
LHS4041 351.44 -74.66 1.4 0.3 37 13.5 16.2
LHS4042 6.55 -76.61 1.5 0.4 38 13.3 16.6
WD0045-061 118.54 -68.96 1.5 0.3 54 13.5 17.7
F351-50 314.26 -83.50 0.3 0.2 53 14.1 18.1
LP586-51 128.88 -63.30 1.3 0.3 47 14.1 17.4
WD0135-039 149.30 -64.53 1.3 0.2 82 14.4 19.0
LP588-37 150.44 -61.52 1.4 0.2 57 13.6 17.7
LHS147 178.72 -73.56 1.5 0.3 43 13.4 16.8
each epoch of observation, together with the fitted path for the
eleven most significant parallaxes, where π/σπ ≥ 4.
4.2. Comparison with Published Distances
We have compared our results with available data from the lit-
erature, employing both trigonometric and photometric paral-
laxes measured previously. We give in Table 3 the comparison
with published trigonometric parallaxes and in Figure 2 a com-
parison of the parallaxes derived in this work with photomet-
ric parallaxes (from OHDHS, where photometric parallax errors
were 20%). Parameters of a weighted linear fit between pho-
tometric and trigonometric parallaxes are: πtrig = a.πphot + b
with a = 1.08+/-0.08 and b = 3.21+/-1.56 [mas] with a reduced
χ2 =8.06.
Table 3. Comparison of trigonometric parallaxes from this
work (πThiswork) with published data (πext) for LHS 147
(Van Altena et al. 1995), LHS 4033 (Dahn et al. 2004) and LHS
542 (Bergeron et al. 2005).
Target πThiswork πext ∆π
[mas] [mas] [mas]
LHS 542 29.6 +/- 1.8 32.2 +/- 3.7 2.6
LHS 147 14.8 +/- 1.8 14.0 +/- 9.2 –0.8
LHS 4033 30.1 +/-1.8 33.9 +/- 0.6 3.8
Our parallaxes are in excellent agreement with the 3 pre-
viously published trigonometric parallaxes, within the errors
(which are considerably smaller in two cases than published val-
ues). In Fig. 2 one notices a clear systematic tendency of pho-
tometric parallaxes to be underestimated. This overestimation of
OHDHS distances is of importance in the calculation of WD
kinematics and space density.
C. Ducourant et al.: Parallaxes of halo white dwarf candidates 5
Table 2. Proper motion and absolute parallaxes of the fifteen halo white dwarf candidates, where µα∗ = µαcos(δ) and σµ=σµα∗=σµδ ;
π and σπ are the parallax and its precision, Dist the derived distance in parsec and MV the absolute magnitude. No value is given
for Dist and Mv when the parallax is not better than 3 σ. N* is the number of reference stars and Nf the Number of frames. Dphot
is the photometric distance from OHDHS and V is extracted from Bergeron et al. (2005) when available, otherwise (cases marked
by an asterix) it comes from Salim et al. (2004). Note that LHS 4041 is in the OHDHS sample, but is not listed in OHDHS Table 1
(see Table 4 of Salim et al. 2004)
Name α δ Epoch V µα∗ µδ σµ π σπ Dist Mv N* Nf Dphot
[J2000] [yr] [mag] [mas/yr] [mas] [pc] [mag] [pc]
WD2214–390 22 14 34.727 –38 59 07.05 2003.5 15.92 1009 –350 2.9 53.5 2.6 19 14.78 38 28 24
WD2242–197 22 41 44.252 –19 40 41.41 2003.5 19.74 359 +48 3.1 11.1 2.3 90 14.89 97 27 117
WD2259–465 22 59 06.633 –46 27 58.86 2002.9 19.56 402 –153 1.8 22.7 1.3 44 16.49 83 32 71
LHS542 23 19 09.518 –06 12 49.92 2003.5 18.15 –615 –1576 1.8 29.6 1.8 34 15.58 42 33 42
WD2324–595 23 24 10.165 –59 28 07.95 2003.5 16.79 136 –562 1.8 (3.1) 1.5 —- —- 62 25 58
WD2326–272 23 26 10.718 –27 14 46.68 2002.9 ∗19.92 574 –85 2.7 (6.2) 2.4 —- —- 80 17 108
LHS4033 23 52 31.941 –02 53 11.76 2002.9 16.98 631 298 2.5 30.1 1.8 33 14.38 39 26 63
LHS4041 23 54 18.793 –36 33 54.60 2002.9 ∗15.46 21 –662 1.8 13.4 1.5 75 11.10 37 27 59
LHS4042 23 54 35.034 –32 21 19.44 2003.5 17.41 421 –37 2.2 13.9 1.8 72 13.13 38 25 85
WD0045–061 00 45 06.325 –06 08 19.65 2002.9 18.26 111 –668 1.9 30.1 1.9 33 15.59 54 27 44
F351–50 00 45 19.695 –33 29 29.46 2003.5 19.01 1820 –1476 2.1 28.3 1.4 35 16.63 53 34 37
LP586–51 01 02 07.181 –00 33 01.82 2002.9 18.18 350 –118 3.6 (2.4) 2.7 —- —- 47 24 120
WD0135–039 01 35 33.685 –03 57 17.90 2002.9 19.68 456 –180 3.4 13.3 2.9 75 15.26 82 21 146
LP588–37 01 42 20.770 –01 23 51.38 2002.9 ∗18.50 112 –328 3.4 (1.4) 4.5 —- —- 57 17 120
LHS147 01 48 09.120 –17 12 14.08 2002.9 17.62 –115 –1094 2.1 14.8 1.8 68 13.46 43 29 71
Fig. 2. Comparison of parallaxes derived in this work with pho-
tometric parallaxes from OHDHS (errors are assumed 20% for
πphot). Parameters of a weighted linear regression (diagonal line)
between both types of parallaxes are π = 1.08πphot + 3.21 [mas]
with a reduced χ2 = 8.06. The photometric distances are sys-
tematically larger than the astrometric ones.
4.3. Proper Motions
We have compared the proper motions derived here with the
OHDHS proper motions in order to check wether some system-
atic effects could affect our proper motions derived on a 1.5 yr
time span and, as a result, our parallaxes. We present this com-
parison in Fig. 3 and Fig. 4. Error bars are drawn in both co-
ordinates but since the present work has much higher precision
than the photographic astrometry, the error bars in x are not vis-
ible. The slope of a linear regression between proper motions in
α cos(δ) derived in this work with the OHDHS proper motions
is 1.04 ± 0.02 with a reduced χ2 = 3.7. The equivalent linear
fit in proper motions in Declination has a slope of 1.01 ± 0.02
with a reduced χ2 = 0.7. For F351-50 (the largest error bars in
both figures), the accordance in RA and Dec proper motions is
not good. This is due to a known problem of contamination by
a background galaxy of the Schmidt plate measurements used
in the OHDHS work. Nevertheless the accordance is within 2σ.
These comparisons show excellent agreement between both sets
of proper motions, and argue against any systematic effects from
the present work.
4.4. Space Velocities
We derived the Galactic space velocities U, V, W
(Johnson and Soderblom 1987) for the white dwarfs using
the distances and proper motions measured here together with
radial velocities from Salim et al. 2004 (data available for 9
of the 15 white dwarfs treated here). Salim’s observed radial
velocities were corrected for a mean gravitational redshift of
+28km/s as suggested by the authors in their paper except in
the case of the very massive white dwarf LHS4033 were the
correction was taken from Dahn et al. 2004. U is radial toward
the Galactic center, V is in the direction of rotation and W
perpendicular to the Galactic disk. U,V and W were corrected
for the Sun’s peculiar velocity (Mihalas and Binney (1981)).
When no radial velocity was available from other studies, we
assumed Vr = 0 km/s. This approximation is acceptable due to
its minor impact on U,V velocities since the targets are located
6 C. Ducourant et al.: Parallaxes of halo white dwarf candidates
Fig. 3. Comparison of proper motions in RA cos(δ) with the
OHDHS proper motions. Error bars are drawn in both coordi-
nates but since the present work has much higher precision than
the photographic astrometry, error bars in abscissae are not visi-
ble. The slope of a linear regression (dotted line) is 1.04 ± 0.02
indicating good accordance between both proper motion data
sets with a reduced χ2 = 3.7.
close to South Galactic Cap (the effect was investigated in
OHDHS and shown to be negligible).
We present in Figure 5 the distribution of velocities in
the Galactic plane together with the velocity dispersion for
the disk (right most)(1, 2 and 3 σ), thick disk (middle)(1,
2 and 3 σ)( Fuhrmann 2004) and halo (left) (1 and 2σ)
(Chiba and Beers 2000) and in Figure 6 the component of mo-
tion perpendicular to the Galactic plane. These two figures con-
cern the 11 objects with parallax measured at the 4σ level or
better.
In Fig. 5 one notices that 4 of the 11 studied WDs have a
velocity incompatible at the 3σ level with the kinematic of the
disk and of the thick disk and that 6 of them are incompatible at a
2σ level. No star lies within the 1σ ellipse of the disk, primarily
because of selection effects in the original proper motion survey
that OHDHS based is based upon Hambly et al. 2005.
Obviously the choice of the center and dispersions of halo,
thick disk and disk ellipses is critical to classify objects as be-
longing to a particular population. We adopted recent values
which are in in the range of the values cited by Reid 2005
in his review: Disk (Fuhrmann 2004) : (U,V) = (7.7, −18.1)
km/s, (σU , σV ) = (42.6, 22.6) km/s; thick disk (Fuhrmann 2004):
(U, V) = (-18, −63) km/s, (σU , σV ) = (58, 41) km/s; halo
(Chiba and Beers 2000): (U,V) = (0, −180) km/s, (σU , σV ) =
(141, 106) km/s.
5. Discussion
As discussed above, OHDHS sparked a lively debate about
whether stellar remnants contribute to a significant fraction of
the baryonic component of the putative dark matter halo of our
Galaxy. The main criticisms have concerned interpretation, and
we do not address those here. However, the photographic pho-
Fig. 4. Comparison of proper motions in Declination derived
in this work with the OHDHS proper motions. Error bars are
drawn in both coordinates but since the present work has much
higher precision than the photographic astrometry, error bars in
abscissae are not visible. The slope of a linear regression (dotted
line) is 1.01 ± 0.02 indicating a good accordance between both
proper motion data sets with a reduced χ2 = 0.7.
tometry and use of a single photometric parallax relation are also
potential sources of systematic error. Both Salim et al. (2004)
and Bergeron et al. (2005) have shown that the original photom-
etry presented in OHDHS was as accurate as could be expected.
Here, we address the question of the accuracy of photometric
parallaxes directly via trigonometric determination of distances.
In Fig. 2 we compare the trigonometric parallaxes derived
here with the OHDHS photometric parallaxes. Parameters of a
weighted linear regression between both types of parallaxes are
π = 1.08 πphot + 3.21 with a reduced χ2 = 8.06. A clear under-
estimation of photometric parallaxes is visible in this figure with
only one point below the diagonal and three points more than
3σ above the relation. With the usual caveat of small number
statistics, this indicates some level of non–Gaussian scatter, or
at least a mean value for the relation that is not coincident with
π = πphot. The photometric parallax overestimates the distance.
This leads, of course, to an overestimation of tangential space
velocities based on proper motion and distance (as an aside, we
note that the quoted photometric parallax errors of 20% were
conservatively overestimated by OHDHS).
It is interesting to note that the mass distribution of hot
(Teff > 12, 000 K) DA WDs is not Gaussian and has a broad
tail on the high mass side (Należyty et al. 2005). Given that ra-
dius r ∝ m−1/3 for WDs, we would expect photometric paral-
laxes to tend to overestimate rather than underestimate distances
since some of the sample may have higher than average mass,
and correspondingly smaller radii, placing them nearer to the
Sun than typical objects of the same colour. Adding in a sprin-
kling of higher mass WDs with helium–dominated atmospheres
will introduce further systematic overestimation of distances. It
is almost certainly the case that the discrepant photometric par-
allaxes for WD2259–465 and WD0135–039 are caused by these
effects; indeed, this has been shown to be the case for LHS 4033
C. Ducourant et al.: Parallaxes of halo white dwarf candidates 7
Fig. 5. Distribution of velocities in the Galactic plane to-
gether with the velocity dispersion for the disk (right most)(1,
2 and 3 σ), thick disk (middle)(1, 2 and 3 σ)( Fuhrmann 2004)
and halo (left) (1 and 2σ) (Chiba and Beers 2000). Filled
squares correspond to objects with a measured radial velocity
(Salim et al. 2004) while open circles correspond to objects with
no Vr measurement. Only objects with parallax measured at the
4σ level or better are plotted.
which has a mass m ∼ 1.3 M� (Dahn et al. 2004). On the other
hand, the low–mass side of the mass distribution is by no means
perfectly Gaussian (e.g. due to low-mass, helium–core white
dwarfs formed in close binaries). Moreover, any overestimation
in distance leads to a corresponding underestimate of space den-
sity using the 1/Vmax technique. So the interpretation of the re-
sults from this relatively small sub–sample is rather complicated,
and it is only through detailed simulations compared with much
larger samples that significant progress is likely to be made con-
cerning the question of the kinematic population of such objects.
From the comparison of trigonometric and photometric
parallaxes (Fig. 2) we recalibrated photometric distances of
the original OHDHS sample and, using radial velocities from
Salim et al. 2004, we derived their associated recalibrated space
velocities. We present the recalibrated UV plane for the entire
OHDHS sample in Fig. 7.
When compared to Fig. 3 of OHDHS, the number of halo
objects has diminished. From the 38 original OHDHS halo can-
didates, 16 appear compatible with a halo status based on a 2σ
cut with the disk and thick disk velocity distributions (a 3σ cut
would reduce this number to 7), the remaining objects being now
located within the disk and thick disk 2 sigma ellipses. In the lit-
erature there is a large spread of the proposed values to charac-
terise the thick disk and halo populations in terms of kinematics.
For instance in Reid 2005 the velocity dispersions for thick disk
vary from 50 to 69 km/s in the U direction and from 39 to 58
km/s in the V direction. Even the center of velocity ellipsoid
varies from –30 to –63 km/s in the < V > coordinate from one
author to another. All this makes it very difficult to separate ob-
jects into halo and thick disk populations and requires a more
detailed analysis which is beyond the scope of the present paper.
Fig. 6. Component of motion perpendicular to the Galactic
plane (W) as function of
U2 + V2. Only objects with paral-
lax measured at the 4σ level or better and with available radial
velocity (Salim et al. 2004) arre plotted. The vertical line is the
OHDHS
U2 + V2 = 94 km/s cut.
The conclusions of OHDHS about local halo WD density
must be now reanalysed since the volume explored by their sur-
vey has changed (re-calibrated distances) and the number of halo
candidates has also changed. This will be the subject of a forth-
coming paper.
6. Acknowledgements
The authors wish to thank G. Daigne for helpful comments and
CAPES/COFECUB, FAPESP organizations and INR for sup-
porting the project.
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Fig. 8. Observations along the fitted path expressed in mas.
Introduction
Observations
Astrometric Reduction
Measurements
Cross-Identification
Differential Colour Refraction
Impact of Pixel Scale Errors on Parallax
Global Solution: Relative Parallax
Conversion from Relative to Absolute Parallax
Results
Distances of Halo White Dwarf Candidates
Comparison with Published Distances
Proper Motions
Space Velocities
Discussion
Acknowledgements
|
0704.0356 | AMR simulations of the low T/|W| bar-mode instability of neutron stars | AMR simulations of the low T/|W | bar-mode
instability of neutron stars
Pablo Cerdá-Durán, Vicent Quilis, and José A. Font
Departamento de Astronomı́a y Astrof́ısica, Universidad de Valencia, Dr. Moliner
50, 46100 Burjassot (Valencia), Spain
Abstract
It has been recently argued through numerical work that rotating stars with a
high degree of differential rotation are dynamically unstable against bar-mode de-
formation, even for values of the ratio of rotational kinetic energy to gravitational
potential energy as low as O(0.01). This may have implications for gravitational
wave astronomy in high-frequency sources such as core collapse supernovae. In this
paper we present high-resolution simulations, performed with an adaptive mesh re-
finement hydrodynamics code, of such low T/|W | bar-mode instability. The complex
morphological features involved in the nonlinear dynamics of the instability are re-
vealed in our simulations, which show that the excitation of Kelvin-Helmholtz-like
fluid modes outside the corotation radius of the star leads to the saturation of the
bar-mode deformation. While the overall trends reported in an earlier investigation
are confirmed by our work, we also find that numerical resolution plays an impor-
tant role during the long-term, nonlinear behaviour of the instability, which has
implications on the dynamics of rotating stars and on the attainable amplitudes of
the associated gravitational wave signals.
Key words: gravitational waves, hydrodynamics, instabilities, stars: neutron stars:
rotation
PACS: 97.60.Jd, 04.30.-w, 95.30.Lz
1 Introduction
Neutron stars following a core collapse supernova are rotating at birth and
can be subject to various nonaxisymmetric instabilities (see e.g. [1] for a re-
view). Among those, if the rotation rate is high enough so that the ratio of
rotational kinetic energy T to gravitational potential energy W , β ≡ T/|W |,
exceeds the critical value βd ∼ 0.27, inferred from studies with incompressible
Maclaurin spheroids, the star is subject to a dynamical bar-mode (l = m = 2
Preprint submitted to Elsevier 30 July 2021
http://arxiv.org/abs/0704.0356v1
f -mode) instability driven by hydrodynamics and gravity. Its study is highly
motivated nowadays as such an instability bears important implications in
the prospects of detection of gravitational radiation from newly-born rapidly
rotating neutron stars.
Simulations of the dynamical bar-mode instability are available in the litera-
ture, both using simplified models based on equilibrium stellar configurations
perturbed with suitable eigenfunctions [2,3,4,5], and more involved models for
the core collapse scenario [6,7,8,9], and in either case both in Newtonian grav-
ity and general relativity. Due to its superior simplicity the former approach
has received much more attention, notwithstanding that the conclusions drawn
from perturbed stellar models may not be straightforwardly extended to the
collapse scenario.
Newtonian simulations of triaxial instabilities following core collapse were first
performed by [6]. These showed that the bar-mode instability sets in when
β ≫ 0.27 and when the progenitor rotates rapidly and highly differentially.
Such conditions are met when the (artificial) depletion of internal energy to
trigger the collapse is large enough to produce a very compact core for which a
significant spun-up can be achieved. More recently, three-dimensional simula-
tions of the core collapse of rotating polytropes in general relativity have been
performed by [7]. These authors studied the evolution of the bar-mode insta-
bility starting with axisymmetric core collapse initial models which reached
values of β ∼ 0.27 during the infall phase. These simulations showed that the
maximum value of β achieved during collapse and bounce depends strongly
on the velocity profile, the total mass of the initial core, and on the equa-
tion of state. In agreement with the findings from the Newtonian simula-
tions of [6], the bar-mode instability sets in if the progenitor rotates rapidly
(0.01 ≤ β ≤ 0.02) and has a high degree of differential rotation. In addition,
the artificial depletion of pressure and internal energy to trigger the collapse,
leading to a compact core which subsequently spins up, also plays a key role
in general relativity for a noticeable growth of the bar-mode instability.
Whether the requirements inferred from numerical simulations are at all met
by the collapse progenitors remains unclear. As shown by [10] magnetic torques
can spin down the core of the progenitor, which leads to slowly rotating neu-
tron stars at birth (∼ 10 − 15ms). The most recent, state-of-the-art compu-
tations of the evolution of massive stars, which include angular momentum
redistribution by magnetic torques and spin estimates of neutron stars at
birth [11,12], lead to core collapse progenitors which do not seem to rotate
fast enough to guarantee the unambiguous growth of the canonical bar-mode
instability. Rapidly-rotating cores might be produced by an appropriate mix-
ture of high progenitor mass (M > 25M⊙) and low metallicity (N. Stergioulas,
private communication). In such case the progenitor could by-pass the Red
Supergiant phase in which the differential rotation of the core produces a
magnetic field by dynamo action which couples the core to the outer layers
of the star, transporting angular momentum outwards and spinning down the
core. According to [13] about 1% of all stars with M > 10M⊙ will produce
rapidly-rotating cores.
On the other hand, Newtonian simulations of the bar-mode instability from
perturbed equilibrium models of rotating stars have shown that βd ∼ 0.27
independent of the stiffness of the equation of state provided the star is not
strongly differentially rotating. The relativistic simulations of [5] yielded a
value of β ∼ 0.24 − 0.25 for the onset of the instability, while the dynamics
of the process closely resembles that found in Newtonian theory, i.e. unstable
models with large enough β develop spiral arms following the formation of
bars, ejecting mass and redistributing the angular momentum. As the degree
of differential rotation becomes higher Newtonian simulations have also shown
that βd can be as low as 0.14 [14]. More recently [15,16] have reported that
rotating stars with an extreme degree of differential rotation are dynamically
unstable against bar-mode deformation even for values of β of O(0.01).
Given its recent discovery and its potential astrophysical implications for post-
bounce core collapse dynamics and gravitational wave astronomy, we present
in this paper high resolution simulations of such low T/|W | bar-mode in-
stabilities. This work is further motivated in the light of the few numerical
simulations available in the literature. Our main goal is to revisit the simula-
tions by [15] on the low T/|W | bar-mode instability, and particularly to check
how sensitive the onset and development of the instability is to numerical
issues such as grid resolution. To this aim we perform Newtonian hydrody-
namical simulations of a subset of models analyzed by [15] using an adaptive
mesh refinement (AMR) code [17] which allows us to perform such three di-
mensional simulations with the highest resolution ever used. Our simulations
reveal the complex morphological features involved in the nonlinear dynamics
of the instability, where the excitation of Kelvin-Helmholtz-like fluid modes
influences the saturation of the bar-mode deformation. We advance that while
the overall trends found by [15] are confirmed by our work, the resolution
employed in the simulations does play a key role for the long-term behaviour
of the instability and for the nonlinear dynamics of rotating stars, which has
implications on the attainable amplitudes of the associated gravitational wave
signals. We note that we plan to upgrade the existing AMR code to account
for the effects of magnetic fields in order to attempt the current study in a
more realistic setup. The present work is a step towards that goal.
The paper is organized as follows: Section 2 gives a brief overview of the
equations to solve. Their solution is outlined in Section 3 which also contains
the bare details of the AMR code. The results of the simulations are discussed
in Section 4. Finally Section 5 presents our conclusions.
2 Mathematical framework
The evolution of a self-gravitating ideal fluid in the Newtonian limit is de-
scribed by the hydrodynamics equations and Poisson’s equation:
+∇ · (ρv) = 0 (1)
+ (v · ∇)v = −1
∇p−∇φ (2)
+∇ · [(E + p)v] = −ρv∇φ (3)
∇2φ = 4πGρ (4)
where x, v = dx
= (vx, vy, vz), and φ(t,x) are, respectively, the Eulerian
coordinates, the velocity, and the Newtonian gravitational potential. The total
energy density, E = ρǫ + 1
ρv2 , is defined as the sum of the thermal energy,
ρǫ, where ρ is the mass density and ǫ is the specific internal energy, and the
kinetic energy (where v2 = v2x + v
y + v
z). Pressure gradients and gravitational
forces are the responsible for the evolution. An equation of state p = p(ρ, ǫ)
closes the system. We use an ideal gas equation of state p = (Γ − 1)ρǫ with
Γ = 2.
The hydrodynamics equations, Eqs. (1–3), can be rewritten in flux-conservative
form:
∂f(u)
∂g(u)
∂h(u)
= s(u) (5)
where u is the vector of unknowns (conserved variables):
u = [ρ, ρvx, ρvy, ρvz, E] . (6)
The three flux functions Fα ≡ {f , g,h} in the spatial directions x, y, z, respec-
tively, are defined by
f(u) =
ρvx, ρv
x + p, ρvxvy, ρvxvz, (E + p)vx
g(u)=
ρvy, ρvxvy, ρv
y + p, ρvyvz, (E + p)vy
h(u) =
ρvz, ρvxvz, ρvyvz, ρv
z + p, (E + p)vz
and the source terms s are given by
Table 1
Overview of the initial models and results of the simulations. The rows report
the name of the model, the ratio of equatorial-to-polar radii (re/rp), the degree of
differential rotation (Â), the ratio of kinetic to potential energy (T/|W |), the size of
the computational grid (L) and the location of the corotation radius (rc) for the two
resolutions used: high (AMR H) and low (AMR L). In models R1H and R2H the
corotation radius lies outside the star. The real (frequency) and imaginary (growth
rate) parts of the bar-mode σ2 are shown, for the low and high resolution simulation
in comparison with the numerical results and linear analysis by [15]. Note that for
model D3 no linear analysis results are available.
Model D1 D2 D3 R1 R2
re/rp 0.805 0.605 0.305 0.305 0.255
 0.3 0.3 0.3 1.0 1.0
T/|W | 0.039 0.085 0.149 0.253 0.275
L/re 4.06 3.73 3.21 4.25 4.03
rc/re AMR L 0.38 0.47 0.58
AMR H 0.36 0.48 0.56 - -
Re(σ2)/Ω0 AMR L 0.76 0.58 0.41
AMR H 0.81 0.55 0.43 - 0.82
Shibata 0.80 0.60 0.45 0.92 0.75
linear 0.80 0.58 - 0.92 0.75
Im(σ2)/Ω0 AMR L 0.0042 0.0154 0.0200
AMR H 0.0089 0.0190 0.0240 0.0005 0.1960
Shibata 0.009-0.013 0.019-0.021 0.013 <0.002 0.23
linear 0.015 0.021 - <0.002 0.20
s(u) =
, ρvx
− ρvy
− ρvz
. (10)
System (5) is a three-dimensional hyperbolic system of conservation laws with
sources s(u).
3 Numerical approach
For our study of the low T/|W | bar-mode instability we perform high-resolution
simulations of rotating neutron stars using a Newtonian AMR hydrodynamics
code called MASCLET [17]. The implementation of the AMR technique in the
code follows the procedure developed by [18]. The hydrodynamics equations
are solved using a high-resolution shock-capturing scheme based upon Roe’s
Riemann solver and second-order cell reconstruction procedures, while Pois-
son’s equation for the gravitational field is solved using multigrid techniques.
The accuracy and performance of the MASCLET code has been assessed in a
number of tests [17]. We note that the code was originally designed for cos-
mological applications, and here it is applied to simulations of self-gravitating
stellar objects for the first time.
The simulations are performed with two different grid resolutions. The low
resolution grid consists of a box of size L with 1283 zones, yielding a fixed
resolution of L/128. We note that the effective resolution of our coarse grid
is comparable to that used by [15]. Correspondingly, the high resolution grid
consists of a base coarse grid of 1283 cells, and one level of refinement composed
of patches with maximum size of 643 cells (323 coarse cells). This yields a grid
resolution on the finest grid of L/256. This resolution is enough to resolve the
structures simulated, and hence no deeper refinement levels are needed. The
patches are dynamically allocated covering those regions of the star where the
highest resolution is required (highest densities). Typically only one patch is
needed for spheroidal models, and 4-8 in models with toroidal topology. The
use of AMR techniques in our high resolution simulations, allows us to save
about a factor 4 in CPU time and memory with respect to a unigrid simulation
with 2563 cells. No symmetries are imposed in the simulations. To the best of
our knowledge, in the investigations of the bar-mode instability performed by
previous groups, grid resolutions as high as the ones we use here were never
employed.
As customary in grid-based codes [19,20] the vacuum surrounding the star
is filled with a tenuous numerical atmosphere with density ρ/ρmax ≈ 10−12
and zero velocities, ρmax being the maximun density. Every grid cell with
ρ/ρmax < 10
−6 is reset to the atmosphere values. A correct treatment of the
atmosphere is essential for an accurate description of the stellar dynamics and
correct computation of the growth rates of unstable modes. We have checked
that values for the atmosphere higher than those we chose or a free evolution of
the atmosphere altogether, lead to remarkable changes in the mode behaviour,
growth rates, and frequencies. We have also checked that lower values for the
atmosphere do not produce those changes, which ensures that our evolutions
are not affected by the atmosphere values used in the simulations.
4 Results
4.1 Initial data
Differentially rotating stellar models in equilibrium are built according to the
method of [21], and used as initial data for the AMR evolution code. The stars
obey a polytropic equation of state P = KρΓ with index Γ = 2. As [15] the
profile of the angular velocity Ω is given by
(̟/re)2 + Â2
, (11)
where re is the equatorial radius of the star, Ω0 is the central angular ve-
locity, ̟ is the distance to the rotation axis, and  parametrizes the degree
of differential rotation, from  ≪ 1 for highly differentially rotating stars to
 → ∞ for rigidly rotating stars. For comparison purposes these parameters
are chosen as in some of the models of [15], and are summarized in Table 1.
Models labelled D rotate with a high degree of differential rotation, as  = 0.3,
and may therefore be subject to the low T/|W | bar-mode instability. We also
consider models almost rigidly rotating, labelled R, prone to experience the
“classical” bar-mode instability. Labels L and H in the models refer to low
and high resolution respectively.
Following [15] we perturb the initial density profile ρ(0) according to
ρ = ρ(0)
1 + δ
x2 − y2
, (12)
the perturbation of the pressure given by the equation of state accordingly.
A perturbation amplitude δ = 0.1 is used in all our simulations. As we show
below this form of the perturbation excites the l = m = 2 bar-mode. In
addition, grid discretization can leak small amounts of energy to all other
possible modes, which could in principle grow provided they were unstable
and the simulations were carried on for sufficiently long times.
4.2 Stability analysis
To compare with [15] we calculate the distortion parameters η+ and η× (and
η = (η2+ + η
)1/2) defined as
0 0.5 1 1.5 2
Ω / Ω0
0 0.21 0.43 0.65 0.86 1.1 1.3 1.5 1.7 1.9
Fig. 1. Power spectra of Am from m = 1 to m = 8 for model D3H.
Ixx − Iyy
Ixx + Iyy
, η× ≡
Ixx + Iyy
, (13)
where Iij(i, j = x, y, z) is the mass-quadrupole moment
Iij =
dx3ρ xixj . (14)
For the study of the growth rate and interaction of the different angular modes
within the star is useful to calculate the global quantity
dx3ρ(x) e−imϕ, (15)
and Am ≡ Am/A0. We follow the time evolution of modes with m ranging
from 1 to 8. Since our initial equilibrium models are axisymmetric and have
equatorial plane symmetry, all Am are zero initially, but once perturbed all
0 20 40 60 80 100
Fig. 2. Evolution of η for models R1H (upper panel) and R2H (lower panel). Expo-
nential fits to the peaks in the growing phase are overplotted as solid lines.
initial models exhibit a dominant m = 2 component. Assuming that the modes
behave as e−i(σmt−mϕ), the real part of σm can be obtained by Fourier trans-
forming Am. In particular Re(σ2), the bar-mode frequency, can be extracted
from either A2 or η as both represent the same mode. This is the dominant
mode in all our simulations and its frequency and growth rate are given in
Table 1. The latter corresponds to the imaginary part of σ2, which is calcu-
lated fitting an exponential to the peak values of η in the growing phase of
the evolution until the modes saturate. Other modes are also identified in the
simulations for values of Am with lower amplitudes. We have checked that
these modes are harmonics of the l = m = 2 mode so that they follow to good
accuracy the relation σm = mσp, σp being the pattern frequency, calculated
as σp = σ2/2. This is shown for model D3H in Fig. 1 which displays the spec-
trum of Am from m = 1 to m = 8 (in arbitrary units). The vertical dashed
lines in this figure indicate the location of the integer multiples of the pattern
frequency σp, their values indicated on the axis at the top of the figure. Each
spectrum for each mode is normalized to its own maximum for plotting pur-
poses. Note that the lower the mode amplitude the noisier the spectrum and
the less accurate the relation σm = mσp.
For the models of our sample subject to the “clasical” bar-mode deformation
(R1H and R2H), our simulations yield a value of β between 0.253 and 0.275,
in good agreement with the critical value for the onset of the dynamical bar-
mode instability. Model R1H is stable and model R2H is unstable. The growth
rates and frequencies reported in Table 1 agree with those of [15]. Note that
0 100 200 300
Fig. 3. Evolution of η for models D1 (upper panel), D2 (central panel) and D3 (lower
panel). Dashed lines correspond to low resolution and solid lines to high resolution.
Exponential fits to the peaks in the growing phase are overplotted as solid lines.
for model R1H, which is stable, the frequency for the m = 2 mode cannot be
computed. The time evolution of η for these two models is displayed in Fig. 2.
For the unstable model R2H, our simulations show the formation of a bar
which saturates for values of η+ and η× close to 1, i. e. in the full nonlinear
regime.
Fig. 3 shows the time evolution of η for models D in our sample, prone to suffer
the low T/|W | bar-mode instability. Solid lines correspond to high resolution
simulations and dashed lines to low resolution. For all three models the pattern
frequencies σp are such that there exists a corotation radius inside the star,
i.e. a radius at which the bar-mode rotates with the same angular velocity as
the fluid. The location of the corotation radius for all models of our sample
is reported in Table 1. As recently discussed by [22] the existence of such
corotation radius is a potential requirement for the ocurrence of the instability.
As becomes clear from Fig. 3, all models are unstable but grid resolution has
an important effect on the saturation of the instability once the nonlinear
phase has been reached, as well as in the long-term dynamics of the stars.
In the linear phase of models D1H and D2H, the growth rates and frequencies
agree with the results of [15] in both, the numerical simulations and the linear
analysis (see Table 1). In the linear phase of model D3H, our frequencies are
similar to the numerical results of [15], although our growth rates are about
a factor two larger. We emphasize that no results are reported in the linear
analysis for this model in the work of [15], and therefore this discrepancy can
be an effect of the resolution used or of the characteristics of each numerical
code. Increasing resolution leads to similar results in the frequencies but to
higher growth rates.
In the nonlinear phase, models D1 and D3 behave similarly for the two resolu-
tions used (see Fig. 3), and also similarly to the results by [15] (compare with
Fig. 3 of that paper). For model D2 we observe a radical change of behavior in
the nonlinear phase of the mode evolution depending on the grid resolution.
This has implications on the long-term dynamics of the star and, in particu-
lar, on the attainable amplitudes of the gravitational radiation emitted, as we
discuss below.
It is worth mentioning the possibility that the unstable mode at the start of
model D2H might excite some other mode in the corotation band, which could
not otherwise be excited for lower grid resolution. As discussed by [23,24] in
their study of differentially rotating shells, there are many zero-step modes in
the band, so that the whole continuous spectrum could potentially be excited.
In such case these modes would have very slow power-law growth.
For all our models we have checked mass conservation along the evolution. The
worst results are obtained for model D3H, for which mass is conserved within
2.5% error when the instability saturates. At the end of the simulation (after
48 orbital periods and 25000 iterations in the coarsest grid) the error has grown
to only 6%. For all other models mass conservation is even more accurate. Note
that these errors are within the round-off error of the code, and it is not related
to the conservation properties of the numerical scheme itself. For a regular grid
with 1283 cells and a simulation employing 25000 iterations, the accumulated
round-off error (binomial distribution) using single-precision arithmetics, is
about
1283 × 25000 × 10−8 = 0.0023 = 0.23%. Correspondingly, for a 2563
grid (with twice the number of iterations for the simulation) the error is about
0.9%. Taking into account that this error affects the nonlinear evolution of the
system, it is not surprising to have an error at the level of a few percent by
the end of our high resolution simulations, for all conserved quantities.
Figure 4 shows the evolution of Am for model D3 and for m ranging from
1 to 8 for our two resolutions. According to this figure, the only two modes
1e-06
0.0001
0 100 200
1e-06
0.0001
m=1, 3, 5, 7
m=1, 3, 5, 6, 7, 8
Fig. 4. Evolution of |Am| for model D3 with low resolution (top) and high resolution
(bottom). The m = 2 mode is represented with thick solid line, m = 4 with thin
solid line, m = 6 with dashed line, m = 8 with dot-dashed line, and all other odd
m with dotted lines.
relevant for the dynamics of the star are m = 2 and m = 4. All other modes
have smaller amplitudes and play no role in the dynamics. Note that for odd
m modes, the value of the integrated quantity Am, if close to zero, is extremely
sensitive to very small numerical asymmetries, which are induced by the patch
creation scheme of our AMR code. This explains the resolution differences in
the initial values for odd m modes in Fig. 4 (at t = 0 they start off at 10−8
level for the low resolution simulation), although they saturate at the same
value irrespective of the resolution.
An important diagnosis for the accuracy of the results is the location of the
center of mass during an evolution. The round-off error of the numerical code
imposes controlled errors in mass and linear momentum, which results in tiny
displacements of the center of mass. However small (one numerical cell in
our runs) this unphysical displacement may hinder the correct analysis of the
0 100 200 300
1e-06
0.0001
Fig. 5. Effects of the artificial displacement of the center of mass (of only one
numerical cell) on the time evolution of |A1| for model D2H. The thin solid line
shows a fictitiuos evolution resulting from the numerical artifact originated by the
center of mass displacement.
mode growth rates. For this reason all integrated quantities shown in Fig. 4
are computed after correcting for the displacement of the center of mass,
xnew = xold − xCM, in a post-processing stage of the data analysis. Were this
not done, a one-armed m = 1 mode would grow much faster than it should to
bring up fictitious features in the plots. This is shown for model D2H in Fig. 5.
The thick solid line in this figure corresponds to the evolution of the m = 1
mode taking into account the correction for the center of mass displacement,
while the thin solid line is the corresponding evolution of this mode without
the correction.
4.3 Gravitational waves
The growth and saturation of the instability is also imprinted on the gravita-
tional waves emitted. The gravitational waveforms h+ and h× for models D1,
D2, and D3, computed using the standard quadrupole formula, are shown in
Fig. 6. For a source of mass M located at a distance R those waveforms can
be calculated from the dimensionless waveform amplitudes a+ and a× as
h+,× = a+,×
sin2 θ
, (16)
using G = c = 1 units. The resulting chirp-like signal in all the models, partic-
ularly apparent for model D2L, indicates the presence of a bipolar distribution
of mass within the star (see Sec. 4.4).
-0.05
0 100 200 300
Fig. 6. Gravitational waves for models D1 to D3 extracted using the standard
quadrupole formula. Thick (thin) solid lines correspond to low (high) resolution.
Only the dimensionless waveform amplitude a+ is plotted.
As mentioned before, the effects of grid resolution on the evolution of the
nonlinear phase of the bar-mode are imprinted on the gravitational waveforms.
Thick solid lines in Fig. 6 are the waveforms which correspond to the low-
resolution models, and thin solid lines to the high-resolution counterparts.
The evolution of η for model D3, displayed in Fig. 3, shows little deviations
with grid resolution, and this translates into very similar gravitational wave
patterns (bottom panel of Fig. 6), the differences becoming more noticeable in
the nonlinear phase following saturation (Ω0t ≥ 75). For model D1 (top panel),
the differences also become more apparent at later times during the evolution,
in good agreement with the dissimilar behaviour of the matter dynamics in
this model, as encoded in the evolution of η in Fig. 3. As happens for model
D3 the first few cycles of the gravitational waveform, when the mode is still
in the linear phase, are accurately captured for both resolutions.
The major dependence of the waveform on the grid resolution is found for
model D2. Again, the linear phase for the growth of the bar deformation is
Ω 0 t=
70.40.0
Fig. 7. Snapshots of the density, vorticity, and specific angular momentum, for
model D3H, at three representative instants of the evolution. All snapshots show
slices of the stars in the equatorial plane. Quantities are normalized as follows:
max, rew
s , and l
ϕ/(rev
s ), where v
s is the initial velocity at the surface
of the star.
accurately captured irrespective of the resolution (and agrees with the per-
turbative results of [15]). This is signalled in the perfect overlapping of both
gravitational waveforms during the first three cycles (see the middle panel of
Fig. 6). However, the different nonlinear dynamics of the bar-mode deforma-
tion for this model, shown in the middle panel of Fig. 3, is severely imprinted
on the gravitational waveform. Model D2H emits gravitational waves which
have roughly one order of magnitude smaller amplitude than those computed
for the corresponding low resolution model.
4.4 Morphology
We next describe the morphological features encountered during the evolution
of some representative models. Fig. 7 shows three snaphsots of the evolution
of model D3H for the density (top), the azimuthal component of the vorticity,
~wϕ = (∇ × ~v)ϕ (middle), and the specific angular momentum, ~l = ~r × ~v
(bottom). From left to right the snapshots correspond to the initial time (Ω0t =
0), a time when the bar-mode instability is growing (Ω0t = 33.6), and the time
when the instability saturates (Ω0t = 70.4). Only the equatorial plane of the
stars is shown in all these plots. Animations of all simulations performed are
available at www.uv.es/∼cerdupa/bars/. We note that our AMR code is able
to dynamically place patches (e. g. between 4 and 8 in the D3H model) and
evolves the system with continuous matching between patches, as exemplified
in Fig. 7.
The evolution of model D3H shows that as the m = 2 mode grows the star
develops an ellipsoidal shape which remains spinning beyond saturation. Since
the low β m = 2 mode saturates at lower values (η ∼ 0.1) than the classical
bar-mode instability (η ∼ 1), no clear bars are visible in the density plot. At
late times (Ω0t > 100) a “boxy” structure becomes apparent as the m = 4
mode has grown to almost similar amplitude as the m = 2 mode (see anima-
tions and Fig. 4). No other global features can be seen, consistent with the
fact that |Am| ≪ 1 for all modes other than m = 2 and 4. The vorticity plot
shows that the m = 2 mode at Ω0t = 33.6 adopts the form of a two-armed
spiral winding up around the central parts of the star. As the mode begins to
saturate (Ω0t = 70.4) the spirals break apart into the outer layers in a turbu-
lent flow reminiscent of the (shear) Kelvin-Helmholtz instability, and shock as
they reach the atmosphere. These trends are also visible in the specific angular
momentum plot.
The presence of a corotation radius, at r/re = 0.56 for model D3H, seems
to play a role in the growth and saturation of the instability, in agreement
with the recent findings of [25]. As the bar-mode grows, pressure waves carry
angular momentum outside the corotation radius, which is deposited in the
outer layers of the star. This excites Kelvin-Helmholtz-like instabilities in the
fluid that break the mode outside the corotation radius. When this happens the
m = 2 instability stops growing and no more angular momentum is extracted.
Figure 8 shows late-time snapshots of the equatorial plane distribution of the
density perturbation, i.e. (ρ−ρ(0))/ρ(0)max, for models D2H and D3H. The times
are chosen well inside the nonlinear and saturation phase of the instability.
This figure helps to interpret the mode dynamics and its saturation along the
lines mentioned before: During the evolution the density perturbations are
shed in waves from the center towards the outer layers of the star. At late
times, when the instability saturates, such shedding stops, and the density
Fig. 8. Snapshots of the density perturbation at the equatorial plane for models D2H
and D3H. The white solid curves indicate the location of the corotation radius. The
white dashed boxes indicate the location of the patches for model D3H.
perturbation reaches the largest values outside the corotation radius (depicted
with white solid lines in Fig. 8), for either model.
We note in passing that the corotation radius in all our high resolution models
lies well inside the outer boundary of the finest box set up by the AMR
refinement pattern. (see, e.g. the white dashed boxes depicted in the right
panel of Fig. 8 indicating the location of the AMR patches for model D3H)
This rules out the possibility of a numerical artifact resulting from the patch
creation scheme of our AMR code being the cause for the different long-term
evolution between low and high resolution models, particularly noticeable for
model D2 in Fig. 3.
Finally, Fig. 9 shows a comparison between models D2L and D2H at Ω0t = 101
(i.e. well within the nonlinear phase), to highlight the effects of the numer-
ical resolution on the morphology. From top to bottom this panel shows a
schlieren plot (|∇ log ρ|) , ~wϕ, and ~l. The resolution differences in the evolu-
tion of model D2 become apparent from this figure. In particular, the “boxy”
structure becomes much more clearly visible in the low resolution simulation
(D2L), indicating an excessive growth rate of the m = 4 mode. The presence
of pressure waves is emphasized in the schlieren plot, very accurately captured
in model D2H. Those waves, once the flow is driven to turbulence past the
corotation radius, redistribute the angular momentum in the outer layers of
model D2L in a much more pronounced way than for model D2H.
Ω 0 t=
D2L D2H
Fig. 9. Resolution comparison between models D2L and D2H once the instability
has saturated. Only slices of the stars in the equatorial plane are shown.
5 Summary and outlook
We have presented AMR high-resolution simulations of the low T/|W | bar-
mode instability of extremely differentially rotating neutron stars. Our main
motivation has been to revisit the simulations by [15] on such instability,
assessing how sensitive the onset and development of the instability is to
numerical issues such as grid resolution. We have addressed the importance
of a correct treatment of delicate numerical aspects which may spoil three-
dimensional simulations in (Cartesian) grid-based codes, always hampered by
insufficient resolution, namely the handling of the low-density atmosphere sur-
rounding the star, the correction for the center of mass displacement, and the
mass and momentum conservation properties of the numerical scheme. Our
simulations have revealed the complex morphological features involved in the
nonlinear dynamics of the instability. We have found that in the nonlinear
phase of the evolution, the excitation of Kelvin-Helmholtz-like fluid modes
outside the corotation radii of the stellar models leads to the saturation of the
bar-mode deformation. While the overall trends reported in the investigation
of [15] are confirmed by our work, the resolution used to perform the simu-
lations may play a key role on the long-term behaviour of the instability and
on the nonlinear dynamics of rotating stars, which has only become apparent
for some specific models of our sample (namely model D2). This, in turn, has
implications on the attainable amplitudes of the associated gravitational wave
signals.
The work reported in this paper is a first step in our ongoing efforts of study-
ing the dynamical bar-mode instability within the magnetized core collapse
scenario.
Acknowledgements
The authors thank Harry Dimmelmeier, Nick Stergioulas, and Anna Wats
for useful comments. Research supported by the Spanish Ministerio de Edu-
cación y Ciencia (MEC; grants AYA2004-08067-C03-01, AYA2003-08739-C02-
02, AYA2006-02570). VQ is a Ramón y Cajal Fellow of the Spanish MEC.
Computations performed at the Servei d’Informática de la Universitat de
València (CERCA-CESAR).
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Introduction
Mathematical framework
Numerical approach
Results
Initial data
Stability analysis
Gravitational waves
Morphology
Summary and outlook
References
|
0704.0357 | Evolutionary games on minimally structured populations | Evolutionary games on minimally
structured populations
Gergely J. Szöllősi∗ and Imre Derényi†
Biological Physics Department Eötvös University, Budapest
Abstract
Population structure induced by both spatial embedding and more general networks of interaction, such
as model social networks, have been shown to have a fundamental effect on the dynamics and outcome of
evolutionary games. These effects have, however, proved to be sensitive to the details of the underlying
topology and dynamics. Here we introduce a minimal population structure that is described by two distinct
hierarchical levels of interaction, similar to the structured metapopulation concept of ecology and island
models in population genetics. We believe this model is able to identify effects of spatial structure that
do not depend on the details of the topology. While effects depending on such details clearly lie outside
the scope of our approach, we expect that those we are able to reproduce should be generally applicable
to a wide range of models. We derive the dynamics governing the evolution of a system starting from
fundamental individual level stochastic processes through two successive meanfield approximations. In
our model of population structure the topology of interactions is described by only two parameters: the
effective population size at the local scale and the relative strength of local dynamics to global mixing. We
demonstrate, for example, the existence of a continuous transition leading to the dominance of cooperation
in populations with hierarchical levels of unstructured mixing as the benefit to cost ratio becomes smaller
then the local population size. Applying our model of spatial structure to the repeated prisoner’s dilemma
we uncover a novel and counterintuitive mechanism by which the constant influx of defectors sustains
cooperation. Further exploring the phase space of the repeated prisoner’s dilemma and also of the “rock-
paper-scissor” game we find indications of rich structure and are able to reproduce several effects observed
in other models with explicit spatial embedding, such as the maintenance of biodiversity and the emergence
of global oscillations.
PACS numbers: 87.10.+e 87.23.-n
∗[email protected]; angel.elte.hu/˜ssolo
†[email protected]; angel.elte.hu/˜derenyi
http://arxiv.org/abs/0704.0357v3
mailto:[email protected]
angel.elte.hu/~ssolo
mailto:[email protected]
angel.elte.hu/~derenyi
I. INTRODUCTION
The dynamics of Darwinian evolution is intrinsically frequency dependent, the fitness of in-
dividuals is tightly coupled to the type and number of competitors. Evolutionary dynamics acts,
however, on populations, not individuals and as a consequence depends on not only population
composition, but also population size and structure. Evolutionary game theory came about as the
result of the realization that frequency dependent fitness introduces strategic aspects to evolution
[1, 2, 3]. More recently the investigation of the evolutionary dynamics of structured populations,
where individuals only compete with some subset of the population, e.g. their neighbors in space
or more generally in some graph [4, 5], has lead to the recognition that the success of different
strategies can be greatly influenced by the topology of interactions within the population. Funda-
mental differences were found – compared to well-mixed populations, where individuals interact
with randomly chosen partners – in models that describe the evolution of cooperation (variants of
the prisoner’s dilemma game [4, 6, 7, 8, 9]) or deal with the maintenance of biodiversity in the
context of competitive cycles (variants of the rock-paper-scissors game [3, 10, 11, 12, 13, 14]).
In order to investigate the coevolutionary dynamics of games on structured populations the
full set of connections between a potentially very large number of individuals must be specified.
This is only possible by reducing the number of degrees of freedom considered, either through
postulating a highly symmetrical (such as lattices [4, 8, 16, 17, 18, 19, 20]) or fundamentally
random connection structure (such as some random graph ensemble [21, 22]). The question of how
one goes about the task of reducing the number of degrees of freedom – of choosing the relevant
parameters to describe the population structure constrained to which individuals undergo evolution
– is not trivial. Both the explicit spatial as well as the random graph ensemble approach have clear
precedents in condensed matter physics and network theory, respectively. It is not, however, clear
which – if either – approach best describes natural populations of cyclically competing species or
societies composed of individuals playing the prisoner’s dilemma game.
As an example let us consider colicin producing bacteria, that play the so called ”rock-paper-
scissors” (RPS) game (for details see below). This system has recently been the subject of two
experimental studies aimed at demonstrating the role of structured populations in the maintenance
of diversity. In the first study [10, 11] bacteria were cultured in vitro in Petri dishes, effectively re-
stricting competition between bacteria to neighbors on the (2D) Petri dish surface (Fig.1 top left),
while in the second experiment [12] in vivo bacterial colonies were established in co-caged mice
and their development was subsequently followed. In the case of the first experiment the analogy
with explicit 2D spatial embedding (present by construction) is clear (Fig.1 bottom left). The pop-
ulation structure of the second experiment is, however, clearly different. The bacteria in individual
mice can be readily considered as locally well-mixed populations, the coevolutionary dynamics
of which reduces in the standard meanfield limit to a system of non-linear differential equations
(the adjusted replicator equations [24]). As the experiments show, however, migration of bacteria
between mice may also occur – resulting in the observed cyclic presence of the three strains in
individuals. There are two distinct scales of mixing present in the system. Bacteria within each
mice compete with each other forming local populations – an unstructured neighborhood com-
posed of individual bacteria, while also being exposed to migrants from mice with whom they
share the cage, together forming a global population – an unstructured neighborhood composed of
individual local populations (Fig.1 top and bottom right). This setup is referred to in the ecology
literature – albeit in significantly different contexts – as a ”structured metapopulation” [26, 27]
where structured here refers to the detailed consideration of the population dynamics of the indi-
vidual populations (often called ”patches”) comprising the metapopulation and is also related to
the finite island models of population genetics [28].
The above example of co-caged mice is not unique, we may readily think of other ecological or
sociological examples where an approximation with hierarchical scales of mixing with no internal
structure can be relevant (such as human societies with two distinct scales of mixing present, the
first within individual nations the between them at an international level). We have, also, recently
used a similar approach to construct a model of genetic exchange among bacteria of the same
species (the bacterial equivalent of sex) with which we were able to take into account the effects of
spatial and temporal fluctuations in a manner that can explain the benefit of such genetic exchange
at the level of the individual [31].
In this paper we construct a hierarchical meanfield theory where the two distinct (i.e. local and
global) scales of mixing are each taken into account in terms of two separate meanfield approx-
imations and fluctuations resulting from finite population size on the local scale of mixing are
also considered. We subsequently explore the similarities and differences between this and other
models of structured populations in the case of the ”rock-paper-scissors” and prisoner’s dilemma
games. Through these examples we suggest that our approach allows the separation of the effects
of structured populations on coevolutionary dynamics into effects which are highly sensitive to
and dependent on the details of the topology and those which only require the minimal structure
Explicit spatial embeding Two distinct scales of mixing
Petri dish experiment Co-caged mice experiment
FIG. 1: (Color online) In the colicin version of the RPS game, strains that produce colicins (red/dark
grey) kill sensitive (green/light grey) strains, that outcompete resistant (blue/black) strains, that outcom-
pete colicin producing strains (toxin production involves bacterial suicide). Experiments [10] show that
colicin-producing strains cannot coexist with sensitive or resistant strains in a well-mixed culture, yet all
three phenotypes are recovered in natural populations. Two recent experiments have examined the role of
population structure in the maintenance of diversity among colicin-producing bacteria. In the first [10] in
vitro colonies were established on an agar substrate in Petri dishes, a setup which effectively limits compe-
tition to neighbors on the petri dish in analogy with explicit spatial embedding in 2D. In the second [12] in
vivo colonies were established in the intestines of co-caged mice, a setup which has two distinct scales of
mixing, with no explicit structure on either scale.
present in our approximation and can consequently (in terms of sensitivity to the details of the
topology) be considered more robust.
II. HIERARCHICAL MEANFIELD THEORY FOR TWO DISTINCT SCALES
Let us consider an evolutionary game between d types (strategies) described by the d×d payoff
matrix A with elements αkj . Assuming finite and constant population size, natural selection can
be described at the level of the individual by the so called the Moran process [30], during which at
each time step an individual is selected randomly from the population to be replaced (death) by the
offspring of an individual that is chosen proportional to its fitness to reproduce (birth). This models
a population in equilibrium, where the time scale of the population dynamics is set by the rate at
which ”vacancies” become available in the population. The fitness of each individual depends on
the payoff received from playing the game described by A with competitors (an individual of type
k receiving a payoff αkj when playing with an individual of type j). In well-mixed populations,
individuals can be considered to come into contact (compete) with equal probability with any
member of the population excluding themselves – this allows one to calculate the fitness of an
individual of type k in a meanfield manner, yielding
πk = πbase +
αkj(nj − δkj)
N − 1
, (1)
where nk is the number of individuals of type k in the population,
k=1 nk = N is the size of
the population, πbase is some baseline fitness and the Kronecker delta symbol δkj is equal to unity
if k = j and is zero otherwise. From this we may calculate the transition probabilities of our
stochastic process, i.e., the probability of an individual of type i being replaced by an offspring of
an individual of type k is given by
Tik =
, (2)
where π̄ =
k=1 πknk/N . The state of any population is completely described by the frequency
of the different strategies xk = nk/N . Due to the normalization
k=1 xk = 1, the values of xk
are restricted to the unit simplex Sd [3]. For d = 2 this is the interval [0, 1], S3 is the triangle with
vertices {(1, 0, 0), (0, 1, 0), (0, 0, 1)} while S4 is a tetrahedron etc.
As Traulsen et al. have recently shown [24, 25] for sufficiently large, but finite populations
the above stochastic process can be well approximated by a set of stochastic differential equa-
tions combining deterministic dynamics and diffusion (population drift) referred to as Langevin
dynamics:
ẋk = ak(x) +
ckj(x)ξj(t), (3)
where the effective deterministic terms ak(x) are given by
ak(x) =
(Tjk − Tkj) = xk
πk(x)− π̄(x)
π̄(x)
, (4)
ckj(x) are effective diffusion terms, that can also be expressed in terms of the transition probabil-
ities as described in [25], and ξj are delta correlated 〈ξk(t)ξj(t′)〉 = δkjδ(t − t′) Gaussian white
noise terms. As N → ∞ the diffusion term tends to zero as 1/
N and we are left with the
modified replicator equation.
In the context of our hierarchical mixing model the topology of connections can be described
by two parameters, the populations size at the local scale of mixing N , and a second parameter
µ, which tunes the strength of global mixing relative to the local dynamics. We take into account
the second (global) scale of mixing – mixing among local populations – by introducing a modified
version of the Moran process. In the modified process a random individual is replaced at each
time step either with the offspring of an individual from the same population (local reproduction)
or with an individual from the global population (global mixing). This is equivalent to considering
the global population to be well-mixed at the scale of local populations.
Let us consider a global population that is composed of M local populations of size N . In each
local population vacancies become available that local reproduction and global mixing compete to
fill. In any local population l the probability of an individual of some type k filling a new vacancy
due to local reproduction must be proportional to the number of individuals of type k multiplied by
their fitness i.e. πlkn
k, where we consider π
k to be determined only by interactions with individuals
in the same local population according to equation (1). To describe the tendency of individuals of
some type k in local population l to contribute to global mixing we introduce the parameters σlk.
The choice of appropriate σlk depends on the details of the global mixing mechanism, for systems
where only the offspring of individuals mix globally it is proportional to the fitness of a given type,
while for mechanisms such as physical mixing, by e.g. wind or ocean currents, it may be identical
for each type. Irrespective of the details, however, the probability of an individual of some type
k filling in a new vacancy due to global mixing should be proportional to the global average of
the number of individuals of type k multiplied by their mixing tendency, which we denoted as
〈σknk〉 =
l=1 σ
k/M , and the strength of global mixing µ. These consideration lead to the
new transition probabilities:
T̂ lik =
k + µ〈σknk〉
k=1(π
+ µ〈σknk〉)
k + µ〈σknk〉
N(π̄l + µ ¯〈σ〉)
, (5)
where π̄l =
k/N and
¯〈σ〉 =
〈σknk〉/N .
We have found that the results presented below are qualitatively the same for both the fitness
dependent choice of σlk = π
k and the fitness independent choice of σ
k = 1. Therefore, in the fol-
lowing we restrict ourselves to the somewhat simpler fitness independent choice of σlk = 1, which
can be considered to correspond to some form of physical mixing mechanism. The transition
probabilities (5) then reduce to:
T̄ lik =
π̄l + µ
π̄l + µ
. (6)
We can see that after a vacancy appears either local reproduction occurs, with probability π̄l/(π̄l+
µ), or global mixing, with probability µ/(π̄l + µ). From (6) we may derive the Langevin equation
describing the coevolutionary dynamics of population l from the
ẋlk = âk(x
l, 〈x〉) +
ĉkj(x
l, 〈x〉)ξj(t), (7)
with the modified deterministic terms given by
âk(x
l, 〈x〉) = x
k(πk(x
l)− π̄(xl)) + µ(〈xk〉 − xlk)
π̄(xl) + µ
, (8)
where the vector 〈x〉 =
l=1 x
l/M with components 〈xk〉 =
l=1 x
k/M describes the fre-
quencies of the individual types in the global population and the diffusion terms ĉ(xl, 〈x〉) can be
expressed in terms of the modified transition probabilities T̂ lik as above.
Equations (7) describe the coevolutionary dynamics of the global population through the cou-
pled evolution of the {x1, . . . ,xM} local populations. In the limit of a large number of local popu-
lations (M → ∞) the distribution of the local populations over the space of population states (the
simplex Sd) is described by a density function ρ(x) that is normalized over Sd, i.e.,
ρ(x) = 1.
The time evolution of ρ(x) follows a d−1 dimensional advection-diffusion equation – the Fokker-
Planck equation corresponding to eq. (7):
ρ̇(x) = −∇{â(x, 〈x〉)ρ(x)}+ 1
b̂(x, 〈x〉)ρ(x)
, (9)
with the global averages 〈xk〉 =
xkρ(x) coupled back in a self-consistent manner into the
deterministic terms âk(x, 〈x〉) and the diffusion matrix b̂kj(x, 〈x〉) =
i=1 ĉki(x, 〈x〉)ĉij(x, 〈x〉).
For large local populations (N → ∞) the diffusion term vanishes as 1/N .
The above advection-diffusion equation (9) presents an intuitive picture of the coevolutionary
dynamics of the population at a global scale. We can see that local populations each attempt to
follow the trajectories corresponding to the deterministic replicator dynamics, while under the
influence of two additional opposing forces: (i) global mixing, which attempts to synchronize
local dynamics and (ii) diffusion resulting from finite population size effects, which attempts to
smear them out over the simplex. The strength of these forces are tuned by two parameters µ and
N , respectively.
If, further, the effects of synchronization are irrelevant, as for example in the case of populations
where selection is externally driven by independent environmental fluctuations, we may replace
the global population average with the time average of any single population. This is the approach
we used in our study of genetic mixing in bacteria [31].
During our numerical investigations we found solving the advection-diffusion equation (9) nu-
merically challenging, particularly in the N → ∞ limit. We resorted instead to solving the coupled
Langevin equations (7) for large M = 104 − 105 to simulate the time evolution of ρ(x).
III. COOPERATION IN POPULATIONS WITH HIERARCHICAL LEVELS OF MIXING
The evolution of cooperation is a fundamental problem in biology, as natural selection under
most conditions favors individuals who defect. Despite of this, cooperation is widespread in na-
ture. A cooperator is an individual who pays a cost c to provide another individual with some
benefit b. A defector pays no cost and does not distribute any benefits. This implies the payoff
matrix
b− c −c
, (10)
where b is the benefit derived from playing with a cooperator while c is the cost for cooperation.
From the perspective of evolutionary game theory, which equates payoff with fitness, the apparent
dominance of defection is simply the expression of the fact that natural selection a priori selects
for fitness of individuals and not the fitness of groups.
Defection dominates cooperation in any well-mixed population [3]. Population structure in-
duced by spatial structure [4, 18] and more general networks of interactions [21, 22, 23]) has,
however, been found to facilitate the emergence and maintenance of cooperation. The mecha-
nism responsible, termed spatial, or more generally, network reciprocity[32] depends strongly on
biased infux
movement of global average
density of
local populations
drift
deterministic dynamics
x = 0 x = 1
ρ(x, t)
ρ(x, t + ∆t)biased influx
0 0.1 0.2 0.3 0.4 0.5 0.6 0.7 0.8 0.9 1
b/c = 95, 97, 100, 111, 125
1000
10 100 1000
d b/c
60 70 80 90 100 110
FIG. 2: a In an infinitely large well-mixed population evolutionary dynamics is deterministic and leads
to the extinction of cooperators as average fitness monotonically declines. The only stable fixed point
corresponds to the point where the fraction of cooperators is zero (x = 0). To understand qualitatively the
mechanism favoring cooperation in hierarchically mixed populations let us consider some density of local
populations (ρ(x, t)) that is symmetric around its mean at time t. Due to global mixing all local populations
are being driven toward the global average. Due to the influx bias, however, populations with a lower than
average number of cooperators will be driven stronger (faster) than those on the other side of the average.
Examining the density of local populations at some time t+∆t, this results in a net movement of the global
average toward a larger fraction of cooperators. This is, of course, opposed by local reproduction that favors
an increase in the number of defectors. For the global average to keep moving toward a higher number of
cooperators and eventually to keep balance with local reproduction bias a density of local population with
finite width is needed over which the effect of the influx bias can exert itself. It is drift caused by local
population size that maintains this finite width, and this is the reason that the b/c threshold above which
cooperation dominates depends on local population size. b Stationary density of local populations ρ(x) for
different values of b/c with N = 100, µ = 0.1. c Transition toward a global dominance of cooperation for
µ = 10. (triangles), µ = 1 (crosses), µ = 0.1 (squares), µ = 0.01 (circles) with N = 100. The critical value
of b/c depends only weakly on µ changing by 20% over four orders of magnitude d Critical values of b/c
as a function of N for different values of µ (notation as before). The dashed line corresponds to b/c = N .
The critical b/c values were determined by numerically finding the inflection point of the transition curves.
M = 103 was used throughout.
the details of local topology. In particular, it seems that lattice like connectivity structures where
three-site clique percolation occurs [17] and more general interaction graphs where the degree of
nodes k does not exceed the ratio of benefit to cost (i.e. k < b/c) [22] are required for cooperation
to be favored.
Examining the effects of hierarchical mixing on the evolutionary dynamics of cooperation we
found that a sharp, but continuous transition leads to the dominance of cooperation as the benefit
to cost ratio becomes smaller then the local population size, i.e. b/c < N . If the benefit to cost
ratio is larger then the local population size the global population is dominated by defectors. The
mechanism leading to the dominance of cooperation arises due to the competition between local
reproduction and global mixing. In local populations with lower average fitness – larger number of
defectors – the influx of individuals from the global scale will be larger than in local populations
with higher average fitness (cf. eq. (6) where the relative strength of the two terms on the left hand
side depends on the sum of the average fitness of population l and µ). The crucial ingredient for
cooperation to be successful is population drift introduced by finite local population size. It is bi-
ased influx coupled with drift that can result in cooperation being favored in the global population
(Fig 2.).
IV. THE RPS GAME
To explore the effects of hierarchical mixing in the context of games with three strategies we
first turn to the case of the so called ”rock-paper-scissors” (RPS) game. In the original popular
version of the game two players are afforded the chance to simultaneously display either rock
(fist), paper (flat hand) or scissors (two fingers). If player one displays a flat hand while player
two displays a fist, player one wins as paper wraps rock. Similarly scissors cut paper, and rocks
smashes scissors. Several examples of this game have been found in nature (e.g. among lizards
[33] ), but it is bacteria that have received the most experimental and theoretical attention.
In ecology the often high diversity among microbial organisms in seemingly uniform envi-
ronments, referred to as the ”paradox of the plankton”, has been difficult to understand. Several
models based on spatially explicit game theoretical models have been proposed to explain this di-
versity [10, 11, 13, 14]. These models are all variants of the RPS game played by colicin producing
bacteria. Colicins are antibiotics produced by some strains of Echerichia coli. In experiments (see
Fig.1) typically three strains are used: colicin producing (C), sensitive (S) and resistant (R). The
(a) (b) ALLD
µ = 0
µ = 0.01
µ = 0.2
µ = 0
µ = 0.1
FIG. 3: a Deterministic replicator dynamics (the N → ∞ limit) of the symmetric RPS game consists of
neutrally stable orbits along which the product of the strategy frequencies xRxPxS is conserved. If global
mixing is present (µ > 0) local populations deviate from these neutral orbits toward the global average
〈x〉. Considering the simplest system with global mixing, that consisting of M = 2 local populations we
see that in the presence of global mixing population x1 and population x2 move toward each other, respec-
tively moving closer and further from the barycentre of the triangle until they become synchronized and
subsequently pursue a common orbit. For deterministic local dynamics (N → ∞) such synchronization
invariably occurs for any M if µ > 0 and typically converges to the barycentre of the simplex for suffi-
ciently homogeneous initial conditions. b The deterministic replicator dynamics of the repeated PD game is
markedly different from that of the RPS game in that the internal fixed point is unstable and in the absence
of global mixing only ALLD survives. Again turning to the simplest scenario with M = 2 we see that
if µ = 0 any pair of populations x1 and x2 (gray and black lines) converge to the to the ALLD corner.
As µ is increased above a critical value a second, stable configuration emerges: for a large subset of the
possible initial conditions (all, but the left most x2) we see that one of the populations (x1) converges to
ALLD , while the second (x1) approaches a limit cycle. If µ is increased further, the above configuration
ceases to be stable, the population which initially converges to ALLD (x1) is subsequently ”pulled out”
by global mixing, following which the two populations synchronize and are finally absorbed together in
ALLD. Simulations, however, show that synchronization may be avoided for M > 2 if µ is not too large.
coevolutionary dynamics of the three strains can be cast in terms of an RPS game, C strains kill S
strains, but are outcompeted, by R strains, because toxin production involves the suicide of bacte-
ria. The cycle is closed by S strains that outcompete R strains, because resistance requires mutant
versions of certain membrane protein, which are less efficient than the wild type [10]. Despite
the cyclic dynamics colicin-producing strains cannot coexist with sensitive or resistant strains in a
well-mixed culture, yet all three phenotypes are recovered in natural populations. Local dispersal
(modeled as explicit spatial embedding) has widely been credited with promoting the maintenance
of diversity in this system [10, 11, 13, 14].
In its most symmetric form the RPS game is described by the payoff matrix
0 −ǫ ǫ
ǫ 0 −ǫ
−ǫ ǫ 0
, (11)
and some πbase > ǫ. The dynamics of this game in an infinitely large well mixed population
consists of neutral orbits along which the product xRxPxS is conserved. For any finite N , however,
fluctuations lead to the inevitable extinction of all but one of the strategies [15]. Spatial population
structure can avert this reduction in diversity [10, 13] through the emergence of a stable fixed
point at the barycentre of the simplex . The effect of the gradual randomization of different lattice
topologies (where a small number of edges are randomly rewired) on the dynamics of the game
has also been investigated. A Hopf bifurcation leading to global oscillations was observed [34, 35]
as the fraction of rewired links was increased above some critical value.
Examining the dynamics of the symmetric RPS game in terms of our hierarchical meanfield
approximation we observed that an internal fixed point emerged for N → ∞ (Fig.3a). More
importantly, diversity was also maintained for finite local population sizes if global mixing was
present. Simulations of the time evolution of ρ(x) also revealed a Hopf bifurcation leading to
the oscillation of the global average as µ was increased above a critical value µc depending on
N (Fig.4a). These results show that previous results obtained from simulations of populations
constrained to different lattice topologies can be considered universal in the sense that not only
lattices, but any population structure that can be approximated by two distinct internally unstruc-
tured scales of mixing are sufficient for their existence. In the context of the ”paradox of the
plankton” these results imply that aside of local dispersal (modeled as explicit spatial embedding)
a minimal metapopulation structure (with local competition and global migration) can also facili-
tate the maintenance of diversity in cyclic competition systems.
0.045 0.055 0.05 0.06 0.065 0.07
0.05 0.07 0.09
b d f
µ = 0.05
1600 1400 1200 1000 800 600 400
µ = 0.05
0.16 0.12 0.08 0.04 0
0.08
0.06
0.08 0.06 0.04 0.02
N = 1000
N = 2000
N = 4000
N = 1200
small N
cyclic
TFTALLC
N = 1000
µ = 0.05µ = 0.15
N = 4000
FIG. 4: (Color online) a In the case of the rock-paper-scissors game a Hopf bifurcation similar to that
observed for populations evolving on gradually randomized lattices [34, 35] leads to the emergence of
global oscillations (the red line indicates the trajectory of 〈x〉) if µ is larger than a critical value µc(N) (see
video S1 [36]). The density ρ(x) is indicated with a blue color scale. b The ratio A of the area of the global
limit cycle and the area of the simplex is plotted as a function of µ for three different values of N . For the
repeated prisoner’s dilemma game the combination of finite local population size and global mixing µ > 0
can lead to a stationary solution (c) qualitatively similar to that observed for explicit spatial embedding e.
This state is characterized by a stable global average (large dot), just as the lattice system (data not shown)
and sustained local cycles of cooperation, defection and reciprocity, also similar to the lattice case where
groups of ALLD (red, dark grey) individuals are chased by those playing TFT (blue, black), which are
gradually outcompeted by ALLC (green, light grey). d As µ is decreased a discontinuous transition can be
observed to the ALLD phase. The ratio I of populations on the internal cycle is plotted as a function of
µ. The inset shows the transition for different values of N . f The same critical line in the µ-N plane can
be approached by increasing N with µ fixed. A large hysteresis can be observed as N is decreased below
the critical value indicating the discontinuous nature of the transition. We numerically simulated the time
evolution of ρ(x) by integrating the stochastic differential equation system defined by eq. (7) for large M
(104−105) throughout. For the RPS game we used πbase = 1 and ǫ = 0.5, while in the case of the repeated
PD game we followed ref. [38], setting T = 5, R = 3, P = 1, S = 0.1,m = 10 and c = 0.8. Lattice
simulations (e) where performed on 1000 × 1000 square lattice with an asynchronous local Moran process
between neighbors and periodic boundary conditions.
V. THE REPEATED PRISONER’S DILEMMA GAME
In the general formulation of the prisoner’s dilemma (PD) game, two players have the choice
to cooperate or to defect. Both obtain some payoff R for mutual cooperation and some lower
payoff P for mutual defection. If only one of the players defects, while the other cooperates,
the defector receives the highest payoff T and the cooperator receives the lowest payoff S. That
is T > R > P > S and defection dominates cooperation in any well-mixed population. New
strategies become possible, however if the game is repeated, and players are allowed to chose
whether to defect or cooperate based on the previous actions of the opponent. In the following
we consider, similar to refs. [37] and [38] that recently examined the role of finite population
size and mutation and finite population size, respectively in terms of the repeated PD game with
three strategies: always defect (ALLD), always cooperate (ALLC), and tit-for-tat (TFT). TFT
cooperates in the first move and then does whatever the opponent did in the previous move. TFT
has been a world champion in the repeated prisoner’s dilemma ever since Axelrod conducted his
celebrated computer tournaments [7], although it does have weaknesses and may be defeated by
other more complex strategies [39].
Previous results indicate that if only the two pure strategies are present (players who either
always defect or ones who always cooperate) explicit spatial embedding [4] and some sufficiently
sparse interaction graphs [22, 40] allow cooperation to survive and the behavior of populations is
highly sensitive to the underlying topology of the embedding [17]. We have found that introducing
global mixing into the PD game with only the two pure strategies present also allows cooperation
to survive. The mechanism responsible for favoring cooperation in this case, however, depends
on the details of the competition between local reproduction and global mixing. For more than
two strategies these details are much less relevant and do not qualitatively influence the dynamics.
We will, therefore, consider the delicate issues concerning the PD game with only the two pure
strategies in a separate publication, and concentrate here on the repeated PD game with three
strategies.
To investigate the effect of global mixing on the repeated PD game with three possible strate-
gies: ALLD, ALLC and TFT following Imhof et al. [38] we considered the payoff matrix:
large N
cyclic
small N
cyclic
TFTALLC
TFTALLC
TFTALLC
TFTALLC
N = 10
N = 10
µ = 0.05
µ = 0.05
(iii)
(iii)
deterministic local populations
(i) (iii) (i)(ii)
effective single population
N = 10
N = 10
µ = 0.15
µ = 0.15
(iii)
FIG. 5: (Color online) Phase space for the repeated prisoner’s dilemma game on a population structure
with two distinct scales (see video S2 [36]). Three different phases are possible depending on the values of
µ and N : (i) only ALLD survives (ii) an internal limit cycles is maintained by global mixing due to a large
density of local populations around the ALLD corner (iii) a globally oscillating self maintaining limit cycle
is formed. For extreme values of µ the global dynamics reduces to that of some well-mixed population
where only ALLD survives: As µ becomes negligible (µ ≪ πk for all k) we approach the limit of isolated
local populations, while for µ ≫ πk we are left with a single synchronized population. Similarly for N = 2
– the smallest system with competition – the system can be described as a single well mixed population for
any µ and ALLD again prevails. In the limit of deterministic local populations (N → ∞) all three phases
can be found depending on the value of µ. The density ρ(x) is indicated with the color scale. A figure
illustrating the phase space of the repeated prisoner’s dilemma game with fitness dependent global mixing
is included in the supplementary material [36].
Rm Sm Rm
Tm Pm T + P (m− 1)
Rm− c S + P (m− 1)− c Rm− c
, (12)
where the strategies are considered in the order ALLC, ALLD, TFT, m corresponds to the number
of rounds played and c to the complexity cost associated with conditional strategies (TFT). The
dynamics of this game has a single unstable internal fixed point and the state where each member
of the population plays ALLD is the only nontrivial stable equilibrium (Fig.3b).
Introducing global mixing, between local well-mixed populations, however, causes new sta-
tionary states to emerge . Three phases can be identified: (i) ALLD wins (ii) large fraction of local
populations in the ALLD corner maintains local cycles of cooperation defection and reciprocity
through providing an influx of defectors that prevent TFT players from being outcompeted by
ALLC playing individuals (iii) a self maintaining internal globally oscillating cycle emerges. The
simplest scenario of two (M = 2) deterministic (N → ∞) local populations coupled by global
mixing (µ > 0) already leads to the emergence of phase (ii) as demonstrated in Fig.3b while phase
(iii) only emerges for larger M . For larger M simulations show that in the limit of large local pop-
ulations all global configurations with less than some maximum ratio of the populations I on the
internal cycle are stable in phase (ii). A transition from phase (ii) to (i) happens as µ is decreased
below a critical value µii→ic and I approaches zero as I = (1−µii→ic /µ) (data not shown). This can
be understood if we considered that near the transition point a critical proportion C = µ(1− I) of
ALLD individuals needs to arrive to stabilize local cycles of cooperation defection and reciprocity.
At the critical point I = 0 and µ = µii→ic which implies C = µ
c giving I = (1− µii→ic /µ)
Exploring the N − µ phase space (Fig.5) we see that the transition from phase (i) to (ii) be-
comes discontinuous for finite N (Fig.4d,e). Further, for any given value of N and µ the global
configuration is described by a unique I due to the presence of diffusion. For appropriate values
of the parameters the global average converges to a stationary value in phase (ii) similarly to case
of explicit spatial embedding (Fig.4c).
For very small (µ ≪ πk for all k) and very large (µ ≪ πk) values of µ the global dynamics can
be reduced to that of some well-mixed population where only ALLD persists (Fig 5.). For small N
we again have an effective well-mixed population – the only limit were defectors do not dominate
is N → ∞. In comparison with previous results of Imhof et al. we can see that evolutionary
cycles of cooperation defection and reciprocity can be maintained not only by mutation, but also
by population structures with hierarchical levels of mixing.
VI. DISCUSSION
While it is, of course, clear that the reduction of any realistic population structure to a manage-
able construction is always an approximation, it has not been clearly established what the relevant
degrees of freedom are in terms of evolutionary dynamics. Meanfield approximations are a classic
method of statistical and condensed matter physics and are routinely used to circumvent intractable
combinatorial problems which arise in many-body systems. Cluster-meanfield approximations of
sufficient precision [18, 19] have been developed that adequately describe the evolutionary dynam-
ics of explicitly structured populations through systematically approximating the combinatorial
complexity of the entire topology with that of small motif of appropriate symmetry. The effects of
more minimal effective topologies have, however, not been investigated previously. In the above
we have shown that straightforward hierarchical application of the meanfield approximation (the
assumption of a well-mixed system) surprisingly unveils a new level of complexity.
In the broader context of ecological and population genetics research on structured populations
our model can be described as a metapopulation model. The term ’metapopulation’ is, however,
often used for any spatially structured population [27], and models thereof. More restrictive defi-
nitions of the term are often implied in the context of ecology and population genetics literature.
The foundations of the classic metapopulation concept where laid down by Levin’s vision of
a ”metapopulation” as a population of ephemeral local populations prone to extinction. A classic
metapopulation persists, like an ordinary population of mortal individuals, in a balance between
’deaths’ (local extinctions) and ’births’ (establishment of new populations at unoccupied sites)
[27]. This classic framework is most wide spread in the ecology literature, a less often employed
extension is the concept of a structured metapopulation where the state of the individual popula-
tions is considered in more detail, this is more similar to our concept of hierarchical mixing, but
differs in considering the possibility of local extinctions.
The effects of finite population size and migration, which our model considers, has been of
more central concern in the population genetics literature. The analog of Levin’s classic metapop-
ulation concept is often referred to as the ’finite-island’ model [28] the effective population genetic
parameters describing which, have been explored in detail[29]. The study of the population ge-
netics of spatially subdivided populations in fact predates Levin, Wright having emphasised the
capacity of drift in small populations to bring about genetic differentiation in the face of selection
and/or migration several decades prior[28].
Our hierarchical mixing model treats the coevolutionary dynamics of evolutionary games on
structured populations in a manner similar to the most simple population genetic models of spa-
tially subdivided populations, focusing on the parallel effects of selection, drift and migration. It
goes beyond these models both in considering the effects of frequency dependent selection (and
the strategic aspects of the evolutionary dynamics this implies) and in using a self-consistent ap-
proach to describe the global state of the subdivided population. Also, in order to maintain a
connection with previous work on the effects of spatial structure on evolutionary games, which
rely on Nowak’s concept of spatial games [4], with individuals restricted to interact, and hence
compete, only with neighbours as defined by some topology of interaction, we develop our model
from the level of the individual by introducing a modified version of the Moran process – and not
by extending the Wright-Fisher process (which considers discrete generations and binomial sam-
pling to account for finite population size). The effective population structure described by our
hierarchical mixing model can be thought of as a population of individuals, interactions among
which are specified by the edges of a hierarchically organized random graph. The fundamental
difference in our picture is that the edges of this graph of interactions are not considered to be
fixed, but are instead in a constant state of change, being present with a different probability be-
tween pairs of individuals who share the same local population and between pairs of individuals
who do not (Fig.1.). We consider annealed randomness, which in contrast to the usual quenched
picture of fixed edges is insensitive to the details of topology. Our approach we believe best facil-
itates the exploration of the effects of changing the relative strengths of drift and migration in the
context of evolutionary games on structured populations.
Examining the effects of hierarchical mixing in the context of the evolution of robustness we
demonstrated that biased influx coupled with drift can result in cooperation being favored, pro-
vided the ratio of benefit to cost exceeds the local population size. This result bears striking
resemblance to that of Ohtsuki et al. [22], who were able to calculate the fixation probability of a
randomly placed mutant for any two-person, two-strategy game on a regular graph and found that
cooperation is favored provided the ratio of benefit to cost exceeds the degree of the graph. Our
results demonstrate that this rule extends to the minimal spatial structure induced by hierarchical
levels of mixing.
Applying our model of spatial structure to the repeated prisoners dilemma revealed that a con-
stant influx of defectors can help to stabilize cycles of cooperation, defection, and reciprocity
through preventing the emergence of an intermittent period of ALLC domination in the popula-
tion, which would present a situation that ”leaves the door wide open” to domination by defectors.
While previous work has been done on the effects of ”forcing” cooperation [41] the idea that an
influx of defectors can in fact stabilize the role of reciprocity in promoting cooperation has not
been proposed previously. It seems highly unlikely that this mechanism can be explained in terms
of kin or multilevel (group) selection, the similarities between which in structured populations
have recently been the subject of intensive debate (see e.g. [42] and [43] or [44] and [45]). Kin
selection can operate whenever interactions occurring among individuals who share a more recent
common ancestor than individuals sampled randomly from the whole population [45] are relevant.
In our case it is the interaction between defectors, arriving from the global scale, and TFT players
present at the local scale that is important, and not the interaction between individuals in the local
population, who may be thought of as sharing a recent common ancestor due to local dispersal.
Also, while the concept of multilevel selection presents a promising framework for the study of
evolution of cooperation, it must nonetheless be possible to derive it from ”first principles” – just
as kin selection can be cast as an emergent effect of local dispersal.
While there has been considerable work on studying the evolutionary games on graphs and
highly symmetric spatial structures very little attention has been paid to the effects of more min-
imal effective population structures, despite their widespread application in ecology and popula-
tion genetics, fields from which evolutionary game theory was born and must ultimately reconnect
with. We believe that the minimal population structure that such a hierarchical meanfield theory
describes is potentially more relevant in a wide range of natural systems, than more subtle setups
with a delicate dependence on the details and symmetries of the topology. We showed through
two examples that such structure is sufficient for the emergence of some phenomena previously
only observed for explicit spatial embedding, demonstrating the potential of our model to identify
robust effects of population structure on the dynamics of evolutionary games that do not depend
on the details of the underlying topology. The practical advantage of our approach, lies in its
ability to readily determine whether or not some feature of a structured population depends on the
topological details of local interactions.
Recent simulation result concerning the dynamics of public goods games on different popula-
tion structures [9, 46] and experiments where global mixing in an RPS like bacteria-phage system
lead to the emergence of a ”Tragedy of the commons” scenario [47] should all be amicable to
analysis in terms of our method.
VII. ACKNOWLEDGMENTS
This work was partially supported by the Hungarian Scientific Research Fund under grant No:
OTKA 60665.
VIII. APPENDIX
Our approach readily generalizes for an arbitrary number of hierarchical mixing levels. For
three levels of mixing we may consider the global population to be comprised of M subpopula-
tions each of which is in turn subdivided into M local populations. With m ∈ {1, · · · ,M} running
over subpopulations and l ∈ {1, · · · ,M} over local populations the transition probabilities can be
written as:
T̂mlik =
πmlk n
k + µ
(1)〈(1)σml′k 〉l′ + µ(2)〈〈(2)σm
k 〉l′〉m′
k=1(π
+ µ(1)〈(1)σml′
〉l′ + µ(2)〈〈(2)σm
〉l′〉m′)
, (13)
where primed indices indicate the scale of mixing over which the average is taken, µ(1) describes
the strength of mixing, and the (1)σmlk the tendencies of mixing among local populations within a
subpopulation, while µ(2) describes the strength of mixing, and the (2)σmlk the tendencies of mixing
among subpopulations in the global population.
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Introduction
Hierarchical Meanfield Theory for Two Distinct Scales
Cooperation in populations with hierarchical levels of mixing
The RPS Game
The Repeated Prisoner's Dilemma Game
Discussion
Acknowledgments
Appendix
References
|
0704.0358 | Flavor Physics in SUSY at large tan(beta) | Flavor Physics in SUSY at large tan β
Paride Paradisi
Departament de F́ısica Teòrica and IFIC, Universitat de València-CSIC, E-46100, Burjassot, Spain.
We discuss the phenomenological impact of a particularly interesting corner of the MSSM: the
large tanβ regime. The capabilities of leptonic and hadronic Flavor Violating processes in shedding
light on physics beyond the Standard Model are reviewed. Moreover, we show that tests of Lepton
Universality in charged current processes can represent an interesting handle to obtain relevant
information on New Physics scenarios.
I. INTRODUCTION
Despite the great phenomenological success of the
Standard Model (SM), it is natural to consider this the-
ory only as the low-energy limit of a more general model.
The direct exploration of New Physics (NP) particles
at the TeV scale will be performed at the upcoming
LHC. A complementary strategy in looking for NP is
provided by high-precision low-energy experiments where
NP could be detected through the virtual effects of NP
particles. In particular, flavor-changing neutral-current
(FCNC) transitions may exhibit a sensitivity reach even
beyond that achievable by the direct searches at the LHC
while representing, at the same time, the best (or even
the only) tool to extract information about the flavor
structures of NP theories.
In view of the above considerations, it is clear that
flavor physics provides necessary and complementary in-
formation to those obtainable by the LHC.
Besides FCNC decays, also the Lepton Flavor Univer-
sality (LFU) tests (Kℓ2 and πℓ2) offer a unique opportu-
nity to probe the SM and thus, to shed light on NP: the
smallness of NP effects is more than compensated by the
excellent experimental resolution and the good theoreti-
cal control.
II. LFV IN SUSY
The discovery of neutrino masses and oscillations has
unambiguously pointed out the existence of the Lepton
Flavor Violation (LFV) thus, we expect this phenomenon
to occur also in the charged-lepton sector.
Within a SM framework with massive neutrinos,
FCNC transitions in the lepton sector like ℓi → ℓjγ are
strongly suppressed by the GIM mechanism at the level
of B(ℓi → ℓjγ) ∼ (mν/mW )
4 ∼ 10−50 well beyond any
realistic experimental resolution [1]. In this sense, the
search for FCNC transitions of charged leptons is one of
the most promising directions where to look for physics
beyond the SM.
Within a SUSY framework, LFV effects originate from
any misalignment between fermion and sfermion mass
eigenstates. In particular, if the light neutrino masses
are obtained via a see-saw mechanism, the radiatively
induced LFV entries in the slepton mass matrix (m2
are given by [2]:
)i6=j ≈ −
ν )i6=j ln
, (1)
where MX denote the scale of SUSY-breaking media-
tion andm0 the universal supersymmetry breaking scalar
mass. Since the see–saw equation 1 allows large (YνY
entries, sizable effects can stem from this running [2].
The determination of (m2
)i6=j would imply a complete
knowledge of the neutrino Yukawa matrix (Yν)ij , which
is not possible even if all the low-energy observables from
the neutrino sector were known. As a result, the predic-
tions of leptonic FCNC effects will remain undetermined
even in the very optimistic situation where all the rele-
vant NP masses were measured at the LHC.
This is in contrast with the quark sector, where similar
RGE contributions are completely determined in terms
of quark masses and CKM-matrix elements.
More stable predictions can be obtained embedding
the SUSY model within a Grand Unified Theory (GUT)
where the see-saw mechanism can naturally arise (such
as SO(10)). In this case the GUT symmetry allows us to
obtain some hints about the unknown neutrino Yukawa
matrix Yν . Moreover, in GUT scenarios there are other
contributions stemming from the quark sector [3]. These
effects are completely independent from the structure of
Yν and can be regarded as new irreducible LFV contribu-
tions within SUSY GUTs. For instance, within SU(5),
as both Q and ec are hosted in the 10 representation, the
CKM matrix mixing the left handed quarks will give rise
to off diagonal entries in the running of the right-handed
slepton soft masses [3].
There exist to different classes of LFV contributions to
rare decays:
i) Gauge-mediated LFV effects through the exchange
of gauginos and sleptons,
ii) Higgs-mediated LFV effects through effective non-
holomorphic Yukawa interactions [4] .
1 The effective light-neutrino mass matrix obtained from a see-
saw mechanism is mν = −YνM̂
ν 〈Hu〉
2, where M̂R is the
3 × 3 right-handed neutrino mass matrix and Yν are the 3 ×
3 Yukawa couplings between left- and right-handed neutrinos
(the potentially large sources of LFV), and 〈Hu〉 is the vacuum
expectation value of the up-type Higgs.
http://arxiv.org/abs/0704.0358v1
The above contributions decouple with the heaviest mass
in the slepton/gaugino loops mSUSY (case i)) or with the
heavy Higgs mass mH (case ii)).
In principle, mH and mSUSY refers to different mass
scales. Higgs mediated effects start being competitive
with the gaugino mediated ones when mSUSY is roughly
one order of magnitude heavier then mH and for tanβ ∼
O(50) [5].
While the appearance of LFV transitions would un-
ambiguously signal the presence of NP, the underlying
theory generating LFV phenomena will remain undeter-
mined, in general.
A powerful tool to disentangle among NP theories is
the study of the correlations of LFV transitions among
same families [5, 6, 7].
Interestingly enough, the predictions for the correla-
tions among LFV processes are very different in the
gauge- and Higgs-mediated cases [5]. In this way, if sev-
eral LFV transitions are observed, their correlated anal-
ysis could shed light on the underlying mechanism of
LFV. In the case of gauge-mediated LFV amplitudes the
ℓi → ℓjℓkℓk decays are dominated by the ℓi → ℓjγ
∗ dipole
transition, which leads to the unambiguous prediction:
B(ℓi → ℓjℓkℓk)
B(ℓi → ℓjγ)
B(µ− e in Ti)
B(µ→eγ)
≃αel . (3)
If some ratios different from the above were discovered,
then this would be clear evidence that some new process
is generating the ℓi → ℓj transition, with Higgs mediation
being a potential candidate 2.
As regards the Higgs mediated case, Br(τ → ljγ) still
gets generally the largest contribution among all the pos-
sible LFV decay modes [5]. The following approximate
relations hold [5]:
Br(τ → ljγ)
Br(τ → ljη)
& 1 ,
Br(τ → ljη)
Br(τ → ljµµ)
3+5δjµ
. (4)
Br(τ → ljee)
Br(τ → ljµµ)
3+5δjµ
. (5)
Br(µ → eγ)
Br(µAl → eAl)
∼ 10 ,
Br(µ → eee)
Br(µ → eγ)
∼ αel . (6)
On the other hand, a correlated study of processes of the
same type but relative to different family transitions, like
2 As recently shown in [7], a powerful tool to disentangle between
Little Higgs models with T parity (LHT) and SUSY theories
is a correlated analysis of LFV processes. In fact, LHT and
SUSY theories predict very different correlations among LFV
transitions [7].
Br(µ → eγ)/Br(τ → µγ) ∼ [(m2
)21/(m
2, provides
important information about the unknown structure of
the LFV source, i.e. (m2
)i6=j .
III. LFU IN SUSY
High precision electroweak tests, such as deviations
from the SM expectations of the LFU breaking, represent
a powerful tool to probe the SM and, hence, to constrain
or obtain indirect hints of new physics beyond it. Kaon
and pion physics are obvious grounds where to perform
such tests, for instance in the π → ℓνℓ and K → ℓνℓ
decays, where l = e or µ. In particular, the ratios
B(P → µν)
B(P → eν)
can be predicted with excellent accuracies in the SM,
both for P = π (0.02% accuracy [8]) and P = K (0.04%
accuracy [8]), allowing for some of the most significant
tests of LFU.
As recently pointed out in Ref. [9], large departures
from the SM expectations can be generated within a
SUSY framework with R-parity only once we assume i)
LFV effects, ii) large tanβ values.
Denoting by ∆r
NP the deviation from µ−e universal-
ity in RK due to NP, i.e.: R
K = (R
K )SM
1 + ∆r
it turns out that [9]:
|∆31R |
2 tan6β. (8)
The deviations from the SM could reach ∼ 1% in the
K case [9] (not far from the present experimental res-
olution [10]) and ∼ few × 10−4 in the R
π case while
maintaining LFV effects in τ decays at the 10−10 level. In
the pion case the effect is quite below the present experi-
mental resolution [11], but could well be within the reach
of the new generation of high-precision πℓ2 experiments
planned at TRIUMPH and at PSI. Larger violations of
LFU are expected in B → ℓν decays, with O(10%) devi-
ations from the SM in R
B and even order-of-magnitude
enhancements in R
B [12].
IV. FLAVOR PHYSICS AT LARGE tanβ AND
DARK MATTER
Within the MSSM, the scenario with large tanβ and
heavy squarks is particularly interesting. On the one
hand, values of tanβ ∼ 30–50 can allow the unification
of top and bottom Yukawa couplings, as predicted in
well-motivated grand-unified models [13]. On the other
hand, a Minimal Flavor Violating (MFV) structure [14]
with heavy (∼ TeV ) soft-breaking terms in the quark
sector and large tanβ ∼ 30 − 50 values leads to in-
teresting phenomenological virtues [12, 15]: the present
(g − 2)µ anomaly and the upper bound on the Higgs
boson mass can be easily accommodated, while satisfy-
ing all the present tight constraints in the electroweak
and flavor sectors. Additional low-energy signatures of
this scenario could possibly show up in the near future
in B(Bu → τν), B(Bs,d → ℓ
+ℓ−) and B(B → Xsγ). In
the following, as discussed in [16], we analyze the above
scenario under the additional assumption that the relic
density of a Bino-like lightest SUSY particle (LSP) ac-
commodates the observed dark matter distribution
0.094 ≤ ΩCDMh
2 ≤ 0.129 at 2σ C.L. . (9)
In the regime with large tanβ and heavy squarks, the
relic-density constraints can be easily satisfied mainly in
the so called A-funnel region [17] where MB̃ ≈ MA/2.
The combined constraints from low-energy observables
and dark matter in the tanβ–MH plane are illustrated
in Figure 1 (left). The light-blue areas are excluded since
the stau turns out to be the LSP, while the yellow band
denotes the allowed region where the stau coannihilation
mechanism is also active. The remaining bands corre-
spond to the following constraints/reference-ranges from
low-energy observables:
• B → Xsγ [1.01 < RBsγ < 1.24]: allowed region
between the two blue lines.
• aµ [2 < 10
−9(aexpµ − a
µ ) < 4 [18]]: allowed region
between the two purple lines.
• B → µ+µ− [Bexp < 8.0×10−8 [19]]: allowed region
below the dark-green line.
• ∆MBs [∆MBs = 17.35 ± 0.25 ps
−1 [20]]: allowed
region below the gray line.
• B → τν [0.8 < RBτν < 0.9]: allowed region be-
tween the two black lines [ red (green) area if all
the other conditions (but for aµ) are satisfied].
From Figure 1 (right), we deduce that there is a quite
strong correlation between ∆aµ and B(Bu → τν) thanks
to the A-funnel region condition MH ≈ 2M1. A SUSY
contribution to aµ of O(10
−9) generally implies a sizable
effect in 0.7 < B(Bu → τν) < 0.9. A more precise de-
termination of B(Bu → τν) is therefore a key element to
test this scenario.
The interplay of B physics observables, dark-matter
constraints, ∆aµ of O(10
−9), and LFV rates is shown
in Figure 2. For a natural choice of |δ12LL| = 10
B(µ → eγ) is in the 10−12 range, i.e. well within the
reach of MEG [21] experiment. On the other hand,
B(τ → µγ) lies within the 10−9 range for a |δ23LL| = 10
that is a natural size expected in many models.
Acknowledgments
I wish to thank the conveners of WG3 for the kind
invitation and G. Isidori, F. Mescia and D. Temes for
collaborations on which this talk is partly based. I also
acknowledge support from the EU contract No. MRTN-
CT-2006-035482, ”FLAVIANET” and from the Spanish
MEC and FEDER FPA2005-01678.
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http://arxiv.org/abs/hep-ph/9408406
http://arxiv.org/abs/hep-ph/9501334
http://arxiv.org/abs/hep-ph/0605139
http://arxiv.org/abs/hep-ph/0206310
http://arxiv.org/abs/hep-ph/0508054
http://arxiv.org/abs/hep-ph/0601100
http://arxiv.org/abs/hep-ph/0401100
http://arxiv.org/abs/hep-ph/0702136
http://arxiv.org/abs/hep-ph/0511289
http://arxiv.org/abs/hep-ph/0605012
http://arxiv.org/abs/hep-ph/9308333
http://arxiv.org/abs/hep-ph/0605177
http://arxiv.org/abs/hep-ph/0703035
http://arxiv.org/abs/hep-ph/0611102
http://arxiv.org/abs/hep-ph/0509372
http://arxiv.org/abs/hep-ex/0508058
http://arxiv.org/abs/hep-ex/0606027
FIG. 1: Left plot: Combined constraints from low-energy observables and dark matter in the tan β–MH plane setting [µ,Mℓ̃] =
[0.5, 0.4] TeV. The light-blue area is excluded by the dark-matter conditions [16]. Within the red (green) area all the reference
values of the low-energy observables (but for aµ) are satisfied. The yellow band denote the area where the stau coannihilation
mechanism is active (1 < Mτ̃R/MB̃ < 1.1); in this area the A-funnel region (where MH ≈ 2M1) and the stau coannihilation
region overlap. Right plot: ∆aµ = (gµ−g
µ )/2 vs. the slepton mass within the funnel region taking into account the B → Xsγ
constraint and setting RBτν > 0.7 (blue), RBτν > 0.8 (red), RBτν > 0.9 (green) [16]. The supersymmetric parameters have
been varied in the following ranges: 200 GeV ≤ M2 ≤ 1000 GeV, 500 GeV ≤ µ ≤ 1000 GeV, 10 ≤ tan β ≤ 50. In both plots,
we have set AU = −1 TeV, Mq̃ = 1.5 TeV, and imposed the GUT relation M1 ≈ M2/2 ≈ M3/6.
FIG. 2: Isolevel curves for B(µ → eγ) and B(τ → µγ) assuming |δ12LL| = 10
−4 and |δ23LL| = 10
−2 in the tan β–MH plane [16].
The green/red areas correspond to the allowed regions for the low-energy observables illustrated in Figure 1 for [µ,M
[0.5, 0.4] TeV.
|
0704.0359 | Some properties of the complex Monge-Ampere operator in Cegrell's
classes and applications | arXiv:0704.0359v1 [math.CV] 3 Apr 2007
Some properties of the complex Monge-Ampère operator in
Cegrell’s classes and applications
NGUYEN VAN KHUE and PHAM HOANG HIEP
Abstract. In this article we will first prove a result about convergence in capacity. Using
the achieved result we will obtain a general decompositon theorem for complex Monge-
Ampère measues which will be used to prove a comparison principle for the complex
Monge-Ampère operator.
2000 Mathematics Subject Classification: Primary 32W20, Secondary 32U15.
Key words and phrases: complex Monge-Ampère operator, plurisubharmonic function.
This work was supported by the National Research Program for Natural Sciences, Vietnam.
1. Introduction
Let Ω be a bounded hyperconvex domain in Cn. By PSH(Ω) we denote the set of plurisub-
harmonic (psh) functions on Ω. In [BT 1,2] the authors established and used the compari-
son principle to study the Dirichlet problem in PSH∩L∞loc(Ω). Recently, Cegrell introduced
a general class E of psh functions on which the complex Monge-Ampère operator (ddc.)n
can be defined. He obtained many important results of pluripotential theory in the class E .
For example, the ones on the comparison principle and solvability of the Dirichlet problem
(see [Ce 1-3]).
The main result of our paper are Theorem 4.1 and some Xing type comparision principles.
Theorem 4.1 is generalize Lemma 5.4 in [Ce1], Lemma 7.2 in [Åh] and Lemma 3.4 in [Ce3].
For definitions of Cegrell’s classes see Section 2. After giving some preliminaries, we start in
Proposition 3.1 with a comparison principle, which is analogous to a comparison principle
due to Xing (Lemma 1 in [Xi1]). It should be observed that our proof is quite different
from Xing’s proof, and the inequality we obtain is slightly stronger than Xing’s inequality,
even in the case of bounded psh functions. Using Proposition 3.1, we give in Theorem
3.5 a sufficient condition for Cn-capacity convergence of a sequence of psh functions in
the class F . This result should be compared to Theorem 3 of [Xi1] where the situation
of bounded psh functions was studied. Applying Theorem 3.5 we give generalizations of
recent results in [Cz] and [CLP] about convergences of multipole Green functions and a
criterion for pluripolarity, respectively. Section 4 focuses on Theorem 4.1 and Theorem
4.9. By applying Theorem 4.1 we give some results on class Cegrell’s classes. We prove in
Proposition 4.4 a local estimate for the Monge-Ampère measure in terms of the Beford-
Taylor relative capacity. As an application, we give in Theorem 4.5 a decomposition result
for Monge-Ampère measure, which is similar in spirit to Theorem 6.3 in [Ce1]. From
Proposition 3.1 and Theorem 4.1 we obtain easily a Xing type comparison principle for
functions in classes F and E .
Acknowledgment. We are grateful to Professor Urban Cegrell for useful discussions that
helped to improve the paper. We are grateful to Per Åhag for fruitful comments. This
work is supported by the National Research Program for Natural Sciences, Vietnam.
http://arxiv.org/abs/0704.0359v1
2. Preliminaries
First we recall some elements of pluripotential theory that will be used throughout the
paper. All this can be found in [BT2], [Ce1], [Ce2], [Le].
2.1. We will always denote by Ω a bounded hyperconvex domain in Cn unless other wise
stated. The Cn-capacity in the sense of Bedford and Taylor on Ω is the set function given
Cn(E) = Cn(E,Ω) = sup{
(ddcu)n : u ∈ PSH(Ω), −1 ≤ u ≤ 0}
for every Borel set E in Ω. It is proved in [BT2] that
Cn(E) =
(ddch∗E,Ω)
where h∗E,Ω is the upper regularization of the relative extremal function hE,Ω for E (relative
to Ω) i.e.,
hE,Ω(z) = sup{u(z) : u ∈ PSH
−(Ω), u ≤ −1 on E}.
The following concepts are taken from [Xi1] and [Xi2]
∗A sequence of functions uj on Ω is said to converge to a function u in Cn-capacity on a
set E ⊂ Ω if for every δ > 0 we have Cn({z ∈ E : |uj(z) − u(z)| > δ}) → 0 as j → ∞.
∗A family of positive measures {µα} on Ω is called uniformly absolutely continuous with
respect to Cn-capacity in a set E ⊂ Ω if for every ǫ > 0 there exists δ > 0 such that for
each Borel subset F ⊂ E with Cn(F)< δ the inequality µα(F)< ǫ holds for all α. We write
µα ≪ Cn in E uniformly for α.
2.2. The following classes of psh functions were introduced by Cegrell in [Ce1] and [Ce2]
E0 = E0(Ω) = {ϕ ∈ PSH
−(Ω) ∩ L∞(Ω) : lim
ϕ(z) = 0,
(ddcϕ)n < +∞},
F = F(Ω) = {ϕ ∈ PSH−(Ω) : ∃ E0(Ω) ∋ ϕj ց ϕ, sup
(ddcϕj)
n < +∞},
E = E(Ω) = {ϕ ∈ PSH−(Ω) : ∃ ϕK ∈ F(Ω) such that ϕK = ϕ on K, ∀K ⊂⊂ Ω},
Ea = Ea(Ω) = {u ∈ E(Ω) : (ddcu)n(E) = 0 ∀ E is pluripolar in Ω}.
For each u ∈ F(Ω), we set
e0(u) =
(ddcu)n.
2.3. Let A = {(wj , νj)}j=1,...,p be a finite subset of Ω × R
+. According to Lelong (see
[Le]), the pluricomplex Green function with poles in A is defined by
g(A)(z) = sup{u(z) : u ∈ LA}
where
LA = {u ∈ PSH
−(Ω) : u(z) − νj log |z − wj | ≤ O(1) as z → wj , j = 1, ..., p}
ν(A) =
νnj , Â = {wj}j=1,...,p.
2.4. We write lim
[u(z) − v(z)] ≥ a if for every ǫ > 0 there exists a compact set K in Ω
such that
u(z) − v(z) ≥ a− ǫ for z ∈ (Ω\K) ∩ {u > −∞}
v(z) = −∞ for z ∈ (Ω\K) ∩ {u = −∞}.
2.5. Xing’s comparison principle (see Lemma 1 in [Xi1]). Let Ω be a bounded open subset
in Cn and u, v ∈ PSH∩L∞(Ω) satisfy lim
[u(z)− v(z)] ≥ 0. Then for any constant r ≥ 1
and all wj ∈ PSH(Ω) with 0 ≤ wj ≤ 1, j = 1, 2, ..., n we have
(n!)2
{u<v}
(v − u)nddcw1 ∧ ... ∧ dd
{u<v}
(r − w1)(dd
v)n ≤
{u<v}
(r − w1)(dd
3. Some convergence theorems
In order to study the convergence of a sequence of psh functions in Cn-capacity, we start
with the following.
3.1. Proposition. a) Let u, v ∈ Fsuch that u ≤ v on Ω. Then for 1 ≤ k ≤ n
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
for all wj ∈ PSH(Ω), 0 ≤ wj ≤ 1, j = 1, ..., k, wk+1, ..., wn ∈ F and all r ≥ 1.
b) Let u, v ∈ E such that u ≤ v on Ω and u = v on Ω\K for some K ⊂⊂ Ω. Then for
1 ≤ k ≤ n
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
for all wj ∈ PSH(Ω), 0 ≤ wj ≤ 1, j = 1, ..., k, wk+1, ..., wn ∈ E and all r ≥ 1.
We proceed through some lemmas.
3.2. Lemma. Let u, v ∈ PSH ∩ L∞(Ω) such that u ≤ v on Ω and lim
[u(z) − v(z)] = 0.
(v − u)kddcw ∧ T ≤ k
(1 − w)(v − u)k−1ddcu ∧ T
for all w ∈ PSH(Ω), 0 ≤ w ≤ 1 and all positive closed currents T .
Proof. First, assume u, v ∈ PSH∩L∞(Ω), u ≤ v on Ω and u = v on Ω\K, K ⊂⊂ Ω. Then,
using the Stokes formula we obtain
(v − u)kddcw ∧ T =
(v − u)kddc(w − 1) ∧ T
(w − 1)ddc(v − u)k ∧ T
= −k(k − 1)
(1 − w)d(v − u) ∧ dc(v − u) ∧ T
(1 − w)(v − u)k−1ddc(u− v) ∧ T
(1 − w)(v − u)k−1ddc(u− v) ∧ T
(1 − w)(v − u)k−1ddcu ∧ T.
General case, for each ǫ > 0 we set vǫ = max(u, v − ǫ). Then vǫ ր v on Ω, vǫ ≥ u on Ω
and vǫ = u on Ω\K for some K ⊂⊂ Ω. Hence
(vǫ − u)
kddcw ∧ T ≤ k
(1 − w)(vǫ − u)
k−1ddcu ∧ T.
Since 0 ≤ vǫ − u ր v − u as ǫ ց 0, letting ǫ ց 0 we get
(v − u)kddcw ∧ T ≤ k
(1 − w)(v − u)k−1ddcu ∧ T.
3.3. Lemma. Let u, v ∈ PSH ∩ L∞(Ω) such that u ≤ v on Ω and lim
[u(z) − v(z)] = 0.
Then for 1 ≤ k ≤ n
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cv)k ∧ T
(r − w1)(dd
cu)k ∧ T.
for all w1, ..., wk ∈ PSH(Ω), 0 ≤ wj ≤ 1 ∀ j = 1, ..., k, wk+1, ..., wn ∈ E and all r ≥ 1.
Proof. To simplify the notation we set
T = ddcwk+1 ∧ ... ∧ dd
First, assume that u, v ∈ PSH ∩ L∞(Ω), u ≤ v on Ω, and u = v on Ω\K, K ⊂⊂ Ω. Using
Lemma 3.2 we get
(v − u)kddcw1 ∧ ... ∧ dd
cwn ≤ k
(v − u)k−1ddcw1 ∧ ... ∧ dd
cwk−1 ∧ dd
cu ∧ T
≤ ...
(v − u)ddcw1 ∧ (dd
cu)k−1 ∧ T
(v − u)ddcw1 ∧ [
(ddcu)i ∧ (ddcv)k−i−1] ∧ T
(w1 − r)dd
c(v − u) ∧ [
(ddcu)i ∧ (ddcv)k−i−1] ∧ T
(r − w1)dd
c(u− v) ∧ [
(ddcu)i ∧ (ddcv)k−i−1] ∧ T
(r − w1)[(dd
cu)k − (ddcv)k] ∧ T.
General case, for each ǫ > 0 we put vǫ = max(u, v − ǫ). Then vǫ ր v on Ω, vǫ ≥ u on Ω
and vǫ = u on Ω\K for some K ⊂⊂ Ω. Hence
(vǫ − u)
kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
k ∧ T
(r − w1)(dd
u)k ∧ T.
Observe that 0 ≤ vǫ − u ր v− u and (dd
k ∧ T → (ddcv)k ∧ T weakly as ǫ ց 0, r−w1
is lower semicontinuous, by letting ǫ ց 0 we have
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cv)k ∧ T
(r − w1)(dd
cu)k ∧ T.
The proof is finished.
Proof of Proposition 3.1. a) Let E0 ∋ uj ց u and E0 ∋ vj ց v as in the definition of F .
Replace vj by max(uj , vj) we may assume that uj ≤ vj for j ≥ 1. By Lemma 3.3 we have
(vj − ut)
kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
k ∧ ddcwk+1 ∧ ... ∧ dd
for t ≥ j ≥ 1. By Proposition 5.1 in [Ce2] letting t → ∞ in the above inequality we have
(vj − u)
w1 ∧ ... ∧ dd
(r − w1)(dd
k ∧ T
(r − w1)(dd
cu)k ∧ T
for j ≥ 1. Next letting j → ∞ again by Proposition 5.1 in [Ce2] we get the desired
conclusion.
b) Let G,W be open sets such that K ⊂⊂ G ⊂⊂ W ⊂⊂ Ω. According to the remark
following Definition 4.6 in [Ce2] we can choose a function ṽ ∈ F such that ṽ ≥ v and ṽ = v
on W . Set
u on G
ṽ on Ω\G
Since u = v = ṽ on W\K we have ũ ∈ PSH−(Ω). It is easy to see that ũ ∈ F , ũ ≤ ṽ and
ũ = u on W . By a) we have
(ṽ − ũ)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cṽ)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
ũ)k ∧ ddcwk+1 ∧ ... ∧ dd
Since ũ = ṽ on Ω\G we have
(ṽ − ũ)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cṽ)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
cũ)k ∧ ddcwk+1 ∧ ... ∧ dd
Since ũ = u, ṽ = v on W and u = v on Ω\K we obtain
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
3.4. Proposition. Let u, v ∈ F and u ≤ v on Ω. Then
(v − u)nddcw1 ∧ ... ∧ dd
(−w1)[(dd
u)n − (ddcv)n]
for all wj ∈ PSH(Ω), −1 ≤ wj ≤ 0, j = 1, ..., n.
Proof. The proposition follows from Proposition 3.1 with k = n, r = 1 and wj are replaced
by wj + 1.
3.5. Theorem. Let u, uj ∈ F and uj ≤ u for j ≥ 1. Assume that sup
(ddcuj)
n < +∞
and ||(ddcuj)
n − (ddcu)n||E → 0 as j → ∞ for all E ⊂⊂ Ω. Then uj → u in Cn-capacity
on every E ⊂⊂ Ω as j → ∞.
Proof. Let Ω′ ⊂⊂ Ω and δ > 0. Put
Aj = {z ∈ Ω′ : |uj − u| ≥ δ} = {z ∈ Ω′ : u− uj ≥ δ}.
We prove that Cn(Aj) → 0 as j → ∞. Given ǫ > 0. By quasicontinuity of u and uj , there
is an open set G in Ω such that Cn(G) < ǫ, and uj |Ω\G, u|Ω\G are continuous. We have
Aj = Bj ∪ {z ∈ G : u− uj ≥ δ}.
where Bj = {z ∈ Ω′\G : u− uj ≥ δ} are compact sets in Ω and
Cn(Aj) ≤ lim
Cn(Bj) + ǫ
We claim that lim
Cn(Bj) = 0. By Proposition 3.4 we have
Cn(Bj) =
(ddch∗Bj )
(u− uj)
n(ddch∗Bj )
(−h∗Bj )[(dd
n − (ddcu)n]
{||(ddcuj)
n − (ddcu)n||K +
(−hΩ′)[(dd
n + (ddcu)n]}
{||(ddcuj)
n − (ddcu)n||K + sup
|hΩ′ |[sup
(ddcuj)
(ddcu)n]}.
As lim
hΩ′(z) = 0 there exists K ⊂⊂ Ω such that
|hΩ′ |[sup
(ddcuj)
(ddcu)n] < ǫ.
By the hypothesis
||(ddcuj)
n − (ddcu)n||K < ǫ for j > j0.
Cn(Bj) < 2ǫ for j > j0.
This proves the claim and hence the theorem.
As an application of Theorem 3.5 we have the following
3.6. Proposition. Let g(Aj) be multipolar Green functions on Ω such that
Âj = {w
1, ..., w
} → ∂Ω and sup
ν(Aj) = sup
)n < +∞
Then g(Aj) → 0 as j → ∞ in Cn-capacity.
Proof. By the hypothesis we have
(ddcg(Aj))
n(Ω) = sup
ν(Aj) < +∞
||(ddcg(Aj))
n||K → 0 as j → ∞ for all K ⊂⊂ Ω.
Theorem 3.5 implies that g(Aj) → 0 as j → ∞ in Cn-capacity.
This section ends up with a criterion for pluripolarity
3.7. Theorem. Let uj ∈ F such that sup
(ddcuj)
n < +∞.
Then there is a constant A > 0 such that
i)( lim
∗ ∈ F .
ii)Cn({z ∈ Ω : ( lim
∗(z) < −t}) ≤ A
iii){z ∈ Ω : lim
uj(z) = −∞} is pluripolar.
Proof. i) For each j ≥ 1 put vj = sup{uj, uj+1, ...}. By [Ce2] v
j ∈ F and
(ddcv∗j )
n ≤ sup
(ddcuj)
n < +∞.
By [Ce2] we have v∗j ց v ∈ F .
ii) By Proposition 3.1 in [CKZ] we have
Cn{z ∈ Ω : ( lim
∗(z) < −t} = Cn{z ∈ Ω : v(z) < −t} ≤
2ne0(v)
where A = 2ne0(v).
iii) According to [BT2] we have
Cn{z ∈ Ω : lim
uj(z) = −∞} = Cn{z ∈ Ω : v(z) = −∞} = 0.
Remark. Theorem 3.7 in the case where uj are multipole Green functions was proved by
D.Coman, N.Levenberg and A.Poletsky in Theorem 4.1 of [CLP].
4. Some properties of the Cegrell’s classes and applications
In this section, first we prove the following
4.1. Theorem. Let u, u1, ..., un−1 ∈ E , v ∈ PSH
−(Ω) and T = ddcu1 ∧ ... ∧ dd
cun−1.
ddc max(u, v) ∧ T |{u>v} = dd
cu ∧ T |{u>v}.
We need the following well-known fact.
4.2. Lemma. Let µ be a measure on Ω and f : Ω → R a measurable function on Ω. The
following are equivalent
i)µ(E) = 0 for all Borell sets E ⊂ {f 6= 0}.
fdµ = 0 for every measurable set E in Ω.
Proof. i)⇒ii) follows from:
fdµ =
E\{f=0}
fdµ +
E∩{f=0}
fdµ = 0
ii)⇒i). It suffices to show that µ = 0 on every Xδ = {f > δ > 0}. By the Hahn
decomposition theorem, there exist measurable subsets X+
and X−
of Xδ such that Xδ =
= ∅ and µ ≥ 0 on X+
, µ ≤ 0 on X−
. We have
δµ(X+
fdµ = 0
δµ(X−
fdµ = 0
Hence, µ(X+
) = µ(X−
) = 0. Therefore, we have µ = 0 on Xδ.
Proof of Theorem 4.1.
a) First we prove the proposition for v ≡ a < 0. According to the remark following
Definition 4.6 in [Ce2], without loss of generality we may assume that u, u1, ..., un−1 ∈ F .
Using Theorem 2.1 in [Ce2] we can find
E0 ∩ C(Ω̄) ∋ u
j ց u, E0 ∩ C(Ω̄) ∋ u
ց uk, k = 1, ..., n− 1.
Since {uj > a} is open we have
ddc max(uj, a) ∧ Tj |{uj>a} = dd
cuj ∧ Tj |{uj>a}.
Thus from the inclusion {u > a} ⊂ {uj > a} we obtain
ddc max(uj , a) ∧ Tj |{u>a} = dd
cuj ∧ Tj |{u>a}.
where Tj = dd
1 ∧ ... ∧ dd
n−1. By Corollary 5.2 in [Ce2], it follows that
max(u− a, 0)ddc max(uj , a) ∧ Tj → max(u− a, 0)dd
c max(u, a) ∧ T.
max(u− a, 0)ddcuj ∧ Tj → max(u− a, 0)dd
cu ∧ T.
Hence
max(u− a, 0)[ddc max(u, a) ∧ T − ddcu ∧ T ] = 0.
Using Lemma 4.2 we have
ddc max(u, a) ∧ T = ddcu ∧ T on {u > a}.
b) Assume that v ∈ PSH−(Ω). Since {u > v} =
{u > a > v}, it suffices to show
ddc max(u, v) ∧ T = ddcu ∧ T on {u > a > v}
for all a ∈ Q−. Since max(u, v) ∈ E , by a) we have
ddc max(u, v) ∧ T |{max(u,v)>a} = dd
c max(max(u, v), a) ∧ T |{max(u,v)>a}
= ddc max(u, v, a) ∧ T |{max(u,v)>a}.
(2) ddcu ∧ T |{u>a} = dd
c max(u, a) ∧ T |{u>a}.
Since max(u, v, a) = max(u, a) on set open {a > v} , we have
(3) ddc max(u, v, a) ∧ T |{a>v} = dd
c max(u, a) ∧ T |{a>v}.
Since {u > a > v} ⊂ {u > a}, {a > v}, {max(u, v) > a} and (1), (2), (3) we have
ddc max(u, v) ∧ T |{u>a>v} = dd
cu ∧ T |{u>a>v}.
The next result is an analogue of an inequality due to Demaily in [De2]
4.3. Proposition. a) u, v ∈ E such that (ddcu)n({u = v = −∞}) = 0. Then
(ddc max(u, v))n ≥ 1{u≥v}(dd
cu)n + 1{u<v}(dd
where 1E denotes the characteristic function of E.
b) Let µ be a positive measure which vanishes on all pluripolar subsets of Ω. Suppose
u, v ∈ E such that (ddcu)n ≥ µ, (ddcv)n ≥ µ. Then (ddc max(u, v))n ≥ µ.
Proof. a) For each ǫ > 0 put Aǫ = {u = v − ǫ}\{u = v = −∞}. Since Aǫ ∩ Aδ = ∅ for
ǫ 6= δ there exists ǫj ց 0 such that (dd
cu)n(Aǫj ) = 0 for j ≥ 1. On the other hand, since
(ddcu)n({u = v = −∞}) = 0 we have (ddcu)n({u = v− ǫj}) = 0 for j ≥ 1. Since Theorem
4.1 it follows that
(ddc max(u, v − ǫj))
n ≥ (ddc max(u, v − ǫj))
n|{u>v−ǫj} + (dd
c max(u, v − ǫj))
n|{u<v−ǫj}
= (ddcu)n|{u≥v−ǫj} + (dd
cv)n|{u<v−ǫj}
= 1{u≥v−ǫj}(dd
cu)n + 1{u<v−ǫj}(dd
≥ 1{u≥v}(dd
cu)n + 1{u<v−ǫj}(dd
cv)n.
Letting j → ∞ and by Remark under Theorem 5.15 in [Ce2] we get
(ddc max(u, v))n ≥ 1{u≥v}(dd
cu)n + 1{u<v}(dd
because max(u, v − ǫj) ր max(u, v) and 1{u<v−ǫj} ր 1{u<v} as j → ∞.
b) Argument as a)
4.4. Proposition. Let u1, ..., uk ∈ PSH(Ω) ∩ L
∞(Ω) and uk+1, ..., un ∈ E . Then
ddcu1 ∧ ... ∧ dd
cun = O((Cn(B))
n ) for all Borel sets B ⊂ Ω′ ⊂⊂ Ω.
B(a,r)
ddcu1 ∧ ... ∧ dd
cun = o((Cn(B(a, r)))
n ) as r → 0 for all a ∈ Ω.
where B(a, r) = {z ⊂ Cn : |z − a| < r}
Proof. We may assume that 0 ≤ uj ≤ 1 for j = 1, ..., k. On the other hand, by the remark
following Defintion 4.6 in [Ce2] we again may assume that uk+1, ..., un ∈ F .
i) For each open set B ⊂⊂ Ω, applying Proposition 3.1 we get
ddcu1 ∧ ... ∧ dd
cun =
(−h∗B)
kddcu1 ∧ ... ∧ dd
(−h∗B)
kddcu1 ∧ ... ∧ dd
(1 − u1)(dd
ch∗B)
k ∧ ddcuk+1 ∧ ... ∧ dd
(ddch∗B)
k ∧ ddcuk+1 ∧ ... ∧ dd
≤ k![
(ddch∗B)
n ∧ [
(ddcuk+1)
n ∧ ... ∧ [
(ddcun)
(by Corollary 5.6 in [Ce2])
≤ k!(e0(uk+1))
n ...(e0(un))
n .[Cn(B)]
≤ constants.[Cn(B)]
Hence
ddcu1 ∧ ... ∧ dd
cun ≤ constants.[Cn(B)]
for all Borel set B ⊂ Ω.
ii) By Proposition 3.1 we have
(−ϕ)kddcu1 ∧ ... ∧ dd
un ≤ k!
(1 − u1)(dd
ϕ)k ∧ ddcuk+1 ∧ ... ∧ dd
(ddcϕ)k ∧ ddcuk+1 ∧ ... ∧ dd
cun < +∞.
Hence (−ϕ)k ∈ L1(dd
cu1 ∧ ... ∧ dd
cun) for all ϕ ∈ F(Ω). Given a ∈ Ω let r0, R0 such that
B(a, r0) ⊂⊂ Ω ⊂⊂ B(a, R0). Then
|z − a|
≤ ga(z) ≤ log
|z − a|
for all z ∈ Ω, where ga denotes the Green function of Ω with pole at a. Since (−ga)
L1(dd
cu1 ∧ ... ∧ dd
cun), it follows that
B(a,r)
(−ga)
kddcu1 ∧ ... ∧ dd
cun → 0 as r → 0
Hence
(log r0 − log r)
B(a,r)
ddcu1 ∧ ... ∧ dd
cun ≤
B(a,r)
(−ga)
kddcu1 ∧ ... ∧ dd
cun → 0
as r → 0. This means that
B(a,r)
ddcu1 ∧ ... ∧ dd
cun = o((
log r0 − log r
)k) as r → 0
Combining this with the inequality
Cn(B(a, r),Ω) ≥ Cn(B(a, r), B(a, R0)) = (
logR0 − log r
)n = O((
log r0 − log r)n
we get
B(a,r)
ddcu1 ∧ ... ∧ dd
cun = o((Cn(B(a, r)))
The next result should be compared with Theorem 6.3 in [Ce1]
4.5. Theorem. Let u1, ..., un ∈ E . Then there exists ũ ∈ E
a such that
ddcu1 ∧ ... ∧ dd
cun = (dd
cũ)n + ddcu1 ∧ ... ∧ dd
cun|{u1=...=un=−∞}.
Proof. First, we write
ddcu1 ∧ ... ∧ dd
cun = µ + dd
cu1 ∧ ... ∧ dd
cun|{u1=...=un=−∞}.
where
µ = ddcu1 ∧ ... ∧ dd
cun|{u1>−∞}∪...∪{un>−∞}.
It is easy to see that µ ≪ Cn in every E ⊂⊂ Ω. Indeed, by Theorem 4.1 we have
ddcu1 ∧ ... ∧ dd
cun|{u1>−j} = dd
c max(u1,−j) ∧ ... ∧ dd
cun|{u1>−j}.
Hence, by Proposition 4.4 (i) it follows that ddcu1 ∧ ... ∧ dd
cun|{u1>−j} ≪ Cn in every
E ⊂⊂ Ω. Next, it remains to show that there exists ũ ∈ Ea such that µ = (ddcũ)n. Let
{Ωj} be an increasing exhaustion sequence of Ω. For each j ≥ 1 put µj = µ|Ωj . By [Åh]
there exists ũj ∈ F such that (dd
cũj)
n = µj . Notice that µj ր µ and
(ddcũj)
n ≤ µ ≤ (ddc(u1 + ... + un))
Applying the comparison principle we obtain
ũj ց ũ ≥ u1 + ... + un ∈ E .
Hence, ũ ∈ Ea and (ddcũ)n = lim
(ddcũj)
n = µ. The proof is thereby completed.
4.6. Corollary. u1, ..., un ∈ E . Then the following are equivalent
i) ddcu1 ∧ ... ∧ dd
cun ≪ Cn in every E ⊂⊂ Ω.
{u1=...=un=−∞}
ddcu1 ∧ ... ∧ dd
cun = 0.
{u1<−s,...,un<−s}∩E
ddcu1 ∧ ... ∧ dd
cun → 0 as s → +∞ for all E ⊂⊂ Ω.
Proof. Direct application of Theorem 4.5.
The comparison principle for class F was studied in [Ce3] and [H1]. By using Proposition
3.1 and Theorem 4.1 we prove a Xing type comparison principle for F
4.7. Theorem. Let u ∈ F , v ∈ E and 1 ≤ k ≤ n. Then
{u<v}
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
{u<v}
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
for all wj ∈ PSH(Ω), 0 ≤ wj ≤ 1, j = 1, ..., k, wk+1, ..., wn ∈ F and all r ≥ 1.
Proof. Let ǫ > 0. We set ṽ = max(u, v − ǫ). By a) in Proposition 3.1 we have
(ṽ − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cṽ)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
Since {u < ṽ} = {u < v − ǫ} and Theorem 4.1 we have
{u<v−ǫ}
(v− ǫ−u)kddcw1 ∧ ...∧ dd
cwn +
{u≤v−ǫ}
(r−w1)(dd
cv)k ∧ ddcwk+1 ∧ ...∧ dd
{u≤v−ǫ}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
Letting ǫ ց 0 we obtain
{u<v}
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
{u<v}
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
4.8. Corollary. Let u ∈ Ea such that u ≥ v for all functions v ∈ E satisfying (ddcu)n ≤
(ddcv)n. Then
{u<v}
(v − u)nddcw1 ∧ ... ∧ dd
{u<v}
(r − w1)(dd
{u<v}
(r − w1)(dd
for all v ∈ E , r ≥ 1 and all w1, ..., wn ∈ PSH(Ω), 0 ≤ w1, ..., wn ≤ 1.
Proof. Let {Ωj} be an increasing exhaustion sequence of relatively compact subdomains
of Ω. Set µj = 1Ωj1{u>−j}(dd
cu)n, where 1E denotes the characteristic function of E ⊂ Ω.
Applying Theorem 4.1 we have
µj = 1Ωj1{u>−j}(dd
c max(u,−j))n ≤ 1Ωj (dd
c max(u,−j))n.
Take φ ∈ E0(Ω) ∩ C(Ω̄). Put
φj = max(u,−j, ajφ)
where aj =
. Then φj = max(u,−j) on Ωj+1, φj ∈ E0 and
µj ≤ 1Ωj (dd
c max(u,−j))n = 1Ωj (dd
n ≤ (ddcφj)
By Ko lodziej’s theorem (see [Ko]) there exists uj ∈ E0 such that
(ddcuj)
n = µj = 1Ωj1{u>−j}(dd
cu)n, ∀ j ≥ 1.
for all j ≥ 1. By the comparison principle we have uj ց ũ ≥ u. On the other hand, since
(ddcu)n({u = −∞}) = 0, it follows that
(ddcuj)
n = 1Ωj1{u>−j}(dd
cu)n → (ddcu)n
weakly as j → ∞. Thus (ddcũ)n = lim
(ddcuj)
n = (ddcu)n. By the hypothesis we have
ũ = u. Applying Theorem 4.7 we get
{uj<v}
(v − uj)
nddcw1 ∧ ... ∧ dd
cwn +
{uj<v}
(r − w1)(dd
{uj<v}
(r − w1)(dd
{uj<v}
(r − w1)(dd
cu)n.
Letting j → ∞ we obtain
{u<v}
(v − u)nddcw1 ∧ ... ∧ dd
cwn +
{u<v}
(r − w1)(dd
Arguing as in Theorem 4.7 we prove a Xing type comparison principle for E .
4.9. Theorem. Let u, v ∈ E and 1 ≤ k ≤ n such that lim
[u(z) − v(z)] ≥ 0. Then
{u<v}
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
{u<v}
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
for all wj ∈ PSH(Ω), 0 ≤ wj ≤ 1, j = 1, ..., k, wk+1, ..., wn ∈ E and all r ≥ 1.
Proof. Let ǫ > 0. We set ṽ = max(u, v − ǫ). By b) in Proposition 3.1 we have
(ṽ − u)kddcw1 ∧ ... ∧ dd
cwn +
(r − w1)(dd
cṽ)k ∧ ddcwk+1 ∧ ... ∧ dd
(r − w1)(dd
u)k ∧ ddcwk+1 ∧ ... ∧ dd
Since {u < ṽ} = {u < v − ǫ} and Theorem 4.1 we have
{u<v−ǫ}
(v− ǫ−u)kddcw1 ∧ ...∧ dd
cwn +
{u≤v−ǫ}
(r−w1)(dd
cv)k ∧ ddcwk+1 ∧ ...∧ dd
{u≤v−ǫ}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
cu)k ∧ ddcwk+1 ∧ ... ∧ dd
Letting ǫ ց 0 we obtain
{u<v}
(v − u)kddcw1 ∧ ... ∧ dd
cwn +
{u<v}
(r − w1)(dd
cv)k ∧ ddcwk+1 ∧ ... ∧ dd
{u<v}∪{u=v=−∞}
(r − w1)(dd
u)k ∧ ddcwk+1 ∧ ... ∧ dd
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|
0704.0360 | Torsional oscillations of longitudinally inhomogeneous coronal loops | Astronomy & Astrophysics manuscript no. 7246 c© ESO 2019
August 20, 2019
Torsional oscillations of longitudinally inhomogeneous coronal
loops
T.V. Zaqarashvili1 & K. Murawski2
1 Georgian National Astrophysical Observatory (Abastumani Astrophysical Observatory), Kazbegi Ave. 2a, Tbilisi 0160, Georgia
2 Group of Astrophysics and Gravity Theory, Institute of Physics, UMCS, ul. Radziszewskiego 10, 20-031 Lublin, Poland
received / accepted
ABSTRACT
Aims. We explore the effect of an inhomogeneous mass density field on frequencies and wave profiles of torsional Alfvén oscillations
in solar coronal loops.
Methods. Dispersion relations for torsional oscillations are derived analytically in limits of weak and strong inhomogeneities. These
analytical results are verified by numerical solutions, which are valid for a wide range of inhomogeneity strength.
Results. It is shown that the inhomogeneous mass density field leads to the reduction of a wave frequency of torsional oscillations, in
comparison to that of estimated from mass density at the loop apex. This frequency reduction results from the decrease of an average
Alfvén speed as far as the inhomogeneous loop is denser at its footpoints. The derived dispersion relations and wave profiles are
important for potential observations of torsional oscillations which result in periodic variations of spectral line widths.
Conclusions. Torsional oscillations offer an additional powerful tool for a development of coronal seismology.
Key words. Magnetohydrodynamics (MHD) – Sun: corona – Sun: oscillations
1. Introduction
Recent space-based observations revealed a presence of vari-
ous kinds of magnetohydrodynamic (MHD) waves and oscilla-
tions in the solar corona. These observations as well as mod-
eling of MHD waves are important as these waves contribute
to the coronal heating problem (Roberts 2000) and they may
consist unique tool of a coronal seismology (Edwin & Roberts
1983, Nakariakov & Ofman 2001). Fast kink (Aschwanden et
al. 1999, Nakariakov et al. 1999, Wang & Solanki 2004) and
sausage (Nakariakov 2003, Pascoe et al. 2007) as well as slow
(de Moortel et al. 2002, Wang et al. 2003) magnetosonic oscil-
lations were observed to be associated either with or without
a solar flare. Analytical studies of these oscillations in coronal
loops were carried on over the last few decades, amongst others,
by Edwin & Roberts (1982, 1983), Poedts & Boynton (1996),
Nakariakov (2003), Van Doorsselaere et al. (2004a,b), Ofman
(2005), Verwichte et al. (2006) and Diáz et al. (2006).
Coronal loops act as natural wave guides for magnetosonic
and torsional Alfvén waves. The later are purely azimuthal os-
cillations in cylindrical geometry. In the linear regime, Alfvén
oscillations do not lead to mass density perturbations. As a re-
sult, contrary to magnetosonic waves, torsional Alfvén waves
can be observed only spectroscopically. While propagating from
the base of the solar corona along open magnetic field lines,
these waves may lead to an increase of a spectral line width with
height (Hassler et al. 1990, Banerjee et al. 1998, Doyle et al.
1998). In closed magnetic field structures, such as coronal loops,
these waves can be observed indirectly as periodic variations of
non-thermal broadening of spectral lines (Zaqarashvili 2003).
Alongside magnetosonic waves, torsional oscillations can be
used to infer, in the framework of coronal seismology, plasma
properties inside oscillating loops. These oscillations are an ideal
Send offprint requests to: T. Zaqarashvili e-mail: [email protected]
tool of coronal seismology as their phase speed depends alone on
plasma quantities within the loop, while wave speeds of magne-
tosonic oscillations are influenced by plasma conditions in the
ambient medium. Having known mass density within a loop,
coronal seismology, that is based on torsional oscillations, en-
ables to estimate a magnetic field strength. Torsional oscilla-
tions are potentially important in the context of rapid attenua-
tion of coronal loop kink oscillations (Aschwanden et al. 1999,
Nakariakov et al. 1999). One of a few suggested mechanisms
of the attenuation is a resonant absorption of fast magnetosonic
kink waves by azimuthal Alfvén waves (Ruderman & Roberts
2002). This process may lead to a formation of torsional oscil-
lations in the outer part of a loop. As a result, spotting torsional
oscillations after the kink mode was attenuated would serve as
an evidence of this attenuation mechanism.
A theoretical study of Alfvén oscillations in a coronal loop
was carried on recently by Gruszecki et al. (2007) who con-
sidered impulsively generated oscillations in two-dimensional
straight and curved magnetic field topologies. They found that
lateral leakage of Alfvén waves into the ambient corona is neg-
ligibly small. However, mass density profiles were adopted ho-
mogeneous within the loop, while the real conditions there are
much more complex.
Despite of significant achievements in a development of re-
alistic models there is still much more effort required to develop
our knowledge of wave phenomena in coronal loops. A goal of
this paper is to study the influence of inhomogeneous mass den-
sity fields on spectrum of torsional oscillations. The paper is or-
ganized as follows. Analytical solutions for torsional oscillations
in a longitudinally inhomogeneous coronal loop are presented in
Sect. 2. The numerical results are showed in Sect. 3. Guidelines
for potential observations of these oscillations are presented in
Sect. 4. This paper is concluded by a discussion and a short sum-
mary of the main results in Sect. 5.
http://arxiv.org/abs/0704.0360v1
2 T. Zaqarashvili & K. Murawski: Torsional oscillations of a coronal loop
2. Analytical model of torsional oscillations
We consider a coronal loop of its inhomogeneous mass density
̺0(z) and length 2L, that is embedded in a uniform magnetic field
B = B0ẑ. Small amplitude torsional Alfvén waves in a cylindri-
cal coordinate system (r, φ, z), in which plasma profiles depend
on a longitudinal coordinate z only, can be described by the fol-
lowing linear equations:
4π̺0(z)
, (1)
, (2)
where uφ and bφ are the velocity and magnetic field components
of Alfvén waves.
These equations can be easily cast into a single wave equa-
V2A(z)
= 0 , (3)
where VA(z) = B0/
4π̺0(z) is the Alfvén speed. Assuming that
uφ ∼ exp(iωt), where ω is a wave frequency, we get the equation
uφ = 0 . (4)
For a trapped solution uφ must satisfy line-tying boundary con-
ditions which are implemented by setting
uφ(z = ±L) = 0 . (5)
Equation (4) with condition (5) consists the well-known Sturm-
Liuville problem which solution depends on the profile of VA(z).
We model the coronal loop by a rarefied plasma at the loop apex
(at z = 0) and by a compressed plasma at the loop footpoints
(z = ±L). Specifically, we adopt
̺0(z) = ̺00
1 + α2
, (6)
where ̺00 is the mass density at the loop apex and α
2 is a param-
eter which defines a strength of the inhomogeneity. For α2 = 0
the above mass density profile corresponds to a homogeneous
loop, while for a larger value of α2 the medium is more inhomo-
geneous. Figure 1 illustrates ̺0(z) for α
2 = 50. The mass density
is described by Eq. (6) with ̺00 = 10
−12 kg m−3 and L = 25 Mm.
Note that plasma is compressed at z = ±L. Substituting Eq. (6)
into Eq. (4), we obtain
1 + α2
uφ = 0 , (7)
where VA0 = B0/
4π̺00. With a use of the notation
y ≡ uφ, x ≡
z, a ≡ −
Eq. (7) can be rewritten in the form of Weber (parabolic cylinder)
equation (Abramowitz & Stegun 1964)
y = 0 . (9)
Fig. 1. Spatial profile of the background mass density, ̺0(z), given by
Eq. (6) with α2 = 50. The mass density and length are expressed in
units of 10−12 kg m−3 and 1 Mm, respectively.
Standard solutions to this equation are called Weber (parabolic
cylinder) functions (Abramowitz & Stegun 1964)
W(a,±x) =
(coshπa)1/4
G1y1(x) ∓
2G3y2(x)
, (10)
where
, G3 =
and y1(x), y2(x) are respectively even and odd solutions to
Eq. (9)
y1(x) = 1 + a
+ · · · ,
y2(x) = x + a
+ · · · .
2.1. Two limiting solutions
Periodic solutions to Eq. (9) can be written analytically in the
limiting cases: (a) for a large value of a but a moderate value
of x; (b) for a large x but a moderate a. The first (second) case
corresponds to α2 ≪ 1 (α2 ≫ 1).
2.1.1. Weakly inhomogeneous plasma
We consider first the case of a weakly inhomogeneous mass den-
sity field, i.e. α2 ≪ 1. In this case we have
a < 0, −a≫ x2, p ≡
−a . (12)
We adopt the following expansion (Abramowitz & Stegun
1964):
W(a, x) + iW(a,−x) =
2W(a, 0) exp [vr + i(px + π/4 + vi)] , (13)
where
W(a, 0) =
, (14)
vr = −
(x/2)2
(2p)2
2(x/2)4
(2p)4
+ · · · , vi =
2/3(x/2)3
+ · · · . (15)
T. Zaqarashvili & K. Murawski: Torsional oscillations of a coronal loop 3
As a result of relation −a≫ x2 we have from Eq. (13)
W(a, x) =
2W(a, 0) exp
cos ζ , (16)
W(a,−x) =
2W(a, 0) exp
sin ζ , (17)
ζ ≡ px + π/4 +
. (18)
The general solution to Eq. (9) is
uφ = c1W(a, x) + c2W(a,−x) , (19)
where c1 and c2 are constants.
For a homogeneous loop, i.e. α2 = 0, we recognize the well
known solution
uφ ∼ c1 cos (kz + π/4) + c2 sin (kz + π/4). (20)
Here wave number k satisfies the following homogeneous dis-
persion relation:
. (21)
Line-tying boundary conditions of Eq. (5) lead then to discrete
values of the wave frequency, viz.
1 + α2/6
, n = 1, 2, 3, . . . . (22)
From this dispersion relation we infer that in a comparison to
the loop with a homogeneous mass density distribution, ̺00, the
weakly inhomogeneous mass density field results in a decrease
of a wave frequency. This reduction is a consequence of the fact
that the inhomogeneous loop is denser at its footpoints, so the
average Alfvén speed is decreased. To show this, we compare the
results for the inhomogeneous loop with the homogeneous loop
with the same average density, so that both loops contain exactly
the same mass (Andries et al. 2005). We introduce a frequency
difference
∆ωn = ωn − ω̄n , (23)
where
ω̄n =
V̄A0 =
4π ¯̺0
corresponds to the average mass density
¯̺0 =
̺0(z)dz = ̺00
. (25)
Substituting Eq. (25) into Eq. (24), we obtain
ω̄n =
1 + α2/3
. (26)
From Eqs. (23) and (26) we find that ∆ωn ≤ 0. Here we infer
that in comparison to the average mass density case the wave
frequency is reduced, but as a result of α2 ≪ 1 the frequency
reduction is small. This is in a disagreement with Fermat’s law
and with the results of Murawski et al. (2004) who showed that
sound waves experience frequency increase in a case of a space-
dependent random mass density field.
2.1.2. Strongly inhomogeneous plasma
We discuss now a strongly inhomogeneous mass density case,
i.e. α2 ≫ 1. This case corresponds to x ≫ |a|. In this limit we
get (Abramowitz & Stegun 1964)
W(a, x) =
2k/x(s1(a, x) cos(ξ) − s2(a, x) sin(ξ)) , (27)
W(a,−x) =
2/kx(s1(a, x) sin(ξ) − s2(a, x) cos(ξ)) , (28)
where
− a ln x +
argΓ(1/2 + ia)
, (29)
1 + e2πa − eπa, (30)
s1(a, x) ∼ 1 +
1!2x2
2!22x4
− · · · , (31)
s2(a, x) ∼ −
1!2x2
2!22x4
+ · · · (32)
ur + ivr = Γ(r + 1/2 + ia)/Γ(1/2 + ia), r = 2, 4, . . . . (33)
The boundary conditions of Eq. (5) lead to the discrete fre-
quency spectrum
. (34)
Here we infer that the strongly inhomogeneous mass den-
sity field results in a significant decrease of a wave frequency in
comparison to the case of the loop with the constant density, ̺00.
This wave frequency decrease is a consequence of the fact that
the inhomogeneous loop is denser at its footpoints. Substituting
Eq. (34) into Eq. (23) we find that ∆ωn > 0. This wave frequency
decrease, in a comparison to the case of an average mass density
is now in an agreement with Fermat’s law and with the results of
Murawski et al. (2004).
3. Numerical results
Numerical simulations are performed for Eqs. (1), (2) with an
adaptation of CLAWPACK which is a software package de-
signed to compute numerical solutions to hyperbolic partial dif-
ferential equations using a wave propagation approach (LeVeque
2002). The simulation region (−L, L) is covered by an uniform
grid of 600 numerical cells. We verified by convergence studies
that this grid does not introduce much numerical diffusion and as
a result it represents well the simulation region. We set reflect-
ing boundary conditions at the left and right boundaries of the
simulation region.
Figure 2 shows a spatial profile of velocity uφ(z) for α
2 = 50,
drawn at t = 1000 s (solid line). This spatial profile results from
the initial Gaussian pulse that was launched at t = 0 in the center
of the simulation region, at z = 0. It is noteworthy that the sine-
wave profile of Eq. (20), which is valid for α2 = 0 (dashed line),
is distorted by the strong inhomogeneity which takes place for
the case of α2 = 50.
As a consequence of the inhomogeneity wave period is al-
tered. Figure 3 displays wave period P vs. inhomogeneity param-
eter α2. Diamonds represent the numerical solutions while the
solid lines correspond to the analytical solution to Eqs. (22) (top
panel) and (34) (bottom panel). Wave periods were obtained by
Fourier analysis of the wave signals that were collected in time
at the fixed spatial location, z = 0. It is discernible that the nu-
merical data fits quite well to the analytical curves. A growth of
4 T. Zaqarashvili & K. Murawski: Torsional oscillations of a coronal loop
Fig. 2. Numerically evaluated velocity profile uφ at t = 1000 s for α2 =
50 (solid line). This profile corresponds to the mode number n = 1. Note
that as a result of strong inhomogeneity, uφ departs from the sine-wave
which corresponds to α2 = 0. The dashed line corresponds to Eq. (20)
with c1 = c2 = 0.5.
Fig. 3. Wave period P = ω/2π vs. α2 for the mode number n = 1.
Diamonds correspond to the numerical solutions to Eqs. (1), (2). Solid
lines are drawn with the use of the analytical solution to Eqs. (22) and
(34). The wave period is expressed in seconds.
wave period P with α2 results from wave scattering on centers
of the inhomogeneity and it can be explained on simple physical
grounds. In an inhomogeneous field wave frequency ωn of the
torsional oscillations can be estimated from the following for-
mula:
V̄A0 , (35)
where V̄A0 is the averaged Alfvén speed that is expressed by Eq.
(24). Using P = 2π/ωn we obtain
4π ¯̺0
. (36)
As ¯̺0 grows with α, the growth of P with α results in.
4. Potential observations of torsional oscillations
Torsional oscillations of a coronal loop may result in periodic
variations of spectral line non-thermal broadening (expressed by
a half line width, ∆λB, hereafter HW) (Zaqarashvili 2003). For a
homogeneous loop, HW can be expressed as
∆λB =
uVA0λ
|sin(ωnt)sin(knz)| , (37)
where u is an amplitude of oscillations, λ is a wave length of
the spectral line and c is the light speed. Periodic variations of
spectral line width depend on a height above the solar surface:
a strongest variation corresponds to the wave antinode and the
place of a lack of line width variation corresponds to the nodes
(loop footpoints). Therefore, time series of spectroscopic ob-
servations may allow to determine a wave period. Knowing a
length of the loop, we may estimate the Alfvén speed, which in
turn gives a possibility to infer the magnetic field strength in the
corona. We estimate the expected value of line width variations
which result from torsional oscillations. For a typical coronal
Alfvén speed of ∼ 800 km/s, an amplitude of linear torsional
oscillation can be ∼ 40 km/s, which consists 5% of the Alfvén
speed. For the ”green” coronal line Fe XIV (5303 Å) from Eq.
(37) we obtain
∆λB ≈ 0.7 Å. (38)
This value is about twice larger than the original thermal broad-
ening of Fe XIV line. As a consequence, torsional oscillations
can be detected in time series of the green coronal line spectra.
For a weakly inhomogeneous distribution of mass density
along a loop, Eq. (22) enables to estimate the Alfvén speed at
the loop apex with the help of the observed period of HW vari-
ation and a loop length. For a strongly inhomogeneous density
profile along a loop, Eq. (34) shows that a wave period of tor-
sional oscillations is not just the ratio of the loop length to the
Alfvén speed, but it strongly depends on the rate of inhomo-
geneity, α2. Therefore, an additional effort is required in order
to apply the method of coronal seismology for torsional oscil-
lations. A spatial variation of mass density along the loop can
be estimated by a direct measurement of spectral line intensity
variation along the loop. Then, the estimated variation can be fit-
ted to Eq. (6), and hence a value of α2 can be inferred. Eq. (34)
provides a value of VA0 at the loop summit. Another possibility
is to collect time series of spectroscopic observations at differ-
ent positions of the loop. A spatial variation of line width along
the loop may be compared to the theoretical plot of uφ (Fig. 2),
which enables to estimate α2 and consequently Alfvén speed at
the loop apex (with a use of Eqs. (22) or (34)).
T. Zaqarashvili & K. Murawski: Torsional oscillations of a coronal loop 5
5. Discussion and summary
It is commonly believed that Alfvén waves are generated in the
solar interior either by convection (granulation, supergranula-
tion) or by any other kinds of plasma flow (differential rotation,
solar global oscillations). Due to their incompressible nature,
these waves may carry energy from the solar surface to the solar
corona and therefore they may significantly contribute to coro-
nal heating and solar wind acceleration. In closed magnetic loops
the Alfvén waves may set up the standing torsional oscillations,
while in opened magnetic structures these waves may propagate
up to the solar wind. As a result, observations of Alfvén waves
can be of vital importance to the problems of plasma heating and
particle acceleration.
The Alfvén waves that propagate along open magnetic field
lines may lead to a growth of a spectral line width with height
(Hassler et al. 1990, Banerjee et al. 1998; Doyle et al. 1998).
However, at some altitudes the spectral line width reveals a sud-
den fall off (Harrison et al. 2002; O’Shea et al. 2003, 2005). This
phenomenon was recently explained by resonant energy transfer
into acoustic waves (Zaqarashvili et al. 2006).
On the other hand, the photospheric motions may set up tor-
sional oscillations in closed magnetic loop systems, which can
be observed spectroscopically as periodic variations of spectral
line width (Zaqarashvili 2003). As a result, the observation of
Alfvén waves can be used as an additional powerful tool of coro-
nal seismology; the observed period and loop mean length en-
ables to estimate the Alfvén speed within a loop, which in turn
makes it possible to infer a mean magnetic field strength.
Besides their photospheric origin, torsional Alfvén waves
can be generated in the solar corona in a process of resonant ab-
sorption of the global oscillations (Ruderman & Roberts 2002,
Goossens et al. 2002, Andries et al. 2005, Terradas et al. 2006).
These oscillations may excite Alfvén waves in the outer inho-
mogeneous part of a loop, leading to attenuation of global oscil-
lations and amplification of torsional oscillations. These Alfvén
oscillations can be detected as periodic variations of spectral line
width. As a consequence, observations of Alfvén waves can be
a key for a determination of a damping mechanism of the loop
global oscillations.
Dynamics of torsional Alfvén waves in a homogeneous loop
can be easily solved. However, real coronal loops are longitudi-
nally inhomogeneous, which leads to alteration of wave dynam-
ics (Arregui et al. 2005, 2007, Van Doorsselaere et al. 2004a,b,
Donnelly et al. 2006, Dymova & Ruderman 2006, McEwan et
al. 2006). Therefore, the dynamics of Alfvén waves in longitudi-
nally inhomogeneous coronal loops must be understood in order
to provide analytical basis for potential observations of torsional
oscillations.
In this paper we discussed by analytical and numerical
means evolution of torsional Alfvén waves in an inhomogeneous
mass density field. The analytical efforts resulted in dispersion
relations which were obtained for a specific choice of an equilib-
rium mass density profile. These dispersion relations were writ-
ten explicitly for two limiting cases: (a) weekly inhomogeneous
and (b) strongly inhomogeneous mass density fields. From these
dispersion relations we inferred that the inhomogeneity results in
a wave frequency reduction in comparison to that of estimated at
the loop summit. This analytical finding is supported by the nu-
merical data which reveals that frequency reduction takes place
outside the region of validity of the analytical approach. As a
result of that we claim that a reduction of wave frequency is
ubiquitous for the inhomogeneous mass density field we consid-
ered. This reduction is a consequence of wave scattering on in-
homogeneity centers and it results from reduction of the average
Alfvén speed within a coronal loop. This frequency reduction
has important implications as far as wave observations are con-
cerned. The analytical formulae can be used for estimation of
coronal plasma parameters and therefore torsional Alfvén waves
consist an additional powerful tool of coronal seismology.
Acknowledgments: The authors express their thanks to the
referee, Prof. S. Poedts, for his stimulating comments. The work
of T.Z. is supported by the grant of Georgian National Science
Foundation GNSF/ST06/4-098. A part of this paper is sup-
ported by the ISSI International Programme ”Waves in the Solar
Corona”.
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Introduction
Analytical model of torsional oscillations
Two limiting solutions
Weakly inhomogeneous plasma
Strongly inhomogeneous plasma
Numerical results
Potential observations of torsional oscillations
Discussion and summary
|
0704.0361 | Pseudo-random Puncturing: A Technique to Lower the Error Floor of Turbo
Codes | Pseudo-random Puncturing: A Technique to Lower
the Error Floor of Turbo Codes
Ioannis Chatzigeorgiou, Miguel R. D. Rodrigues, Ian J. Wassell
Digital Technology Group, Computer Laboratory
University of Cambridge, United Kingdom
Email: {ic231, mrdr3, ijw24}@cam.ac.uk
Rolando Carrasco
School of EE&C Engineering
University of Newcastle, United Kingdom
Email: [email protected]
Abstract— It has been observed that particular rate-1/2
partially systematic parallel concatenated convolutional codes
(PCCCs) can achieve a lower error floor than that of their
rate-1/3 parent codes. Nevertheless, good puncturing patterns
can only be identified by means of an exhaustive search,
whilst convergence towards low bit error probabilities can be
problematic when the systematic output of a rate-1/2 partially
systematic PCCC is heavily punctured. In this paper, we present
and study a family of rate-1/2 partially systematic PCCCs, which
we call pseudo-randomly punctured codes. We evaluate their bit
error rate performance and we show that they always yield
a lower error floor than that of their rate-1/3 parent codes.
Furthermore, we compare analytic results to simulations and
we demonstrate that their performance converges towards the
error floor region, owning to the moderate puncturing of their
systematic output. Consequently, we propose pseudo-random
puncturing as a means of improving the bandwidth efficiency
of a PCCC and simultaneously lowering its error floor.
I. INTRODUCTION
Although in certain applications, such as satellite
communications, link reliability is of essence and low
rate codes are used to support it, bandwidth occupancy is
more important in wireless communications and hence high
rate codes are preferred. A high rate convolutional code can
be obtained by periodic elimination, known as puncturing,
of particular codeword bits from the output of a parent low
rate convolutional encoder. Extensive analyses on punctured
convolutional codes have shown that their performance is
always inferior to the performance of their low rate parent
codes (e.g. see [1], [2]).
The performance of punctured parallel concatenated
convolutional codes (PCCCs), also known as punctured turbo
codes, has also been investigated. Design considerations have
been derived by analytical [3]–[5] as well as simulation-based
approaches [6]–[8], while upper bounds on the bit error
probability (BEP) were evaluated in [5], [9]. Punctured turbo
codes are usually classified as systematic, partially systematic
or non-systematic depending on whether all, some or none of
their systematic bits are transmitted [7]. Recent papers [7]–[9]
have demonstrated that partially systematic PCCCs yield lower
error floors than systematic PCCCs of the same rate.
In [10] we showed that rate-1/2 non-systematic PCCCs
can achieve error floors, which are lower even than those
of their rate-1/3 parent PCCCs. This interesting outcome is
valid when maximum-likelihood (ML) decoding is employed.
When suboptimal iterative decoding is used, the absence of
received systematic bits causes erroneous decisions, which
prohibit the iterative decoder from converging to the error
floor. Nevertheless, we demonstrated that rate-1/2 child codes,
whose BEP performance converges towards an error floor
which is lower than that of their rate-1/3 parent PCCC, can
still be found by means of an exhaustive search. During
this process, the union bound on the BEP of each rate-1/2
punctured PCCC is computed and compared to the union
bound of the rate-1/3 parent PCCC. Note that the union bound
coincides with the error floor of the code for high values of
Eb/N0 [11]. Punctured PCCCs that achieve a bound lower
than that of their rate-1/3 parent PCCC are selected.
Computation of the exact union bound on the BEP of
a punctured PCCC becomes intensive as the interleaver
size increases. In [12] we presented a simple technique
to approximate the union bound of a turbo code and we
demonstrated that this approximation is very accurate when
a large interleaver size is used. We used our technique to
identify a family of rate-1/2 partially systematic PCCCs, which
we called pseudo-randomly punctured PCCCs (PRP-PCCCs).
Although we did not explore their BEP performance in detail,
we observed that particular PRP-PCCC configurations could
achieve a lower error floor than that of their parent codes.
This paper builds upon the work carried out in [10] and [12].
Initially, we provide analytical expressions for the parameters
that influence the bit error performance of PCCCs. We then
evaluate those parameters and compute the union bound
approximations for both rate-1/3 parent PCCCs and rate-1/2
PRP-PCCCs. We demonstrate that the latter always exhibit
a lower error floor than the former, when large interleaver
sizes are considered. In order to verify our theoretical analysis,
we compare analytic results to simulations for specific PCCC
configurations. The paper concludes with a summary of the
main contributions.
II. PERFORMANCE EVALUATION OF PCCCS
Turbo codes, in the form of symmetric rate-1/3 PCCCs,
consist of two identical rate-1/2 recursive systematic
convolutional encoders separated by an interleaver of size N
[13]. The information bits are input to the first constituent
convolutional encoder, while an interleaved version of the
information bits are input to the second convolutional encoder.
http://arxiv.org/abs/0704.0361v1
The output of the turbo encoder consists of the systematic bits
of the first encoder, which are identical to the information bits,
the parity check bits of the first encoder and the parity check
bits of the second encoder.
The bit error probability Pb of a PCCC employing ML
soft decoding, on an additive white Gaussian noise (AWGN)
channel, is upper bounded as follows
Pb ≤ P
b (1)
where the union bound P u
is defined as
P ub =
P (w). (2)
Here, the sum runs over all possible values of input
information weight w, with P (w) being the contribution to
the union bound P u
of only those codeword sequences which
were generated by input sequences of a specific information
weight w. An individual contribution P (w) is given by [11],
P (w) =
Bw,dQ
2R · Eb
, (3)
where N is the interleaver size, R is the code rate of the turbo
encoder and Bw,d denotes the number of codeword sequences
having overall output weight d, which were generated by input
information sequences of weight w.
In [11] it was shown that the union bound on the BEP of
a PCCC using a uniform interleaver of size N coincides with
the average of the union bounds obtainable from the whole
class of deterministic interleavers of size N . For small values
of N , the union bound can be very loose compared with the
actual performance of turbo codes using specific deterministic
interleavers. However, for N≥1000, it has been observed that
randomly generated interleavers generally perform better than
deterministic interleaver designs [15]. Consequently, the union
bound provides a good indication of the actual bit error rate
performance of a PCCC operating in the error floor region,
when long interleavers are considered.
Derivation of all coefficients Bw,d becomes a
computationally intensive process as the interleaver size
increases, especially when punctured PCCCs are considered
[12]. However, the union bound can be approximated as
follows
P ub ≈ P (w=2), (4)
when long interleavers are used. This approximation is based
on a number of observations:
1) Codeword sequences, which were generated by input
sequences having the minimum possible information
weight, become the main contributors to the bit error rate
performance, as the size N of the interleaver increases
[12], [16].
2) Owning to the structure of the constituent encoders, the
minimum information weight of an input sequence is
always equal to two [16].
Therefore, P (w=2) is the dominant contribution to the union
bound over a broad range of bit error probabilities [12], [16]
and can be used to predict the error floor of turbo codes.
Throughout this paper, we use the union bound
approximation as the basis for the analytic performance
comparison of turbo codes. In particular, if P and P ′ are two
PCCCs using long interleavers of identical size, we say that
P yields a lower error floor than that of P ′ when their bound
approximations, PP(2) and PP
(2) respectively, satisfy
PP(2) < PP
(2). (5)
The above condition can be expanded using (3) as follows
2RP Eb
A. (6)
It was demonstrated in [16] that the free effective distance, df,
which conveys the minimum weight of a codeword sequence
for a weight-2 input information sequence, has a major impact
on the performance of a turbo code. Consequently, if dPf
and dP
f denote the free effective distances of P and P
respectively, condition (6) collapses to
2RP Eb
A, (7)
which only considers the first non-zero, that is the most
significant, term of each sum.
Function Q(ξ) is a monotonically decreasing function of
ξ, where ξ is a real number. Therefore, if ξ1 and ξ2 are real
numbers, with ξ1 > ξ2, we deduce that Q(ξ1) < Q(ξ2), and
vice versa, i.e,
Q(ξ1) < Q(ξ2) ⇔ ξ1 > ξ2. (8)
Consequently, inequality (7) reduces to
RPdPf > R
f , (9)
BP2,df ≤ B
. (10)
When the code rates are equal, the free effective distance of
turbo codes plays a role similar to that of the free distance
of convolutional codes, since the performance criterion (9) is
simplified to
dPf > d
f . (11)
Expressions (9) and (10) will be the basis for the comparison
of the BEP performance in the error floor region of two
PCCCs.
III. DETERMINATION OF PARAMETERS THAT INFLUENCE
THE PERFORMANCE OF TURBO CODES
We will now determine the various parameters that affect
performance for two classes of turbo codes: conventional
rate-1/3 PCCCs and pseudo-randomly punctured rate-1/2
PCCCs. The turbo codes considered throughout this paper are
symmetric, i.e., the two constituent encoders are identical.
A. Rate-1/3 PCCCs
Criteria (9) and (10) require knowledge of the free effective
distance df and the coefficient B2,df of each PCCC. In the
remainder of the paper, we use the abbreviation “Par” to denote
a rate-1/3 parent PCCC. Its free effective distance dParf can
be expressed as the sum of the minimum weight dmin of the
codeword sequence generated by the first constituent encoder,
and the minimum weight zmin of the parity check sequence
generated by the second constituent encoder, when a sequence
of information weight w=2 in input to the PCCC
dParf = dmin + zmin. (12)
Taking into account that the turbo codes are symmetric and
the weight umin of the systematic output sequence is always 2
since w=2, we can write
dParf = (umin + zmin) + zmin = 2 + 2zmin. (13)
The number BPar
of codeword sequences, generated by
a turbo encoder using a uniform interleaver of size N , can
be associated with the number B2,dmin of codeword sequences
having weight dmin, generated by the first constituent encoder,
and the number B2,zmin of parity check sequences having
weight zmin, generated by the second constituent encoder, if we
elaborate on the expressions described in [11]. In particular,
we obtain
B2,dmin · B2,zmin
) , (14)
where B2,dmin and B2,zmin return the same value, since they
both consider the same trellis paths. Note that the first index
in the above notations refers to the input information weight,
which is two.
It was shown in [16] that good rate-1/3 PCCCs are obtained
when their feedback generator polynomial GR is chosen to be
primitive, whilst their feedforward generator polynomial GF
is different than GR. The period L of a primitive polynomial
is given by [17]
L = 2ν − 1, (15)
where ν is the order of the polynomial, or equivalently, the
memory size of each constituent code.
We demonstrated in [12] that when a primitive feedback
generator polynomial is used, the minimum weight zmin and
the coefficient B2,2,zmin can be expressed as
zmin = 2
ν−1 + 2,
B2,dmin = B2,zmin = N − L,
respectively. Consequently, expression (13) assumes the form
dParf = 6 + 2
ν , (17)
whilst, if we combine (14) and (16), the coefficient BPar
be expressed as a function of the intlerleaver size N and the
period L, as follows
2(N − L)2
N(N − 1)
. (18)
In the special case when the size N of the interleaver is an
integer multiple of the period L of the feedback generator
polynomial, i.e., N=µL, we can rewrite (18) as
2L(µ− 1)2
µ(µL− 1)
. (19)
B. Rate-1/2 Pseudo-randomly Punctured PCCCs
A high rate PCCC can be obtained by periodic elimination
of specific codeword bits from the output of a rate-1/3 parent
PCCC. A puncturing pattern P can be represented by a 3×M
matrix as follows:
p1,1 p1,2 . . . p1,M
p2,1 p2,2 . . . p2,M
p3,1 p3,2 . . . p3,M
, (20)
where M is the puncturing period and pi,m ∈ {0, 1}, with
i=1, 2, 3 and m=1, . . . ,M . For pi,m=0 the corresponding
output bit is punctured, otherwise it is transmitted. The first
and second rows of the pattern are used to puncture the
systematic and parity check outputs, respectively, of the first
constituent encoder. The third row determines which parity
check bits from the output of the second constituent encoder
will be punctured.
Pseudo-random puncturing has been described in [12], in
detail. It is applied to rate-1/3 PCCCs, which use primitive
feedback generator polynomials, hence the polynomial period
L is also given by (15). The puncturing pattern can
be constructed once the parity check sequence y =
(y0, y1, . . . , yL) for an input sequence x = (1, 0, . . . , 0) of
length L+1, has been obtained at the output of the first
constituent encoder. As long as a trail of zeros follows the
first non-zero input bit, the component encoder behaves like a
pseudo-random generator, hence the parity check bits from y1
to yL form a pseudo-random sequence. We set the elements
of the second row of the puncturing pattern to be equal to
the bits of this pseudo-random sequence, but circularly shifted
rightwards by one, i.e., p2,m+1 = ym for m=1, . . . , L. Note
that in pseudo-random puncturing, the puncturing period M is
equal to the period L of the feedback polynomial, i.e., M=L.
The first row of the pattern is set to be the complement of the
second row, thus p1,m=1− p2,m. In order to achieve a code
rate of 1/2, we do not puncture the parity check output of the
second constituent encoder, hence all the elements of the third
row are set to one, i.e., p3,m=1.
As an example, let us consider a rate-1/3 PCCC with
generator polynomials (GF , GR)=(5, 7)8 in octal form. The
memory size of each constituent encoder is ν = 2, thus the
period of GR is found to be L = 2
2−1 = 3. Consequently,
we set the input sequence to (1, 0, 0, 0) and we obtain the
parity check sequence (1, 1, 1, 0) at the output of the first
constituent encoder. The block of the last L=3 parity check
bits, i.e., (1, 1, 0), forms a pseudo-random sequence. If we
circularly shift the bits of this pseudo-random sequence to the
right by one and map them to the elements of the second row
of the puncturing pattern, we obtain [0 1 1]. Eventually the
puncturing pattern, based on which the rate-1/2 PRP-PCCC is
generated from the rate-1/3 parent PCCC, assumes the form
1 0 0
0 1 1
1 1 1
. (21)
We emphasize that the puncturing pattern depends on the
generator polynomials of the rate-1/3 parent PCCC, hence
different polynomials yield different puncturing patterns.
Furthermore, a rate-1/2 PRP-PCCC can be obtained only if the
parent PCCC uses primitive feedback generator polynomials.
We have previously determined [12] the minimum weight
d′min of the codeword sequence generated by the first
constituent encoder, when a sequence of information weight
w = 2 in input to the rate-1/2 PRP-PCCC. In particular, we
found that
d′min = 2
ν−2 + 2. (22)
The parity check sequence generated by the second constituent
encoder is not punctured, thus its minimum weight is also
given by (16). Therefore, we can compute the free effective
distance dPRPf of a rate-1/2 PRP-PCCC as follows
dPRPf = d
min + zmin
= (2ν−2 + 2) + 2ν−1 + 2
= 4 + 3(2ν−2).
Every time a particular column m of the puncturing
pattern is active during the N time steps of the coding
process, codeword sequences having minimum weight d′min
are generated. Their exact number, Am, can be computed
using the expressions in [12]. In particular, we find that for
M =L the number of minimum-weight codeword sequences
Am, generated when column m is active, is given by
⌊N/M⌋ − 1, if (N mod M)<m
⌊N/M⌋ , otherwise,
where (ξ1 mod ξ2) denotes the remainder of division of ξ1
by ξ2, and ⌊ξ⌋ denotes the integer part of ξ. In order to
facilitate our analysis, we assume that the interleaver size N is
an integer multiple of the puncturing period M , i.e., N=µM ,
where µ is a positive integer. Hence, (24) collapses to
Am = µ− 1, (25)
since (N mod M) is always zero and m>0.
It has been demonstrated in [12] that minimum-weight
codeword sequences can be obtained only when the active
column m is in the range 2 ≤ m ≤ M ; every time one
of these M−1 columns of the puncturing pattern is active,
Am minimum-weight codeword sequences are generated.
Consequently, the total number of codeword sequences having
weight d′min assumes the value
B2,d′
= (M − 1)Am, (26)
or, equivalently
B2,d′min = (L− 1)(µ− 1), (27)
where M has been replaced by L, since they are equal
quantities and they can be used interchangeably.
Similarly to the second constituent encoder of the rate-1/3
parent PCCC, the second constituent encoder of the rate-1/2
PRP-PCCC also generates a total of B2,zmin sequences having
weight zmin, since its parity check output is not punctured.
Consequently, the coefficient BPRP
or a rate-1/2 PRP-PCCC
can be expressed as
BPRP2,df =
B2,d′min · B2,zmin
[(L− 1)(µ− 1)] · (N − L)
2(L− 1)(µ− 1)2
µ(µL− 1)
invoking (14), which can be used when PCCCs employing
uniform interleavers of size N are considered.
IV. PERFORMANCE COMPARISON OF ANALYTIC TO
SIMULATION RESULTS
Having evaluated the parameters that influence the
performance of the PCCCs under investigation, we are now in
the position to explore whether a rate-1/2 PRP-PCCC exhibits
a lower bound approximation than that of its rate-1/3 parent
PCCC. We observe that dPRPf can be expressed in terms of
dParf , if we subtract (17) from (23)
dPRPf = d
f − (2 + 2
ν−2). (29)
Coefficient BPRP
can also be represented in terms of BPar
we divide (28) by (19)
BPRP2,df =
BPar2,df . (30)
According to (9) and (10), if both conditions
dPRPf >
dParf (31)
< BPar
are satisfied, a rate-1/2 PRP-PCCC yields a lower bound
approximation than that of its rate-1/3 parent code. We deduce
from (30) that BPRP
is always less than BPar
, thus the second
condition holds true. The first condition assumes the following
form, if we substitute dPRPf with its equivalent, based on (29),
dParf > 6 + 3(2
ν−2). (33)
Nevertheless, we have shown in (17) that the free effective
distance of the parent PCCC is given by dParf =6+ 2
ν , which
can be rewritten as dParf = 6 + 4(2
ν−2). Therefore, dParf is
always greater than 6 + 3(2ν−2), and hence, both conditions
are satisfied.
The outcome of this investigation reveals that rate-1/2
PRP-PCCCs using long interleavers are always expected to
0 1 2 3 4 5 6
(dB)
Simulation, Rate−1/3 PCCC (ν=2)
Bound Approx., Rate−1/3 PCCC (ν=2)
Simulation, Rate−1/2 PRP−PCCC (ν=2)
Bound Approx., Rate−1/2 PRP−PCCC (ν=2)
Simulation, Rate−1/3 PCCC (ν=3)
Bound Approx., Rate−1/3 PCCC (ν=3)
Simulation, Rate−1/2 PRP−PCCC (ν=3)
Bound Approx., Rate−1/2 PRP−PCCC (ν=3)
Fig. 1. Comparison of bound approximations to simulation results. The
exact log-MAP algorithm is applied over 8 iterations and an interleaver size
of 1, 000 bits is used.
yield a lower bound approximation, or equivalently a lower
error floor, than that of their rate-1/3 parent codes.
Fig.1 compares bound approximations to simulation results
for rate-1/3 parent PCCCs and rate-1/2 PRP-PCCCs of
memory size ν =2 and ν =3, over the AWGN channel. For
ν=2, the generator polynomials of the PCCCs are taken to be
(GF , GR)=(5, 7)8, whilst for ν=3, the PCCCs are described
by (GF , GR) = (17, 15)8. The component decoders employ
the conventional exact log-MAP algorithm [18]. A moderate
interleaver size of 1, 000 bits has been chosen, so as to allow
the bit error rate performance of the PCCCs to approach the
corresponding bound approximations at BEPs in the region of
10−6 to 10−7.
As expected, Fig.1 confirms that for high values of Eb/N0,
the BEP of each rate-1/2 PRP-PCCC is indeed lower than
that of the corresponding rate-1/3 parent code, whilst after 8
iterations the performance curves of all turbo codes approach
the respective bound approximation curves.
V. CONCLUSION
In previous work [9], [10], [12] we introduced techniques
to evaluate the performance of punctured PCCCs and we
observed that, in some cases, the error floor could be lowered
by reducing the rate of a PCCC from 1/3 to 1/2. Nevertheless,
good puncturing patterns were identified by means of an
exhaustive search, whilst convergence towards low bit error
probabilities of those rate-1/2 PCCCs whose systematic output
was heavily punctured, had to be investigated.
In this paper, we established that rate-1/2 pseudo-randomly
punctured PCCCs, which form a subset of rate-1/2 partially
systematic PCCCs, not only approach the error floor region
for an increasing number of iterations but always yield a
lower error floor than that of their rate-1/3 parent codes.
Consequently, pseudo-random puncturing can be used to
reduce the rate of a PCCC from 1/3 to 1/2 and at the same
time achieve a coding gain at low bit error probabilities.
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Introduction
Performance Evaluation of PCCCs
Determination of Parameters that Influence the Performance of Turbo Codes
Rate-1/3 PCCCs
Rate-1/2 Pseudo-randomly Punctured PCCCs
Performance Comparison of Analytic to Simulation Results
Conclusion
References
|
0704.0362 | The Arctic Circle Revisited | arXiv:0704.0362v1 [math-ph] 3 Apr 2007
The Arctic Circle Revisited
F. Colomo and A.G. Pronko
Abstract. The problem of limit shapes in the six-vertex model with domain
wall boundary conditions is addressed by considering a specially tailored bulk
correlation function, the emptiness formation probability. A closed expression
of this correlation function is given, both in terms of certain determinant and
multiple integral, which allows for a systematic treatment of the limit shapes
of the model for full range of values of vertex weights. Specifically, we show
that for vertex weights corresponding to the free-fermion line on the phase
diagram, the emptiness formation probability is related to a one-matrix model
with a triple logarithmic singularity, or Triple Penner model. The saddle-point
analysis of this model leads to the Arctic Circle Theorem, and its generalization
to the Arctic Ellipses, known previously from domino tilings.
1. Introduction
The Arctic Circle has first appeared in the study of domino tilings of large
Aztec diamonds [EKLP, JPS]. The name originates from the fact that in most
configurations the dominoes are ‘frozen’ outside the circle inscribed into the dia-
mond, while the interior of the circle is a disordered, or ‘temperate’, zone. Further
investigations of the domino tilings of Aztec diamonds, such as details of statistics
near the circle, can be found in [CEP, J1, J2]. Here we mention that the Arctic
Circle is a particular example of a limit shape in dimer models, in the sense that
it describes the shape of a spatial phase separation of order and disorder. Apart
from domino tilings, many more examples have been discussed recently, see, among
others, papers [CKP, CLP, KO, KOS, OR].
As long as only dimer models are considered, this amounts to restrict to dis-
crete free-fermionic models, although with nontrivial boundary conditions. Indeed,
many of them can be viewed as a six-vertex model at its Free Fermion point (the
correspondence being however usually not bijective), with suitably chosen fixed
boundary conditions. In particular, this is the case of domino tilings of Aztec dia-
monds [EKLP], and the corresponding boundary conditions of the six-vertex model
are the so-called Domain Wall Boundary Conditions (DWBC). Hence the problem
of limit shapes extends to the six-vertex model with generic weights, and with fixed
boundary conditions, among which the case of DWBC is the most interesting.
Historically, the six-vertex model with DWBC was first considered in paper [K]
within the framework of Quantum Inverse Scattering Method [KBI] to prove the
Gaudin hypothesis for norms of Bethe states. The model was subsequently solved
2000 Mathematics Subject Classification. 15A52, 82B05, 82B20, 82B23.
http://arxiv.org/abs/0704.0362v1
2 F. COLOMO AND A.G. PRONKO
in paper [I] where a determinant formula for the partition function was given; see
also [ICK] for a detailed exposition. Quite independently, the model was later
found, under certain restrictions on the vertex weights, to be deeply related with
enumerations of alternating sign matrices (see, e.g., [Br] for a review) and, as
already mentioned, to domino tilings of Aztec diamonds [EKLP].
Concerning the problem of limit shapes for the six-vertex model with DWBC,
as far as the Free Fermion point is considered, the relation with domino tilings
provided apparently an indirect proof of the corresponding Arctic Circle. The non-
bijective nature of the correspondence between the two models asked for more direct
results, purposely for the free-fermion six-vertex model, see [Zi1, FS, KP]. Out
of the Free Fermion point, however, only very few analytical results are available,
such as exact expressions for boundary one-point [BPZ] and two-point [FP, CP1]
correlation functions. The present knowledge on the subject is based mainly on
numerics [E, SZ, AR]; some steps towards finding the limit shapes of the model
have been done recently in [PR].
In the present note we propose a rather direct strategy to address the problem:
after briefly reviewing the six-vertex model with DWBC, we define a bulk corre-
lation function, the Emptiness Formation Probability (EFP), which discriminates
the ordered and disordered phase regions. We give for this correlation function two
equivalent representations, in terms of a determinant and of a multiple integral.
The core derivation of EFP is heavily based on the Quantum Inverse Scattering
Method [KBI], along the lines of papers [BPZ, CP1]; it is out of the scope of the
present paper, corresponding details being given in a separate publication [CP4].
Here our aim is to demonstrate how the limit shapes for the considered model
can be extracted from EFP in a suitable scaling limit, by making use of ideas and
techniques of Random Matrix Models.
To be more specific, and to establish a contact with previous results, we spe-
cialize here our further discussion to the case of free-fermion six-vertex model. We
show that the asymptotic analysis of multiple integral formula for EFP in the scal-
ing limit reduces to a saddle-point problem for a one-matrix model with a triple
logarithmic singularity, or triple Penner model. We argue that the limit shape cor-
responds to condensation of all saddle-point solutions to a single point. This allows
us to recover the known Arctic Circle and Ellipses.
As a comment to our approach, it is to be stressed that it is directly tailored on
the six-vertex model, rather than domino tilings. For this reason it is not restricted
to the free-fermion models, even if, of course, further significant efforts might be
necessary, essentially from the point of view of Random Matrix Model reformu-
lation, for application to more general situations. On the basis of our previous
results in [CP2], however, the application of the method to the particular case of
the so-called Ice Point of the model should be straightforward. This would provide
the limit shape of alternating sign matrices.
2. The model
2.1. The six-vertex model. The six-vertex model (for reviews, see [LW,
Ba]) is formulated on a square lattice with arrows lying on edges, and obeying the
so-called ‘ice-rule’, namely, the only admitted configurations are such that there are
always two arrows pointing away from, and two arrows pointing into, each lattice
vertex. An equivalent and graphically simpler description of the configurations of
THE ARCTIC CIRCLE REVISITED 3
w1 w2 w3 w4 w5 w6
Figure 1. The six allowed types of vertices in terms of arrows
and lines, and their Boltzmann weights.
Figure 2. A possible configuration of the six-vertex model with
DWBC at N = 4, in terms of arrows and lines.
the model can be given in terms of lines flowing through the vertices: for each arrow
pointing downward or to the left, draw a thick line on the corresponding edge. This
line picture implements the ‘ice-rule’ in an automated way. The six possible vertex
states and the Boltzmann weights w1, w2, . . . , w6 assigned to each vertex according
to its state are shown in Figure 1.
2.2. Domain Wall Boundary Conditions. The Domain Wall Boundary
Conditions (DWBC) are imposed on the N×N square lattice by fixing the direction
of all arrows on the boundaries in a specific way. Namely, the vertical arrows on the
top and bottom of the lattice point inward, while the horizontal arrows on the left
and right sides point outward. Equivalently, a generic configuration of the model
with DWBC can be depicted by N lines flowing from the upper boundary to the
left one. A possible state of the model both in terms of arrows and of lines is shown
in Figure 2.
2.3. Partition function. The partition function is defined, as usual, as a
sum over all possible arrow configurations, compatible with the imposed DWBC,
each configuration being assigned its Boltzmann weight, given as the product of all
the corresponding vertex weights,
arrow configurations
with DWBC
wn11 w
2 . . . w
Here n1, n2, . . . , n6 denote the numbers of vertices with weights w1, w2, . . . , w6,
respectively, in each arrow configuration (n1 + n2 + · · ·+ n6 = N
2.4. Anisotropy parameter and phases of the model. The six-vertex
model with DWBC can be considered, with no loss of generality, with its weights
invariant under the simultaneous reversal of all arrows,
w1 = w2 =: a , w3 = w4 =: b , w5 = w6 =: c .
4 F. COLOMO AND A.G. PRONKO
Under different choices of Boltzmann weights the six-vertex model exhibits different
behaviours, according to the value of the parameter ∆, defined as
a2 + b2 − c2
It is well known that there are three physical regions or phases for the six-vertex
model: the ferroelectric phase, ∆ > 1; the anti-ferroelectric phase, ∆ < −1; the
disordered phase, −1 < ∆ < 1. Here we restrict ourselves to the disordered phase,
where the Boltzmann weights are conveniently parameterized as
a = sin(λ+ η) , b = sin(λ− η) , c = sin 2η . (2.1)
With this choice one has ∆ = cos 2η. The parameter λ is the so-called spectral
parameter and η is the crossing parameter. The physical requirement of positive
Boltzmann weights, in the disordered regime, restricts the values of the crossing
and spectral parameters to 0 < η < π/2 and η < λ < π − η.
The special case η = π/4 (or ∆ = 0) is related to free fermions on a lattice, and
there is a well-known correspondence with dimers and domino tilings. In particular,
at λ = π/2, the ∆ = 0 six-vertex model with DWBC is related to the domino tilings
of Aztec diamond. For arbitrary λ ∈ [π/4, 3π/4], we shall refer to the ∆ = 0 case
as the Free Fermion line.
The case η = π/6 (i.e. ∆ = 1/2) and λ = π/2, where all weights are equal,
a = b = c, is known as the Ice Point; all configurations are given the same weight.
In this case there is a one to one correspondence between configurations of the
model with DWBC and N ×N alternating sign matrices.
2.5. Phase separation and limit shapes. The six-vertex model exhibits
spatial separation of phases for a wide choice of fixed boundary conditions, and,
in particular, in the case of DWBC. Roughly speaking, the effect is related to the
fact that ordered configurations on the boundary can induce, through the ice-rule,
a macroscopic order inside the lattice.
The notion of phase separation acquires a precise meaning in the scaling limit,
that is the thermodynamic/continuum limit, performed by sending the number of
sites N to infinity and the lattice spacing to zero, while keeping the total size of the
lattice fixed, e.g., to 1. On a finite lattice, several macroscopic regions may appear,
which in the scaling limit are expected to be sharply separated by some curves, the
so-called Arctic curves.
For the six-vertex model with DWBC the shape of the Artic curve, or limit
shape, has been found rigorously only on the Free Fermion line, and for the closely
related domino tilings of Aztec diamond [JPS, CEP, Zi1, FS, KP]. For generic
values of weights the limit shapes are not known, but the whole picture is strongly
supported both numerically [E, SZ, AR] and analytically [KZ, Zi2, BF, PR].
3. Emptiness Formation Probability
3.1. Definition. We shall use the following coordinates on the lattice: r =
1, . . . , N labels the vertical lines from right to left; s = 1, . . . , N labels the horizontal
lines from top to bottom. We may now introduce the correlation function FN (r, s),
measuring the probability for the first s horizontal edges between the r-th and
r+1-th line to be all ‘full’ (i.e. thick in the line picture, or with a left arrow in the
THE ARCTIC CIRCLE REVISITED 5
Figure 3. Emptiness Formation Probability. The sum in (3.1) is
performed over all configurations compatible with the drawn ar-
rows.
standard picture of the six-vertex model):
FN (r, s) =
‘constrained’
arrow configurations
with DWBC
wn11 w
2 . . . w
6 . (3.1)
Here the sum is performed over all arrow configurations on the N × N lattice,
subjected to the restriction of DWBC, and to the condition that all arrows on the
first s edges between the r-th and r + 1-th line should point left, see Figure 3.
Although this correlation function may appear rather sophisticated, it is com-
putable in some closed form by means of the Quantum Inverse Scattering Method,
on which DWBC are indeed tailored. It is the natural adaptation of the Empti-
ness Formation Probability of quantum spin chains to the present model. For this
reason, and to link to the common practice in the quantum integrable models com-
munity, even if FN (r, s) actually describes ‘fullness’ formation probability, we shall
call it Emptiness Formation Probability (EFP).
3.2. Qualitative discussion of FN (r, s). Let us restrict ourselves to the dis-
ordered regime, −1 < ∆ < 1, for definiteness. From previous analytical and nu-
merical work, in the large N limit the emergence of a limit shape, in the form of
a continuous closed curve touching once each of the four sides of the lattice, is ex-
pected. It follows that five regions emerge in the lattice: a central region, enclosed
by the curve, and four corner regions, lying outside the closed curve and delimited
by the sides of the lattice. The central region is disordered, while the four corners
are frozen, with mainly vertices of type 1, 3, 2, 4 (see Figure 1) appearing in the
top-left, top-right, bottom-right and bottom-left corner, respectively.
By construction, EFP is expected to be almost one in frozen regions of type 1,
or 3, bordering the top side of the lattice, and to be rather small otherwise. DWBC
exclude a region of type 3 to emerge in the upper part of the lattice. Hence FN (r, s)
describes, at a given value of r, as s increases, a transition from a frozen region of
vertices of type 1, where FN (r, s) ∼ 1, to a generic region where FN (r, s) ∼ 0.
It follows that FN (r, s) can describe only the upper left portion of the closed
curve, between its top and left contact points. Nevertheless, it should be mentioned
that the full curve can be built from the knowledge of its top left portion, just
6 F. COLOMO AND A.G. PRONKO
exploiting the crossing symmetry of the six-vertex model. Hence EFP, FN (r, s), is
well suited to describe limit shapes.
3.3. Some notations. For a given choice of parameters λ, η we define
sin 2η
sin(λ+ η) sin(λ− η)
and the integration measure on the real line
µ(x) := ex(λ−π/2)
sinh(ηx)
sinh(πx/2)
related to ϕ as follows:
µ(x) dx .
Let us introduce the complete set of monic orthogonal polynomial {Pn(x)}n=0,1,...
associated to the integration measure µ(x), with the orthogonality relation
Pn(x)Pm(x)µ(x) dx = hnδnm .
The square norms hn are completely determined by the measure µ(x), and may
be expressed, in principle, in terms of its moments. In the following we shall be
interested in the complete set of orthogonal polynomials {Kn(x)}n=0,1,... defined as
Kn(x) = n!ϕ
n+1 1
Pn(x) .
We moreover define
ω(ǫ) :=
sin(ǫ)
sin(ǫ− 2η)
, ω̃(ǫ) :=
sin(ǫ)
sin(ǫ + 2η)
Note that the following relation holds
a2 ω̃ − 2∆ab ω̃ω + b2 ω = 0 , (3.2)
allowing to express ω̃ in terms of ω.
3.4. Determinant representation. For EFP in the six-vertex model with
DWBC, the following representation holds:
FN (r, s) = (−1)
s det
1≤j,k≤s
KN−k(∂ǫj )
[ω(ǫj)]
[ω(ǫj)− 1]
1≤j<k≤s
[ω̃(ǫj)− 1] [ω(ǫk)− 1]
ω̃(ǫj)ω(ǫk)− 1
ǫ1=0,...,ǫs=0
. (3.3)
This representation has been obtained in the framework of the Quantum Inverse
Scattering Method [KBI], along the lines of analogous derivations worked out for
one-point and two-point boundary correlation functions of the model [BPZ, CP1].
The details of the derivation can be found in [CP4].
THE ARCTIC CIRCLE REVISITED 7
3.5. The boundary correlation function. If we consider expression (3.3)
when s = 1, we recover the boundary polarization, introduced and computed in
[BPZ]. It is convenient to consider the closely related boundary correlation function
HN (r) := FN (r, 1)− FN (r − 1, 1) .
As shown in [BPZ, CP1], the following representation holds:
HN (r) = KN−1(∂ǫ)
[ω(ǫ)]N−r
[ω(ǫ)− 1]N−1
We define the corresponding generating function
hN (z) :=
HN (r) z
r−1 . (3.4)
Noticing that ω(ǫ) → 0 as ǫ → 0, it can be shown that, given any arbitrary function
f(z) regular in a neighbourhood of the origin, the following inverse representation
holds
KN−1(∂ǫ)f(ω(ǫ))
(z − 1)N−1
hN(z)f(z) dz . (3.5)
Here C0 is a closed counterclockwise contour in the complex plane, enclosing the
origin, and no other singularity of the integrand.
3.6. Multiple integral representation. Plugging (3.5) into representation
(3.3), we readily obtain the following multiple integral representation for EFP:
FN (r, s) =
· · ·
dsω det
1≤j,k≤s
hN−k+1(ωj)
ωj − 1
ωN−r−1j
(ωj − 1)N
1≤j<k≤s
(ω̃j − 1)(ωk − 1)
ω̃jωk − 1
. (3.6)
Here ω̃j ’s should be expressed in terms of ωj ’s through (3.2). Indeed, due to (3.5),
relation (3.2) for functions ω(ǫ), ω̃(ǫ), translates directly into the same relation
between ωj and ω̃j , j = 1, . . . , s.
Representation (3.6), and all results in this Section hold for any choice of param-
eters λ and η within the disordered regime. Moreover, by analytical continuation
in parameters λ and η, these results can be easily extended to all other regimes.
The determinant in expression (3.6) is a particular representation of the parti-
tion function of the six-vertex model with DWBC, when the homogeneous limit is
performed only on a subset of the spectral parameters [CP3]. The structure of the
previous multiple integral representation therefore closely recalls analogous ones for
the Heisenberg XXZ quantum spin chain correlation functions [JM, KMT].
For generic values of λ and η, the orthogonal polynomialsKn(x), or the generat-
ing function hN (z), are known only in terms of rather implicit representations. For-
tunately, there are three notable exceptions [CP2]: the Free Fermion line (η = π/4,
−π/4 < λ < π/4, ∆ = 0), the Ice Point (η = π/6, λ = π/2, ∆ = 1/2), and the Dual
Ice Point (η = π/3, λ = π/2, ∆ = −1/2). In these three cases, the Kn(x) turn
out to be classical orthogonal polynomials, namely Meixner-Pollaczek, Continuous
Hahn and Continuous Dual Hahn polynomial, respectively. Correspondingly, the
8 F. COLOMO AND A.G. PRONKO
generating function can be represented explicitly in terms of terminating hyperge-
ometric functions that may simplify considerably further evaluation of EFP. In the
next Section we shall focus on the case of Free Fermion line.
4. Multiple integral representation at ∆ = 0
4.1. Specialization to η = π/4. We shall now restrict ourselves to the case
η = π/4. We have ∆ = 0, and the six-vertex model reduces to a model of free
fermions on the lattice. The parameter λ can still assume any value in the interval
(−π/4, π/4). It is convenient to trade λ for the new parameter
τ = tan2(λ− π/4) , 0 < τ < ∞ .
The symmetric point (related to the domino tiling of Aztec Diamond) corresponds
now to τ = 1. For generic values of τ we have:
ω̃ = −τω .
The generating function (3.4) is known explicitely (see [CP2] for details):
hN (z) =
1 + τz
1 + τ
Plugging this expression into (3.6), we get
FN (r, s) =
· · ·
dsω det
1≤j,k≤s
(1 + τωj)(ωj − 1)
(1 + τ)ωj
ωN−r−1j
(ωj − 1)N
1≤j<k≤s
(1 + τωj)(ωk − 1)
1 + τωjωk
. (4.1)
4.2. Symmetrization. After extracting a common factor
(1 + τωj)(ωj − 1)
(1 + τ)ωj
from the determinant in (4.1), we recognize it to be of Vandermonde type. We can
therefore collect from the integrand of (4.1) the double product
1≤j<k≤s
(1 + τωj)(ωj − 1)
(1 + τ)ωj
(1 + τωk)(ωk − 1)
(1 + τ)ωk
(1 + τωj)(ωk − 1)
1 + τωjωk
Noticing that the integration and the remaining of integrand are fully symmetric
under permutation of variables ω1, . . . , ωj , we can perform total symmetrization of
the previous double product over all its variables, with the result
(−1)s(s−1)/2
ωs−1j
1≤j<k≤s
(ωj − ωk)
THE ARCTIC CIRCLE REVISITED 9
Hence, we finally obtain the following representation for EFP on the Free Fermion
line:
FN (r, s) =
(−1)s(s+1)/2
s!(1 + τ)s(N−s)(2πi)s
· · ·
1≤j<k≤s
(ωj − ωk)
(1 + τωj)
(ωj − 1)s ω
. (4.2)
The appearance of a squared Vandermonde determinant in this expression naturally
recalls the partition functions of s× s Random Matrix Models.
5. Triple Penner model and Arctic Ellipses
5.1. Scaling limit. We shall now address the asymptotic behaviour of expres-
sion (4.2) for EFP in the ∆ = 0 case. We are interested in the limit N, r, s → ∞,
while keeping the ratios
r/N = x , s/N = y ,
fixed. In this limit, x, y ∈ [0, 1] will parameterize the unit square to which the
lattice is rescaled. Correspondingly EFP is expected to approach a limit function
F (x, y) := lim
FN (xN, yN) , x, y ∈ [0, 1] .
We shall exploit the standard approach developed for instance in the investigation
of asymptotic behaviour for Random Matrix Models. Before this let us however
point out some facts which holds already for any finite value of s.
5.2. A useful identity. Consider the quantity
IN (r, s) :=
(−1)s(s+1)/2
s!(1 + τ)s(N−s)(2πi)s
· · ·
1≤j<k≤s
(ωj − ωk)
(1 + τωj)
(ωj − 1)s ω
which differs from (4.2) only in the integration contours. Here C1 is a closed,
clockwise oriented contour (note the change in orientation) in the complex plane
enclosing point ω = 1, and no other singularity of the integrand. We have the
identity
IN (r, s) = 1 (5.1)
for any integer r, s = 1, . . . , N . The simplest way to prove the previous identity is
by shifting ωj → ωj + 1, and rewriting IN (r, s) as an Hankel determinant; indeed
we have
IN (r, s) =
(−1)s(s−1)/2
(1 + τ)s(N−s)
1≤j,k≤s
ωj+k−2−s(1 + τ + τω)N−s
(1 + ω)r
The entries of the Hankel matrix vanish for j+k > s+1, and hence the determinant
is simply given by the product of the antidiagonal entries, j + k = s + 1 (modulo
a sign (−1)s(s−1)/2 emerging from the permutation of all columns). Identity (5.1)
follows immediately.
10 F. COLOMO AND A.G. PRONKO
5.3. Saddle-point evaluation for large N and finite s. When using the
saddle-point method in variables ω1, . . . , ωs to evaluate the behaviour of FN (r, s) for
large N and r, and finite s, it is rather easy to see that the saddle-point equations
decouple at leading order, and that each saddle-point will be on the real axis,
contributing with a factor e−NSj with Sj positive.
If a given saddle-point is smaller than 1, the contour C0 can be deformed
through the saddle-point without encountering any pole, and its contribution will
vanish as e−NSj in the large N limit. If however the saddle-point, still on the real
axis, happens to be larger than 1, the deformation of the contour C0 through the
saddle-point will pick up the contribution of the pole at ω = 1 (with a reversed
orientation of the contour), and the j-th integral will behave as 1 + e−NSj . Hence,
in the large N limit (at fixed s) the quantity FN (r, s) will vanish unless all the
saddle-points are greater than 1, in which case FN (r, s) ∼ IN (r, s) = 1. Note that
in the present situation the s saddle-points coincide. A detailed analysis shows that
in this case the position of the s saddle-points depends on the value x = r/N as
τ(1−x)
. In correspondence to the value x0 =
, for which these saddle-
points are exactly 1, the function F (x, 0) has a step discontinuity. More precisely,
it is easy to show that for x ∈ [0, 1], F (x, 0) = θ(x − x0), where θ(x) is Heaviside
step function. From a physical point of view x0 is the contact point between the
limit shape and the boundary. What have been discussed here can easily be verified
in the case s = 1. The extension to finite s > 1 is rather direct as well.
5.4. Saddle-point equation. Having in mind the analogy with s × s Ran-
dom Matrix Models, and the scaling limit specified in Section 5.1, we rewrite our
expression for FN (r, s) at ∆ = 0 as follows:
FN (r, s) =
(−1)s(s+1)/2
s!(1 + τ)s
2(1/y−1)(2πi)s
· · ·
dsω exp
j,k=1
j 6=k
ln |ωj − ωk|
ln(τωj + 1)− ln(ωj − 1)−
ln(ωj)
. (5.2)
Both sums in the exponent are O(s2). The corresponding (coupled) saddle-point
equations read
ωj − 1
(1/y − 1)τ
τωj + 1
k 6=j
ωj − ωk
. (5.3)
A standard physical picture reinterprets the saddle-point equations as the equi-
librium condition for the positions of s charged particle confined to the real axis,
with logarithmic electrostatic repulsion, in an external potential. In the present
case the latter can be seen as generated by three external charges, 1, x/y, and
−(1/y − 1) at positions 1, 0, and −1/τ , respectively. It is natural to refer to this
model as the triple Penner model. Although the simple Penner [P] matrix model
has been widely investigated, not so much is known about the much more compli-
cate double Penner model [M, PW]. We have not been able to trace any previous
study concerning the triple Penner model.
THE ARCTIC CIRCLE REVISITED 11
5.5. The exact Green function at finite s. To investigate the structure of
solutions of the saddle-point equations (5.3) for large s we first introduce the Green
function
Gs(z) =
z − ωj
which, if the ωj ’s solves (5.3), has to satisfy the differential equation:
z(z − 1)(τz + 1)
sG′s(z) + s
2G2s(z)
− s(αz2 + βz + γ)sGs(z)
τs(s− 1)− αs2
z + (1 − τ)s(s− 1)− βs2 +Ω
2τs(s− 1)− αs2
. (5.4)
The coefficients α, β and γ are readily obtained as the coefficients of the second
order polynomial appearing in the numerator, when setting to common denominator
the left hand side of (5.3). We give them explicitly for later convenience:
α = τ
, β =
+ (1− τ)
, γ = −
The derivation of the differential equation is very standard (see, e.g., [SD]). The
left hand side is built by suitably combining the explicit definition of the Green
function and its derivative. The result has to be a polynomial of the first degree in
z, whose coefficients are constructed by matching the leading and first subleading
behaviour of the left hand side as |z| → ∞.
5.6. The first moment Ω. The quantity Ω appearing in (5.4) is defined as
the first moment of the solutions of the saddle-point equations:
It is related in a obvious way to the first subleading coefficient of Gs(z); indeed,
from the definition of the Green function, it is evident that
Gs(z) =
+O(z−3) , |z| → ∞ .
It is worth to emphasize that Ω is still unknown, and that in principle its value
should be determined self consistently by first working out the explicit solution
of Gs(z) (which will depend implicitly on Ω), from (5.4) and then demanding that
j=1 ωj evaluated from this solution coincides with Ω. The appearance of the un-
determined parameter Ω is a manifestation of the ‘two-cuts’ nature of the Random
Matrix Model related to (5.2), see, e.g., par. 6.7 of [D1].
5.7. The asymptotic Green function. We are now in condition to perform
the large s (and large N , r) limit at fixed x, y. In the limit, we can neglect terms of
order O(s) in the differential equation (5.4), which therefore reduces to an algebraic
equation for the limiting Green function G(z):
z(z−1)(τz+1)[G(z)]2−(αz2+βz+γ)G(z) = (τ−α)z+(1−τ−β)+Ω(2τ−α) . (5.5)
The previous algebraic equation has to be supplemented by the normalization con-
dition
G(z) ∼
, |z| → ∞ . (5.6)
12 F. COLOMO AND A.G. PRONKO
Hence the Green function describing the large s asymptotic distribution of solutions
for the saddle equation (5.3) reads:
G(z) =
2z(z − 1)(τz + 1)
(αz2 + βz + γ)
(αz2 + βz + γ)2 + 4z(z − 1)(τz + 1)[(τ − α)z +Ω(2τ − α) + 1− τ − β]
(5.7)
We have selected the positive branch of the square root, to satisfy the normalization
condition (note that the coefficient of z4 under the square root is (α − 2τ)2, and
α − 2τ is negative for any x, y ∈ [0, 1]). However, the expression for G(z) is not
completely specified yet, because Ω is still undetermined.
5.8. Limit shape and condensation of roots. The polynomial under the
square root is of fourth order, hence G(z) will have in general two cuts in the
complex plane. The emergence of a two-cut problem was already expected from
the appearance of the undetermined first moment Ω in (5.4). The discontinuity of
G(z) across these cuts defines, when positive, the density of solutions of the saddle-
point equations (5.3) when s → ∞. The problem of explicitly finding this density,
for arbitrary α, β, γ (or x, y), is a formidable one, not to mention the evaluation of
the corresponding ‘free energy’, and of the saddle-point contribution to the integral
in (5.2). But our aim is much more modest, since we are presently interested only
in the expression of the limit shape, i.e. in the curve in the square x, y ∈ [0, 1],
delimiting regions where F (x, y) = 0 from regions where F (x, y) = 1. Of course we
are here somehow assuming that the transition of F (x, y) from 0 to 1 is stepwise in
the scaling limit, but this is supported both by the physical interpretation of EFP
(in the disordered region, by definition, the number of ‘thin’ lines is macroscopic,
and the probability of finding no ‘thin’ horizontal edges immediately vanishes in
the scaling limit) and by the discussion of Section 5.3.
As explained in the discussion of the double Penner model in paper [PW], the
logarithmic wells in the potential can behave as condensation germs for the saddle-
point solutions. In our case, this can role can be played only by the ‘charge’ at
ω = 1 in the Penner potential since the charge at ω = −1/τ is always repulsive,
while the one at ω = 0 is larger than 1, at least in the region of interest. [PW]
have shown that condensation can occur only for charges less than or equal to 1,
since this will be the fraction of condensed solutions. This consideration, together
with the expected stepwise behaviour and the discussion in Section 5.3, suggest
the following picture for the evolution of saddle-point solution density from the
disordered region, F (x, y) ∼ 0, to the upper left frozen region, F (x, y) ∼ 1: in the
disordered region there is a macroscopic fraction of solutions which are real and
smaller than 1, while in the upper left frozen region this fraction vanish. On the
basis of the discussion here and in Sections 3.2 and 5.3, we shall assume that at
the transition between the two regions all saddle-point solutions have condensed at
ω = 1.
5.9. Main assumption. We claim that the Arctic curve in the square x, y ∈
[0, 1] separating the disordered phase from the upper left frozen phase is defined by
the condition that all solutions of the saddle-point equation lies at ω = 1.
In the derivation of the limit shape, this is indeed the only assumption to which
we are unable to provide a proof. There is in fact no guarantee, at this level, for
THE ARCTIC CIRCLE REVISITED 13
this possibility to occur, and limit shapes could in principle emerge from a different
condition. But if for some values of x, y ∈ [0, 1] we have all solutions of the saddle-
point equation condensing at ω = 1, then this provides a transition mechanism
from 0 to 1 for F (x, y), and this might correspondingly define some limit shape.
If all saddle-point solutions condensate at ω = 1, then we obviously have:
Ω = 1 .
Moreover, the complicate expression (5.7) for G(z) should simply reduce to
G(z) =
z − 1
, (5.8)
since we expect to have no cuts, and only one pole at z = 1 with unit residue.
5.10. Arctic Ellipses. Consider the quartic polynomial under the square root
in (5.7). It is convenient to rewrite it in terms of
α̃ := 2τ − α = τ
β̃ := 2− β = τ
x+ y − 1
y − x
γ̃ := −γ =
(5.9)
Note that α̃ and γ̃ are always positive for x, y ∈ [0, 1]. When Ω = 1, our quartic
polynomial reads
α̃2z4 + 2α̃β̃z3 + (β̃2 + 2α̃γ̃)z2 + 2β̃γ̃z + γ̃2 ,
which may be equivalently rewritten as
(α̃z2 + β̃z + γ̃)2 .
We see that the quartic polynomial reduces to a perfect square, and hence, when
Ω = 1, the two cuts of G(z) disappear, as expected.
Now, when Ω = 1, in our new notations, the Green function reads:
G(z) =
[(2τ − α̃)z2 + (2− β̃)z − γ̃] +
(α̃z2 + β̃z + γ̃)2
2z(τz + 1)(z − 1)
. (5.10)
We now require the coefficients α̃, β̃, γ̃ to be such that the polynomial under
the square root combines with the first part of the numerator in (5.10) to give
2z(τz+1) and simplify the Green function according to (5.8). Once we have chosen
a given branch of the square root (the positive one, in order to satisfy normalization
condition (5.6)), it is obvious that the required simplification can occur for any z
in the complex plane only if the second order polynomial α̃z2 + β̃z + γ̃ does not
change its sign, i.e. only if its two roots coincide, implying:
β̃2 − 4α̃γ̃ = 0 .
Rewriting the last relation in terms of x, y, through (5.9), we readily get
(1 + τ)2x2 + (1 + τ)2y2 − 2(1− τ2)xy − 2τ(1 + τ)x − 2τ(1 + τ)y + τ2 = 0 .
We have therefore recovered the limit shape, which in this Free Fermion case is the
well-known Arctic Ellipse (Arctic Circle for τ = 1) [JPS, CEP]. We recall that,
as discussed in Section 3.2, F (x, y) is non-vanishing only in the upper left region
14 F. COLOMO AND A.G. PRONKO
of the unit square. Therefore, concerning EFP, only the upper left portion of the
Arctic curve, between the two contact points at ( τ
, 0) and (1, 1
), is relevant.
6. Concluding remarks
Our starting point has been the definition of a relatively simple but relevant
correlation function for the six-vertex model with DWBC, the Emptiness Formation
Probability. We have provided both a determinant representation and a multiple
integral representation for the proposed correlation function. This is the first ex-
ample in literature of a bulk (as opposed to boundary) correlation function for the
considered model, for generic weights.
The multiple integral representation, specialized to the Free Fermion case, has
been studied in the scaling limit. In the standard picture of Random Matrix Mod-
els, we recognize the emergence of a triple Penner model. Assuming condensation of
the roots of saddle point equations in correspondence to a limit shape, we recover
the well-known Arctic Circle and Ellipse. It would be interesting to investigate
whether universality considerations of Random Matrix Models (see, e.g., [D2]) can
be extended to the Penner model in the neighbourhood of its logarithmic singular-
ities. This would imply directly the results of [CEP, J1, J2] on the Tracy-Widom
distribution and the Airy process, emerging in a suitably rescaled neighbourhood
of the Arctic Ellipse.
It is worth to stress that the multiple integral representation for EFP presented
in Section 3 can be studied beyond the usual Free Fermion situation. We expect
that condensation of roots of the saddle point equation in correspondence of the
limit shape is a general phenomenon. We believe that this assumption could be of
importance in addressing the problem of limit shapes in the six-vertex model with
DWBC.
Our derivation of the limit shape in the Free Fermion case uses the explicit
knowledge of function hN (z), standing in the multiple integral representation (3.6).
It is worth mentioning that function hN(z) is also known explicitly at Ice Point,
(∆ = 1/2), and Dual Ice Point, (∆ = −1/2), being expressible in terms of (poly-
nomial) Gauss hypergeometric function [Ze, CP2]. For instance, at Ice Point the
triple Penner model discussed above generalizes to a two-matrix Penner model.
This model can be studied along the lines presented here, thus providing a solution
to the longstanding problem of limit shape for Alternating Sign Matrices.
Acknowledgements
We thank Nicolai Reshetikhin for useful discussion, and for giving us a draft
of [PR] before completion. FC is grateful to Percy Deift, and Courant Institute
of Mathematical Science, for warm hospitality. AGP thanks INFN, Sezione di
Firenze, where part of this work was done. We acknowledge financial support from
MIUR PRIN program (SINTESI 2004). One of us (AGP) is also supported in
part by Civilian Research and Development Foundation (grant RUM1-2622-ST-
04), by Russian Foundation for Basic Research (grant 04-01-00825), and by the
program Mathematical Methods in Nonlinear Dynamics of Russian Academy of
Sciences. This work is partially done within the European Community network
EUCLID (HPRN-CT-2002-00325), and the European Science Foundation program
INSTANS.
THE ARCTIC CIRCLE REVISITED 15
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I.N.F.N., Sezione di Firenze and Dipartimento di Fisica, Università di Firenze, Via
G. Sansone 1, 50019 Sesto Fiorentino (FI), Italy
E-mail address: [email protected]
Saint Petersburg Department of Steklov Mathematical Institute of Russian Acad-
emy of Sciences, Fontanka 27, 191023 Saint Petersburg, Russia
E-mail address: [email protected]
|
0704.0363 | Time and motion in physics: the Reciprocity Principle, relativistic
invariance of the lengths of rulers and time dilatation | arXiv:0704.0363v2 [physics.gen-ph] 10 Feb 2009
Time and motion in physics: the Reciprocity Principle,
relativistic invariance of the lengths of rulers and time
dilatation
J.H.Field
Département de Physique Nucléaire et Corpusculaire Université de Genève . 24, quai
Ernest-Ansermet CH-1211 Genève 4.
E-mail: [email protected]
Abstract
Ponderable objects moving in free space according to Newton’s First Law con-
stitute both rulers and clocks when one such object is viewed from the rest frame
of another. Together with the Reciprocity Principle this is used to demonstrate, in
both Galilean and special relativity, the invariance of the measured length of a ruler
in motion. The different times: ‘proper’, ‘improper’ and ‘apparent’ appearing in
different formulations of the relativistic time dilatation relation are discussed and
exemplified by experimental applications. A non-intuitive ‘length expansion’ effect
predicted by the Reciprocity Principle as a necessary consequence of time dilatation
is pointed out.
PACS 03.30.+p
http://arxiv.org/abs/0704.0363v2
1 Introduction
The standard text-book presentation of special relativity follows closely that of Ein-
stein’s seminal paper of 1905 [1] in basing the theory on the Special Relativity Principle,
classical electromagnetism and the postulate of constant light speed. However an alterna-
tive and conceptually simpler approach to the physics of space and time, in the absence
of gravitational fields, is possible in which it is not necessary to consider light signals,
classical electromagnetism, or indeed, any dynamical theory whatsoever. The Lorentz
transformation (LT) was first derived in this way by Ignatowsky [2] in 1910. Purely
mathematical considerations imply, in such a derivation of the LT, the existence of a
maximum relative velocity, V , of two inertial frames. Use of relativistic kinematics then
shows that V is equal to the speed of light, c, when light in identified as a manifestation
of the propgation in space-time of massless particles –photons [3]. In this way Einstein’s
mysterious second postulate is derived from first principles. The fundamental axiom un-
derlying such an approach is the Reciprocity Principle (RP) [4, 3], discussed in Section
3 below, relating the the relative velocities of two inertial frames. Derivations of the LT
and the parallel velocity addition formula based on the RP and other simple axioms are
given in Ref. [3].
In the present paper the space-time properties of ponderable1 physical bodies in free
space, as described by Newton’s First Law of mechanics, are used together with the RP,
to demonstrate the invariance of the measured length of a ruler in uniform motion. The
proof given is valid in both Galilean and special relativity, since Newton’s First Law and
the RP hold in both theories.
The analysis presented is based on a careful definition of physical time concepts. In
particular, the ‘frame time’ or ‘proper times’ that appear in in the RP, are distinguished
from the ‘improper time’ or ‘apparent time’ (of a moving clock) that appear in the Time
Dilatation (TD) relation of special relativity.
The paper is organised as follows: The following section contains an elementary discus-
sion of the concepts of ‘space’, ‘time’ and ‘motion’ in physics, in connection with Newton’s
First Law. In Section 3, the RP is introduced and discussed in relation to Newton’s First
Law. It is pointed out that, because of the RP, ‘rulers are clocks’ and ‘clocks are rulers’
when the motion of ponderable bodies in free space is considered. In Section 4 the RP is
used to demonstrate the invariance of the measured length of a uniformly moving ruler.
In Section 5 the operational meanings of the time symbols appearing in the TD formula
of special relativity are discussed. This may be done in a ‘clock oriented’ manner in terms
of ’proper’ and ’improper’ times of the observed clock, or in an ‘observer oriented’ manner
in terms of the proper time of the observer’s local clock and the ‘apparent time’, as seen
by the observer, of the moving clock. Two specific experiments are described to exemplify
the operational meanings of the time symbols in the TD formula. A non-intuitive ‘length
expansion’ effect is found to relate similarly defined spatial intervals corresponding to the
observation of an event either in the rest frame of the clock, or in a frame in which it is
in uniform motion.
1That is bodies, with a non-vanishing Newtonian mass, which may be associated with an inertial
frame in which the body is at rest. No such frame may be associated with a massless object.
The results of the present paper show that the ‘length contraction’ effect and the
correlated ‘relativity of simultaneity’ effect of conventional special relativity do not exist.
A detailed discussion of the reason for the spurious nature of these effects of conventional
special relativity theory may be found in Refs. [5, 6, 7, 8, 9, 10] .
However, a genuine ‘relativistic length contraction’ effect does occur when distances be-
tween spatial coincidences of moving objects are observed from different inertial frames [11].
Also a genuine ‘relativity of simultaneity’ effect occurs when clocks at rest in two different
inertial frames are viewed from a third one [12, 13]. An alternative derivation, directly
from the RP, of the invariance of the measured spatial separation of two objects at rest
in the same inertial frame as well as the absence of the conventional ‘relativity of simul-
taneity’ effect is given in Ref. [9].
2 Physical time and Newton’s First Law of Mechan-
In physics the concepts of ‘time’ and ‘motion’ are inseparable. In a world in which
motion did not exist the physical concept of time would be meaningless. Similarly the
physical concepts of ‘space’ and ‘motion’ are inseperable. Without the concept of space,
no operational definition of motion is possible. The concept of historical time –the time
of the everyday world of human existence– requires the introduction of the further, and
equivalent, concepts of ‘uniform motion’ and ‘cyclic motion with constant period’. For
example, the unit of time the ‘year’ is identified with the (assumed constant) period of
rotation of the Earth about the Sun.
The idea of uniform motion entered into physics in a quantitative way with the for-
mulation of Newton’s First Law [14]
Every body continues in its state of rest, or uniform motion in a right
line unless it is compelled to change that state by forces impressed upon
This law gives an operational meaning to the physical concept of ‘uniform motion’.
It is defined by observations of the position of any ponderable object in ‘free space’ i.e.
in the absence of any mechanical interaction of the object with other objects. There is
a one-to-one correspondence between such a ponderable object and an ‘inertial frame’
of relativity theory. As will be discussed in the following section, one such ponderable
object, O, constitutes both a ruler and a clock for an observer in the rest frame of another
such object, O’, and vice versa.
When time is measured by using a cyclic physical phenomenon, e.g. an analogue clock,
time measurement reduces to recording the result of a spatial (or angular) measurement.
There is a one-to-one correspondence between the spatial coincidence of a stationary
‘mark’ on the face of the clock and a moving ‘pointer’, constituted by the hand of the clock,
and the time measurement [6]. A ‘time interval’ is measured by the angular separation of
two such ‘pointer-mark coincidences’. The implicit assumption is that the motion of the
pointer is ‘uniform’. There is an evident logical circularity here since ‘equal’ time intervals
measured by such an analogue clock assume that the angular velocity of the hand is
constant, whereas constant angular velocity is established by observation of equal angular
increments for equal time time intervals (i.e. also equal angular increments) recorded by
a second clock of supposedly known uniform rate. In practice, this conundrum is resolved
by an appeal to physics. For example, an undamped pendulum in a uniform gravitational
field is predicted, by the laws of mechanics, to have a constant period of oscillation.
Quantum mechanics predicts the same transition frequency and mean lifetimes for two
identical atoms in the same excited state, in the same physical environment, etc.
Measurements of ‘time’ are then ultimately observations of spatial phenomena, e.g.
the time measurement corresponding to observation of the number displayed by a digital
clock is a spatial perception. This will also be the case for time measurements related to
observation of two ponderable objects O and O’ in motion in free space that will now be
discussed.
3 The Reciprocity Principle: rulers are clocks, and
clocks are rulers
Consider two non-interacting ponderable objects O and O’, with arbitary motions in
free space. They are placed at the origins of inertial coordinate systems S and S’ with
axes orientated so that the x and x′ axes are parallel to the relative velocity vector of O
and O’. Without any loss of generality for the following discussion, it may be assumed
that O and O’ lie on the common x-x′ axis.
The Reciprocity Principle (RP) [4, 3, 9] is defined by the equation:
v = vO′O =
∂xO′O
= −vOO′ (3.1)
where xO′O ≡ xO′−xO and x
, or in words: ‘If the velocity of O’ relative to O
is ~v, the velocity of O relative to O’ is - ~v’. In many discussions of special relativity, the RP
is taken as ‘obvious’ and is often not even declared as a separate axiom. This is the case,
for example, in Einstein’s 1905 special relativity paper [1]. However, as first demonstrated
by Ignatowsky in 1910 [2], it is sufficient, together with some other weaker axioms such as
the homogeneity of space or single-valuedness of the transfomation equations, to derive [3]
the space-time Lorentz transformation and hence the whole of special relativity theory.
Eqn(3.1) looks very similar to the equation defining the relative velocity of two objects
A and B as observed in a single inertial reference frame (say S):
vAB ≡ vA − vB =
d(xA − xB)
= −vBA (3.2)
The crucial difference is the appearence in the RP, (3.1), of two different times t and t′.
The time t is the ‘frame time’ of S. i.e. the time registered by a synchronised clock at rest,
at any position in S, according to an observer also a rest in S. The frame time t′ is similarly
defined by an array of synchronised clocks at rest in S’. Eqn(3.1) (and its integral) gives a
relation between the times t and t′ Both t and t′ correspond to ‘proper times’ of clocks at
rest, whereas, as explained in Section 4 below, the Lorentz transformation relates instead
a proper time to an ‘improper time’ –the observed time of a clock in uniform motion.
Suppose now that O and O’ are equipped with local clocks that are observed to run at
exactly the same rate when they are both at rest in the same inertial frame. The direction
of the relative velocity vector ~v of O’ relative to O is such that they are approaching each
other at the frame times t and t′. The spatial separations of O and O’ in S and S’ are ℓ(t)
and ℓ′(t′) respectively, at times t and t′. Using the RP, a spatial coincidence of O and O’
will be observed at the time
tOO′ = t+
(3.3)
in S, and
= t′ +
ℓ′(t′)
(3.4)
in S’. The OO’ coincidence event will be mutually simultaneous in the frames S and S’.
Note that the OO’ spatial coincidence that is mutually simultaneous in S and S’
constitutes a pair of reciprocal pointer mark coincidences. In S the mark is at the position
of O and the moving pointer at the position of O’, whereas in S’ the position of O’
constitutes the mark and the position of O the pointer. A corollary is that all such pairs
of reciprocal pointer mark coincidences are mutually simultaneous. This is the basis of the
‘system external synchronistation’ [15] as introduced in Einstein’s first special relativity
paper [1] to synchronise clocks at rest in different inertial frames when they are in spatial
coincidence.
The observation of the OO’ coincidence event in both frames can be used to give a
condition that any other pair of events, one observed in S, the other observed in S’ are
mutually simultaneous. If the time of an event in S is t̃ and another event in S’ is t̃′ they
will be ‘mutually simultaneous’ providing that:
t̃′ − t̃ = t′
− tOO′ (3.5)
Combining (3.3)-(3.5) gives:
t̃′ − t̃ = t′
− tOO′ = t
′ − t+
ℓ′(t′)− ℓ(t)
(3.6)
If now events occuring at times t in S and t′ in S’ are mutually simultaneous, it follows
from (3.5) and (3.6) that ℓ(t) = ℓ′(t′), so that events which occur when O and O’ have
the same spatial separation in S and S’ are mutually simultaneous. A special case occurs
if the clock arrays in S and S’ are mutually synchronised so that ℓ(t) = ℓ′(t′ = t). There
is then a direct correlation between either t or t′ and the spatial separation of O and O’:
When mutually synchronised clocks in the frames S and S’ have the same reading, O and
O’ have the equal spatial separations in S and S’, and conversely, When O and O’ have
equal spatial separations in the frames S and S’, mutually synchronised clocks in S and S’
have the same reading.
The dependence of ℓ on t in Eqn(3.3) and ℓ′ on t′ in Eqn(3.4) means that each of
the objects may be considered to be an ‘inertial clock’ by an observer in the rest frame
of the other one. That is, t is measured by the spatial separation of O’ from O in S
and t′ is measured by the spatial separation of O from O’ in S’. Conversely, after mutual
synchronisation of the clock arrays in S and S’ at the instant when O and O’ are in spatial
coincidence, t measures the spatial separation of O’ and O in S (and so is effectively a
ruler in this frame) while t′ measures the spatial separation of O’ and O in S’, constituting
a ruler in this frame. Matching of these measurements of the separation of O and O’ with
the lengths of physical rulers at rest in S and S’ is now used to demonstrate the invariance
of the measured length of the length of a ruler in uniform motion –that is, the absence of
any relativistic length contraction effect– in this case.
4 Invariance of the measured length of a ruler in uni-
form motion
Figure 1: Rulers attached to objects O and O’ are viewed from the frame S (left) and S’
(right). The equality of the separations of O and O’ in S and S’ at time t = t′ = L/v,
predicted by the RP, is used to establish the invariance of the measured length of the
moving ruler R’ in S, or of the moving ruler R in S’ (see text).
Suppose that O and O’ are equipped with rulers R and R’, parallel to the x-x′ axis
as shown in Fig.1. O coincides with the mark MR(0) of the ruler R and O’ with the
mark MR′(10) of the ruler R’. A t = t′ = 0 (Fig.1a) O and O’ are in spatial coincidence.
The clock arrays in S and S’ are mutually synchronised at this time. The length of each
ruler in its rest frame is L. The object O’ now moves along the ruler R, being in spatial
coincidence with different marks of the ruler at different times. The object O moves in a
similar manner along the ruler R’. At any given time t the separation of O and O’ in S is
given by the corresponding ‘Pointer Mark Coincidence’ (PMC):
PMC(O′, t) ≡ O′(t)@MR(J) (4.1)
where the symbol before the ampersand denotes the moving ‘pointer’, and the symbol
after it the stationary ‘mark’ with which it is spatial coincidence2. Since
PMC(O, t) ≡ O(t)@MR(0) for all t (4.2)
and x[MR(0)] = 0 it follows that the separation of O and O’ in the frame S at time t is
given by:
dO′O(t) = x[MR(J)]− x[MR(0)] = x[MR(J)] (4.3)
where
x[MR(J)] =
and where, in Fig.1, Jmax = 10, is the ordinal number of the mark at the end of the ruler.
Thus the x-coordinate origin is at MR(0). Defining in a similar manner a PMC in the
frame S’:
PMC(O, t′) ≡ O(t′)@MR′(K) (4.4)
and since
PMC(O′, t′) ≡ O′(t′)@MR′(10) for all t′ (4.5)
the separation of O and O’ in S’ at the time t′ is
(t′) = x′[MR′(10)]− x′[MR′(K)] (4.6)
where
x′[MR′(K)] =
and where, in Fig.1, Kmax = 10. The spatial configurations in S and S’ at the times
t = t′ = L/v are shown in Fig.1b. The corresponding PMC are:
PMC(O′, L/v) ≡ O′(L/v)@MR(10) (4.7)
PMC(O, L/v) ≡ O(L/v)@MR′(0) (4.8)
It follows from (4.3) and (4.6) that
dO′O(L/v) = x[MR(10)]−x[MR(0)] = L = x
′[MR′(10)]−x′[MR′(0)] = d′
(L/v) (4.9)
Since O’ coincides with MR′(10) at all times it follows that, at t = L/v
x[MR′(10)] = x[O′] = x[MR(10)] (4.10)
Also, since O is in spatial coincidence with MR′(0) at t = t′ = L/v it follows that at
t = L/v,
x[MR′(0)] = x[O] = x[MR(0)] = 0 (4.11)
2This notation was introduced in Ref. [6]. Note the similarity with an e-mail address
Eqns(4.9)-(4.11) then give at t = L/v:
x[MR′(10)]− x[MR′(0)] = x[MR(10)]− x[MR(0)] = L (4.12)
That is, the measured length of the moving ruler R’ in the frame S, at t = L/v, is the
same as the length of the same ruler at rest –there is no ‘length contraction’ effect. A
similar calculation for the length of the ruler R as measured in the frame S’ gives, at
t′ = L/v:
x′[MR(10)]− x′[MR(0)] = x′[MR′(10)]− x′[MR′(0)] = L (4.13)
The length of the moving ruler R as measured in S’, at t′ = L/v, is the same as the length
of the same ruler at rest. The above calculations have used the equality of the spatial
separations of O and O’ in S and S’ at equal times of mutually synchronised clocks in these
frames, that follows from the RP, to establish, via corresponding PMCs, the equality of
the measured lengths of a ruler at rest, or in motion. Note that nowhere in any of the
calculations was the Lorentz transformation invoked. In fact the calculations are the same
in Galilean and special relativity, since the RP is equally valid for both.
5 The time dilatation effect; proper, improper and
apparent time intervals
All the times considered above were ‘frame times’ i.e. t and t′ are the times recorded
by a synchronised clock at rest at any position in S and S’ as viewed by an observer at rest
in these respective frames. In order to discuss the time dilatation effect it will be found
convenient to use the notation t(S), t′(S ′) for the frame times where the arguments S, S’
specify the reference frame of the observer of the clock. Such times are proper times of
such a clock. The Lorentz transformation relates the space-time coordinates (x′,t′(S ′)) of
an event specified in the frame S’ to those of the same event, (x,t′(S)) as observed in S, or
vice versa. The times t(S ′)[ t′(S)] which are those of clocks at rest in S[S’], as viewed from
S’[S] are called improper times. The space-time LT gives the following invariant interval
relation between corresponding space and time intervals in the frames S and S’:
c2(∆τ ′)2 = c2(∆t′(S))2 − (∆x)2 = c2(∆t′(S ′))2 − (∆x′)2 (5.1)
where ∆x ≡ x2 − x1 etc, while the inverse LT gives:
c2(∆τ)2 = c2(∆t(S ′))2 − (∆x′)2 = c2(∆t(S))2 − (∆x)2 (5.2)
In order to use the general interval relation (5.1) to derive the time dilatation effect it is
necessary to identify the time interval ∆t′(S ′) with the proper time interval of a clock at
rest in S’ (∆x′ = 0), and with equation of motion in S: ∆x = v∆t′(S). Using the latter
equation to eliminate ∆x from (5.1) and setting ∆x′ = 0 yields the time dilatation (TD)
relation:
∆t′(S) = γ∆t′(S ′) (5.3)
Figure 2: An experiment to illustrate the TD effcet viewed from S (left) and S’ (right). a)
The pulsed lamp PL at rest in S flashes at time t(S) = L/v and PL’ at rest in S’ flashes
at time t′(S ′) = L/v. b) The light signal from PL is observed at time t(S ′) = γL/v in the
frame S’, that from PL’ at time t′(S) = γL/v in the frame S. The PMCs corresponding to
the positions of observation of the signals in the different frames are indicated. See text
for discussion.
Figure 3: Spatial configurations in the frame S (left) and the frame S’ (right) are viewed
at different times. a) t(S) = t′(S ′) = 0; the Λ is created and moves to the right in the
plane of the figure with speed v =
3c/2. b) t(S) = t′(S ′) = T ′; the Λ is observed to decay
in the frame S’. The decay products move in the plane of the figure perpendicular to the
direction of motion of the Λ. c) t(S) = t′(S ′) = γT ′; the Λ is observed to decay in the
frame S. See text for discussion. The momentum vectors of the p and π− are drawn to
scale in the different reference frames. The spatial position of each particle is at the tail
of the corresponding momentum vector.
where γ ≡ 1/
1− (v/c)2, relating the improper to the proper time of a clock at rest in
S’. In a similar manner the interval relation (5.2) gives the TD relation for a clock at rest
in S and observed from S’:
∆t(S ′) = γ∆t(S) (5.4)
It is important to note the existence of four different time symbols, with different opera-
tional meanings in Eqns(5.3) and (5.4). The proper times t(S) and t′(S ′) (corresponding
to the ‘frame times’ t and t′ of the previous sections) and the improper times t(S ′) and
t(S ′). The notation for these times just introduced may be called ‘clock oriented’ since
only the readings of a single clock (observed either at rest, or in motion) appear in the TD
relations. In any actual experiment where the TD effect in measured, two clocks are nec-
essary, the observed moving clock, and another one at rest to measure the corresponding
time interval in the observer’s proper frame. If a clock at rest in S’ is observed from S as
in Eqn(5.3), the time interval ∆t′(S) is actually that, ∆τ , recorded by a similar clock, at
rest in S while ∆t′(S ′) is the corresponding time interval recorded by the (slowed-down)
moving clock. Since the observed rate of the moving clock depends on its motion, ∆t′(S ′)
is not a proper time interval for the observer in S. From the view-point of the latter this
is an ‘apparent’ (velocity-dependent) time interval that may be denoted simply as ∆t′,
to distinguish it from the observer’s proper time interval ∆τ . This gives an alternative
‘observer oriented’ time notation for the TD relations (5.3) and (5.4) above:
∆τ = γ∆t′ (5.5)
∆τ ′ = γ∆t (5.6)
This alternative notation has beeen employed in several previous papers by the present
author [6, 8, 11, 12, 13, 16].
In order to apply the TD relations (5.3) and (5.4), or (5.5) and (5.6), to any actual
or imagined experiment an operational definition must be given to the improper time
intervals of Eqns(5.3) and (5.4) or the apparent time intervals of (5.5) and (5.6). Two
examples of such definitions will be given, the first in a thought experiment to illustrate
the physical meaning of the TD effect, the second in an actual experiment typical of
many performed in particle physics, where the TD effect is used to measure the proper
decay time of an unstable particle. However as will be seen, the thought experiment and
actually realisable (and many times realised) one are similar in all essential features.
What notation is most convenient depends on the experiment considered. In the
observation of the TD effect in the last CERN muon g-2 experiment [17] where the time
interval ∆τ was directly measured by clocks in the laboratory frame, and ∆t′ was the
known muon rest-frame lifetime, it was natural to use Eqn(5.5). For the second of the
two experiments considered below where ∆τ is not directly measured but inferred from
spatial measurements in the frame S, the relation (5.3) relating connecting a proper time
in the frame S’ to an improper time in the frame S, is used.
In the thought experiment it is imagined that the objects O, O’ are each equipped
with local pulsed lamps PL, PL’. The objects O, O’ are in spatial coincidence at times
t(S) = t′(S ′) = 0 and are attached to rulers of length 2L in similar spatial configurations to
that shown in Fig.1a. The objects move apart with relative velocity v =
3c/2. As shown
in Fig.2a, at the times t(S) = t′(S ′) = L/v, PL and PL’ both flash, producing an isotropic
pulse of photons. The observation times in S of the photon signal produced by PL’, and
in S’ of the photon signal produced by PL, are given by Eqns(5.3) and (5.4) respectively.
Since γ = 2, these observations occur at the times t(S) = t′(S ′) = γL/v = 2L/v. The
corresponding spatial configurations of O and O’ at these times shown in Fig.2b. It can
be seen that the observation times of the light flashes in S and S’ correspond to different
PMCs of the objects O and O’ and to different spatial separations of the objects:
In S PL : PMC(MR′(10), L/v) ≡ MR′(10)@0 = MR′(10)@MR(0) (5.7)
PL′ : PMC(O′, γL/v) ≡ O′@MR(20) = MR′(20)@MR(20) (5.8)
In S′ PL′ : PMC(MR(10), L/v) ≡ MR(10)@O′ = MR(10)@MR′(20) (5.9)
PL : PMC(O, γL/v) ≡ O@MR′(0) = MR(0)@MR′(0) (5.10)
ℓ(γL/v)
ℓ(L/v)
ℓ′(γL/v)
ℓ′(L/v)
vγL/v
vL/v)
= γ (5.11)
The relations in (5.11) follow directly from the RP, while the PMC in (5.7)-(5.10) are
obtained from the geometry of Fig.2 and the invariance of the lengths of the moving
rulers derived in Section 3 above.
The different PMC corresponding to observations of the light flashes emitted by PL
and PL’ in different frames in (5.7)-(5.10) is deeply perplexing for common-sense con-
cepts of space and time. For example the photon bunches emitted by PL’ correspond
to MR(10)@MR′(20) in S’ and to MR′(20)@MR(20) in S. In some discussions of time
dilatation this apparent paradox is avoided by invoking a hypothetical contraction of a
moving ruler by a factor 1/γ [18]. This has the effect of shortening the moving ruler
R by a factor 1/2 in the right hand figure in Fig.2a, so that the PMC corresponding
to the flashing of PL’ becomes MR(20)@MR′(20), the same as in S with inversion of
pointer and mark. However, as demonstrated above, there is no such length contraction
effect, which, as pointed out elsewhere [5, 6, 7, 8, 9] is a spurious consequence of misin-
terpreting the space-time Lorentz transformation. Indeed the possibility of such a length
contraction effect is already excluded by inspection of Fig.2a. In the right hand figure, the
PMC correponding to the moving object O considered as a pointer is MR(0)@MR′(10).
Since O is in motion and R’ at rest no hypothetical length contraction effect operates
here. In the left hand figure the mutually simultaneous PMC in S is MR′(10)@MR(0)
so that at t(S) = t′(S ′) = L/v observers in S and S’ see reciprocal PMCs, i.e. ones
related by exchange of the pointer and mark symbols. If however the length contrac-
tion effect exists, the observer in S will see instead that the PMC corresponding to O is
MR′(0)@MR(0) at time t(S) = L/v. But from the RP this PMC must correspond to the
times t(S) = t′(S ′) = 2L/v (see Fig.2b) contrary to the assumption that t(S) = L/v. The
length contraction hypothesis therefore contradicts the corollory of the RP that states
that mutually simultaneous events in two frames have reciprocal PMCs, since it implies
that the reciprocal PMCs MR′(0)@MR(0) and MR(0)@MR′(0) are not mutually simul-
taneous.
The second example of a TD experiment illustrates a typical application of the effect
in particle physics (see Fig.3). A π− meson interacts with a proton in a thin plastic
target T to produce a Λ hyperon via the reaction3 π−p → ΛK0 The hyperon moves with
3The results of an actual such experiment constructed to test the ∆S = ∆Q rule in semileptonic
neutral kaon decays are described in Ref. [19].
velocity v =
3c/2 perpendicular to the plane of the target in the laboratory frame S.
After the time t′(S ′) = T ′ in its rest frame S’, it decays to a proton and a negative pion:
Λ → pπ−. These decay products are observed in the laboratory system. The experiment
is in every way similar to that shown in Fig.2. The object O is replaced by the target T,
the object O’ by the undecayed Λ or the kinematical system constructed from its decay
products. The photon pulse emitted by PL’ is replaced by the decay products of the
Λ. By reconstructing the trajectories of the decay p and π− in a particle detector the
position of the decay event and hence the decay length lD –the distance between the point
of production and decay of the Λ– in the frame S can be measured. Identification of the
p and π− and measurement of their momenta (typically by measurement of the curvature
of their trajectories in a known magnetic field ) enables the momentum P and the energy
E of the Λ to be determined. Since v = Pc2/E and γ = E/(mΛc
2) where mΛ is the mass
of the Λ, the proper decay time of the Λ is given by Eqn(5.3) as:
T ′ = ∆t′(S ′) =
∆t′(S)
(5.12)
The spatial configurations of T and the Λ at different times in the frames S and S’ are
shown in Fig.3. The spatial separations of T and the Λ at the observed instant of decay
in S and S’ obey the relation (5.11). This implies that this separation, in changing the
frame of observation from the rest frame of the Λ to the laboratory system in which it is
motion, undergoes a ‘length expansion’ by the factor γ. In accordance with Eqn(5.11), it
can be seen that this is a necessary consequence of the RP, given the existence of the TD
effect.
The mutally simultaneous events in S and S’ shown in Fig.3c, correspond, as they
must, to equal spatial separations of T and the physical object constituted by the decay
products, p and π−, of the Λ. However, in the frame S, these particles have just been
created and have vanishing spatial separation, whereas in S’ they are spatially separated by
a distance corresponding to a time-of-flight (γ− 1)T ′. This also seems highly paradoxical
when interpreted by commonsense classical concepts of space and time.
Acknowledgement
I thank the referee of the journal that rejected Ref. [11] for publication for correspon-
dence that was important for the clarificatiion of the ideas expressed in both the latest
version of Ref. [11] and the present paper.
Added Note
The calculations presented in the present paper are flawed by a major conceptual
misunderstanding which is rectified in later papers [20, 21] treating similar subjects.
At the time of writing the present paper, the author had correctly understood the
spurious nature of the ‘relativity of simultaneity’and ‘length contraction’ effects of con-
ventional special relativity theory [5, 7, 8, 10] but had not yet drawn the simple conclusion
that the existence of the genuine and experimentally-confirmed time dilatation effect then
necessarily implies that the Reciprocity Principle, as generally understood, also breaks
down in special relativity. This point is easily understood by considering the first member
of Eqn(3.1), written in a simplified notation as:
dxO′O
Transforming into the frame S’, the invariance of length intervals implies that
dxO′O = −dx
Since the time dilatation relation gives dt = γdt′, the Reciprocity Principle of (3.1) is
replaced by:
dxO′O
so that
= −γv
to be compared with v′ = −v given by (3.1).
The detailed calculations presented in Section 4 are correct and logically coherent
given the initial assumptions, but the configurations shown in the frame S’ in Fig.1 do not
correspond to observations in this frame of the coincidence events specified in the frame
S in the same space-time experiment. If this were the case, in the S’ frame configurations
in Fig.1 v should be replaced by γv and t and t′ should be related by time dilatation
relation t = γt′. In fact, what are shown in Fig.1 and considered in Section 4 are the
configurations in S of a primary experiment and in S’ of the corresponding but physically
independent reciprocal experiment [20, 22].
Nevertheless, the invariance of corresponding length intervals can be derived [21] by
considering the configurations in S and S’ in Fig.1b in the case that they are corresponding
ones, at the same epoch, in the same space-time experiment. In this case, as explained
above, the speed of O in S’ should be γv, not v. Consider, however, an object Õ with the
same x′ coordinate as O that does have the velocity v. The separation L′ of O and O’
in S’ is then equal to that between O’ and Õ. at the epoch of Fig.1b. Compare now the
configuration of O and O’ in S, with separation L with the corresponding one of Õ. and
O’ in S’ with separation L′. From the symmetry of the configurations it can be seen that
both L and L′ can depend only on v: L = L(v), L′ = L′(v). The reciprocity of the two
configurations is now invoked to give the condition, as stated by Pauli [23]:
The contraction of length at rest in S’ and observed from S is equal to
the length at rest in S as observed from S’.
The ‘length at rest in S’ ’ is L′ which ‘as observed from S’ is L, whereas the ‘length
at rest in S ’ is L which ‘as observed fron S’ ’ is L′. Denoting the contraction factor by
α(v), the above condition states that
L = α(v)L′, L′ = α(v)L
which implies that L = α(v)2L or α(v)2 = 1 so that L = L′ and the spatial separation
between O and O’ is the same in S and S’ at corresponding epochs. The same conclusion
is more simply reached by noting the symmetry of the configurations of O,O’ in S and
Õ,O’ in S’. and applying Leibnitz’ Principle of Sufficient Reason [21].
If, therefore, in the primary experiment, shown in S in Fig.2b and S’ in Fig.2a, the
configuration in S’ in Fig.2a is to correctly represent that corresponding to the configura-
tion in S in Fig.2b, the velocity v in S’ should be replaced by γv, so that when PL’ flashes
O’ is aligned with MR(20) in both S and S’. In the reciprocal experiment, shown in S in
Fig.2a and S’ in Fig.2b, v in S in Fig.2a should be replaced by γv so that O is aligned
with MR’(0) in both S and S’ when PL flashes.
Similarly, in the thought experiment of Fig.5, if the S’ frame configurations on the right
side of the figure are to represent observations in this frame of events shown in S by the
configurations on the left side, instead of what are actually shown which are configurations
of the physically independent reciprocal experiment, v should be replaced by γv in all the
S’ frame configurations. In this case, there is no mismatch between the spatial position of
the decay event in the two frames and the claimed ‘length expansion’ effect does not occur.
Indeed the claimed ‘... different PMC corresponding to observations of the light flashes
emitted by PL and PL’ in different frames in (5.7)-(5.10)’ is not only ‘...deeply perplex-
ing for common-sense concepts of space and time.’ it is the absurd (self-contradictory)
consequence of assuming, at the same time, that length intervals are invariant, time di-
latation occurs and the conventional interpretation of the Reciprocity Principle holds. In
conventional special relativity theory time dilatation and the Reciprocity Principle are
reconciled by invoking the spurious ‘length contraction’ effect dxO′O = −γx
[18]. so
that v′ = −v. The correct physical interpretation of the Reciprocity Principle is actually
the definition of the configuration in S’ of the physically-independent experiment that is
reciprocal to the primary one specified by the standard configuration of the frames S and
S’ [20, 22].
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|
0704.0364 | B --> rho K* decays and other rare vector-vector modes | B → ρK∗ decays and other rare vector-vector modes ∗
G. Vasseur†
DSM/DAPNIA/SPP, CEA/Saclay, F-91191 Gif-sur-Yvette, France
The recent analyses of the following rare vector-vector decays of the B meson are presented:
∗, ωK∗, ωρ, ωω, and ωφ charmless final states. The latest results indicate that the fraction of
longitudinal polarization is about 0.5 in penguin-dominated modes and close to 1 for tree-dominated
modes.
I. MOTIVATION
The search for rare charmless hadronic decays of the
B meson to vector-vector final states has become a quite
active field in the experiments at the B factories, Belle
at KEK and BABAR at SLAC. As a lot of these decays
have not yet been seen, the first goal of these studies
is to observe such modes and measure their branching
fraction. The measurements can then be compared to
theoretical predictions.
The direct CP-violation asymmetry in these modes
can also be measured. It is defined as ACP = (Γ
Γ+)/(Γ−+Γ+), where the superscript on the total width
Γ indicates the sign of the b-quark charge in the B me-
son. Some modes can be used for further CP studies. In
fact, the result on B+ → ρ+K∗0 has already been used
to constrain the effect of the penguin amplitude on the
measurement of the angle α of the unitarity triangle from
B0 → ρ+ρ− using SU(3) flavor symmetry [1].
A hot topic is the measurement of the fraction of lon-
gitudinal polarization. The helicity angles θ1 and θ2 of
the two vector mesons are defined as the angles between
the vector meson direction in the B meson rest frame
and the direction of one of its decay product in the vec-
tor meson rest frame, as illustrated on one example in
Fig. 1. Integrating over the φ angle between the decay
planes of the two vector mesons, the fraction of longitu-
dinal polarization fL can be extracted from the angular
θK*0 θρ +K*0
FIG. 1: Definition of the helicity angles in the case of the
vector-vector B+ → ρ+K∗
decay.
∗Presented at the 4th International Workshop on the CKM Uni-
tariry Triangle, Nagoya, Japan, December 12-16, 2006. Preprint
DAPNIA-06-601.
†Electronic address: [email protected]
dependence of the decay rate, which is proportional to
(1− fL) sin
2 θ1 sin
2 θ2 + fL cos
2 θ1 cos
2 θ2.
A value of fL close to unity of order (1−O(
2 )) is ex-
pected for light vector mesons from helicity conservation.
This is expected to be true for both tree and penguin dia-
grams. However the experimental situation is more com-
plex. If fL has indeed been measured close to 1 in the
tree-dominated B → ρρ modes [2], it is surprisingly close
to 0.5 in the penguin-dominated B → φK∗ modes [3].
This effect is not yet understood. There are several possi-
ble explanations, either within the Standard Model, such
as rescattering in the final state, contribution from anni-
hilation or electroweak penguin diagrams, and transverse
gluon [4], or in new physics outside the Standard Model.
To have a better picture, it is important to measure other
vector-vector modes, both tree-dominated, like B → ωρ
and B0 → ωω, and penguin-dominated like B → ρK∗
and B → ωK∗.
The recent studies of the B → ρK∗ modes are reviewed
in section II, the ones involving an ω meson in section III.
Charge-conjugate modes are implied throughout.
II. B → ρK
MODES
A. Introduction
The B → ρK∗ charmless decays proceed through
dominant gluonic penguin loops and doubly Cabibbo-
suppressed tree processes, as shown in Fig. 2. The ex-
ternal tree diagram is only possible with a K∗+, and the
color-suppressed internal tree diagram with a ρ0. Hence
B+ → ρ+K∗0 is pure penguin.
According to isospin symmetry, the two modes with a
charged ρ are expected to have a branching fraction twice
as large as the two modes with a neutral ρ.
FIG. 2: Feynmann diagrams for the B → ρK∗ decay: gluonic
penguin, external tree and internal tree diagrams.
http://arxiv.org/abs/0704.0364v1
mailto:[email protected]
B. Results from Belle
Mbc (GeV/c
∆E (GeV)
FIG. 3: Projections of Mbc for events in the ∆E signal region
(left) and of ∆E in the Mbc signal region (right). The solid
curves show the results of the fit. The dashed curve is the
signal contribution. The hatched histograms represent the
continuum background. The sum of the b → c and continuum
background component is shown as dot-dashed lines.
0.4 0.8 1.2 1.6 2.0
M(π+π0) (GeV/c2)
0.64 0.84 1.04 1.24 1.44
M(K+π-) (GeV/c2)
FIG. 4: Signal yields obtained from the Mbc-∆E distribution
in bins of M(π+π0) (left) for events in the K∗0 region and in
bins of M(K+π−) (right) for events in the ρ+ region. The
points with error bars show the data. Solid curves show the
results of the fit. Hatched histograms are for the nonresonant
component.
Belle was the first experiment in 2005 to publish a re-
sult on the observation of the B+ → ρ+K∗0 mode [5],
on a sample of 275 millions of BB̄ pairs. A signal of
B+ → π+π0K+π− is extracted from the e+e− → qq̄
continuum and BB̄ backgrounds in an extended un-
binned maximum-likelihood fit using the B meson beam-
constrained mass Mbc and energy difference ∆E, as
shown in Fig 3.
The B+ → ρ+K∗0 signal is extracted by fits to Mbc
and ∆E in bins of the vector meson masses M(π+π0)
and M(K+π−), as shown in Fig 4. This is necessary
because there is a large nonresonant ρKπ background,
which gives a continuum in the distribution of M(K+π−).
Nethertheless there is a clear B+ → ρ+K∗0 signal of
85± 16 events with a significance of 5.2 σ.
As for fL, it is obtained by fitting simultaneously the
signal yields obtained from Mbc-∆E fits in bins of the two
helicity angles, assuming an S-wave Kπ system in the
ρKπ background. The results for the branching fraction
and fL in B
→ ρ+K∗0 are:
B = (8.9± 1.7± 1.2) 10−6,
fL = 0.43± 0.11
+0.05
−0.02.
The value found for fL is similar to the one found in φK
and its error is about twice as large as in φK∗.
C. Results from BABAR
(GeV)0π+πm
0.5 1 1.5E
(GeV)0π+πm
0.5 1 1.5E
(GeV)-π+πm
0.6 0.8 1 1.2 1.4E
(GeV)-π+πm
0.6 0.8 1 1.2 1.4E
(GeV)-π+Km
0.8 1 1.2 1.4E
(GeV)-π+Km
0.8 1 1.2 1.4E
(GeV)-π+Km
0.8 1 1.2 1.4E
(GeV)-π+Km
0.8 1 1.2 1.4E
FIG. 5: sPlots for the ππ (top) and Kπ (bottom) invariant
masses in the B+ → ρ+K∗
(left) and B0 → ρ0K∗
/B0 →
f0(980)K
∗0 (right) analyses. The points with error bars show
the data. The solid curve shows the signal and nonresonant
background contribution, the dashed curve is the nonreso-
nant background contribution (ρKπ except for the top right
plot where it represents the sum of f0(1370)K
∗, ππK∗, and
ππKπ). The arrows show the standard mass windows used
in the final fit.
More recently BABAR published an anlysis of all four
B → ρK∗ modes [6], performed on a sample of 232 mil-
lions of BB̄ pairs. It is based on an unbinned maximum-
likelihood fit, using seven variables: the B meson energy-
substituted mass mES and energy difference ∆E, a neural
network output or a Fischer discriminant combining sev-
eral event shape variables, the two vector meson masses,
and the two helicity angle cosines. The fit allows the
simultaneous extraction of the branching ratio and the
fraction of longitudinal polarization.
The major challenge in the analysis comes from the
nonresonant backgrounds, which share the same final
state as the signal. They are studied by enlarging the
vector meson mass windows, as illustrated in Fig. 5. As
in Belle, a large ρKπ background is seen in the mKπ dis-
tribution in the B+ → ρ+K∗0 mode. The Kπ system in
this background is measured to be mostly S-wave. The
situation is even more complex in the B0 → ρ0K∗0 mode,
since in addition to the ρKπ background there are sev-
eral contributions seen in the mππ distribution for a ρ
in contrast to the one for a ρ+. The f0(980) can be seen
clearly. In fact B → f0(980)K
∗, which is a scalar-vector
TABLE I: Results from BABAR on the B → ρK∗ modes: signal yield with its statistical uncertainty, significance (systematic un-
certainties included), branching fraction (90% confidence level upper limit in parentheses), fraction of longitudinal polarization
and direct CP asymmetry. (The numbers in brackets are not quoted as measurements.)
Mode Signal yield Significance (σ) B(×10−6) fL ACP
∗+ 51± 24 2.5 < 6.1 (3.6± 1.7± 0.8) [0.9 ± 0.2]
∗+ 60± 24 1.6 < 12.0 (5.4± 3.6± 1.6)
∗0 194± 29 7.1 9.6± 1.7± 1.5 0.52 ± 0.10± 0.04 −0.01 ± 0.16± 0.02
∗0 185± 30 5.3 5.6± 0.9± 1.3 0.57 ± 0.09± 0.08 0.09 ± 0.19± 0.02
a) b)
c) d)
(GeV)ESm
(GeV)ESm
5.25 5.26 5.27 5.28 )2 (GeV/cESm )2 (GeV/cESm
5.26 5.27 5.28 5.29
FIG. 6: Projections of mES of events passing a signal likeli-
hood threshold for (a) B+ → ρ0K∗
, (b) B+ → ρ+K∗
, (c)
∗+, (d) B0 → ρ0K∗
, (e) B+ → f0(980)K
∗+, and
(f) B0 → f0(980)K
∗0. The points with error bars show the
data. The solid curve is the fit function, the dashed curve
is the total background contribution, and the dotted curve is
the continuum background contribution.
mode, is considered as another signal to be measured in
the same maximum-likelihood fit. Also present are con-
tributions from the f0(1370) and nonresonant ππ. The
yields of the nonresonant backgrounds are fitted in the
enlarged mass windows, then extrapolated to the stan-
dard ones and fixed in the final fit with the standard mass
windows.
The projection plots in the B mass shown in Fig. 6 il-
lustrate the extraction of the signal from the continuum
and BB̄ backgrounds in the four B → ρK∗ channels and
the two B → f0(980)K
∗ modes. Table I summarizes the
results. No significant enough signals are observed for
B0 → ρ−K∗+ and B+ → ρ0K∗+, where upper limits at
the 90 % confidence level are set on the branching ratios.
For the latter a related signal B+ → f0(980)K
∗+ is ob-
served with a significance of 5.0 σ and a measured branch-
ing fraction of (5.2 ± 1.2 ± 0.5) 10−6. In B+ → ρ+K∗0,
the result is in very good agreement with the result from
Belle, with a similar precision. The B0 → ρ0K∗0 mode is
observed for the first time. The ratio between the branch-
ing fractions in these two modes is compatible with the
factor 2 expected from isospin symmetry.
The value of ACP is measured in the two significant
modes to be compatible with 0, as expected since there is
one dominant diagram. Finally fL is found close to 0.5 in
these two modes. It is compatible with the measurement
from Belle and has about the same precision. It is again
similar to the value found for φK∗.
III. MODES WITH ω
40 +ρω
10 ωω
20 *0Kω
20 0ρω
40 +ρω
20 ωω
E (GeV) ∆
-0.2 -0.1 0 0.1 0.2
(GeV/cESM
5.25 5.26 5.27 5.28 5.29
20 0fω
FIG. 7: Projections of ∆E (left) and mES (right) of events
passing a signal likelihood threshold for, from top to bottom,
∗0, B+ → ωK∗
, B0 → ωρ0, B+ → ωρ+, B0 →
ωω, B0 → ωφ, and B0 → ωf0(980). The points with error
bars show the data. The solid curve is the fit function, the
dashed curve is the signal contribution, and the dot-dashed
curve is the background contribution.
TABLE II: Results from BABAR on modes involving an ω meson: signal yield with its statistical uncertainty, significance
(systematic uncertainties included), branching fraction (90% confidence level upper limit in parentheses), fraction of longitudinal
polarization and direct CP asymmetry. (The numbers in brackets are not quoted as measurements.)
Mode Signal yield Significance (σ) B(×10−6) fL ACP
∗0 55± 20 2.4 < 4.2 (2.4± 1.1± 0.7) [0.71 ± 0.25]
∗+ 8± 16 0.4 < 3.4 (0.6± 1.3± 1.0)
−18± 16 0.6 < 1.5 (−0.6± 0.7+0.8−0.3)
+ 156± 32 5.7 10.6± 2.1+1.6−1.0 0.82 ± 0.11± 0.02 0.04 ± 0.18± 0.02
→ ωω 48+24−19 2.1 < 4.0 (1.8
−0.9 ± 0.4) [0.71 ± 0.25]
→ ωφ 3.1±+4.4−8.5 0.3 < 1.2 (0.1± 0.5± 0.1)
ρθcos
-0.5 0 0.5
| ωθ|cos
0 0.5 1
FIG. 8: Projections of the helicity-angle cosines for ω (left)
and ρ+ (right) of events passing a signal likelihood threshold
from the fit for B+ → ωρ+ decays. The points with error bars
show the data. The solid curve is the fit function, the dashed
curve is the signal contribution, and the dot-dashed curve is
the background contribution.
On the same sample of 232 millions of BB̄ pairs,
BABAR has also recently published a search for several
vector-vector modes involving an ω meson [7]: B0 →
ωK∗0, B+ → ωK∗+, B0 → ωρ0, B+ → ωρ+, B0 → ωω,
and B0 → ωφ. The related vector-scalar mode B0 →
ωf0(980) was also searched for. An earlier search for
B → ωK∗ and B → ωρ on 89 millions of BB̄ pairs
resulted in the first observation of the B+ → ωρ+ chan-
nel [8].
The analysis is also based on an extended unbinned
maximim-likelihood fit using the same seven variables as
in the previous section. Nonresonant ππ and Kπ back-
grounds are fixed in the fit as determined from extrap-
olations from higher-mass regions. The projection plots
of ∆E and mES of Fig. 7 illustrate the extraction of the
signal from the continuum and BB̄ backgrounds in all
these modes. In most of them, no significant enough sig-
nal is seen. The only channel where a significant signal is
observed is B+ → ωρ+. Its measured branching fraction
is about 2 standard deviations smaller than the one of
B+ → ρ+ρ0 [2], while these two branching fractions are
naively expected to be equal. Table II summarizes the
results in all the modes. To calculate the branching frac-
tion, fL is left free in the fit for the three modes with a
signal significance greater than 2σ and is fixed otherwise.
Upper limits at the 90 % confidence level are set on the
branching fractions for the modes other than B+ → ωρ+.
The maximum-likelihood fit also provides the value of
fL in B
→ ωρ+, which is found to be 0.82 ± 0.11, a
high value expected for this tree-dominated mode. This
is illustrared in the projection plots of the helicity angle
cosines shown in Fig. 8. The direct CP asymmetry is also
measured and found to be compatible with 0.
IV. CONCLUSION
In summary, improved analyses with explicit consider-
ation of nonresonant backgrounds have been performed
on several charmless hadronic vector-vector decays of
the B meson. The B+ → ωρ+, B+ → ρ+K∗0, and
B0 → ρ0K∗0 modes have been observed and measured in
the past few years. Improved upper limits have been set
on the branching fraction of other vector-vector modes.
The recent results on vector-vector modes have also
brought more pieces to the polarization puzzle. The
penguin-dominated B+ → ρ+K∗0 and B0 → ρ0K∗0
modes have a fraction of longitudinal polarization of
about 0.5 like φK∗, while the tree-dominated B+ → ωρ+
mode has one closer to 1 like ρρ. As a lot of charmless
vector-vector modes have not yet been observed, new re-
sults can be expected with more data.
[1] M. Beneke et al., Phys. Lett. B 638, 68 (2006).
[2] A. Somov, contribution to this conference.
[3] K.F. Chen, contribution to this conference.
[4] G.W.S. Hou, contribution to this conference.
[5] J. Zhang et al., Phys. Rev. Lett. 95, 141801 (2005).
[6] B. Aubert et al., Phys. Rev. Lett. 97, 201801 (2006).
[7] B. Aubert et al., Phys. Rev. D 74, 051102 (2006).
[8] B. Aubert et al., Phys. Rev. D 71, 031103 (2005).
|
0704.0365 | Extending the theory of phonon-mediated superconductivity in quasi-2D | arXiv:0704.0365v1 [cond-mat.supr-con] 3 Apr 2007
7 Extending the theory of phonon-mediated
superconductivity in quasi-2D
J.P.Hague
Department of Physics, Loughborough University, Loughborough, LE11 3TU
Abstract. I present results from an extended Migdal–Eliashberg theory of electron-phonon inter-
actions and superconductivity. The history of the electron-phonon problem is introduced, and then
study of the intermediate parameter regime is justified from the energy scales in the cuprate su-
perconductors. The Holstein model is detailed, and limiting cases are examined to demonstrate the
need for an extended theory of superconductivity. Results of the extended approximation are shown,
including spectral functions and phase diagrams. These are discussed with reference to Hohenberg’s
theorem, the Bardeen–Cooper–Schrieffer theory and Coulomb repulsion. [Published in: Lectures
on the physics of highly correlated electron systems X, p255-264, AIP Conference Proceedings
vol. 846 (2006)]
INTRODUCTION
Over the past half-century, the study of the role of electron-phonon interactions in
condensed matter physics has been an active and controversial field. Initially of interest
from the point of view of thermal properties, early models of the interactions between
lattice vibrations and electrons included the continuum Fröhlich model [1]. Interest in
electron-phonon interactions increased dramatically when in 1957, Bardeen, Cooper and
Schrieffer (BCS) published their famous theory of superconductivity [2], which directly
implicated phonons as the microscopic mechanism for the low temperature absence of
resistivity in a variety of metals. Until the discovery of the cuprate superconductors by
Bednorz and Müller in 1986 [3], the BCS picture was found to account well for all
superconducting materials - a remarkable success for a simple mean-field theory which
is only applicable at weak coupling!
Soon after the realisation that phonons were responsible for superconductivity, Eliash-
berg extended the theoretical description beyond the absolute weak coupling theory with
the famous Eliashberg equations [4]. In doing this, he built on the earlier work of Migdal,
who argued that a simple resummation of a certain class of Feynmann diagrams should
be sufficient to describe the limit of low phonon frequency [5]. Eliashberg’s theory can
be argued to be one of the first applications of the dynamical mean-field theory (DMFT)
[6], since (in its original sense) it ignores spatial fluctuations (momentum dependence)
in the self-energy, while keeping frequency dependent (dynamical) effects.
The purpose of this paper is to describe an extension to the theory of superconductivity
from electron-phonon interactions. The approach goes beyond the Eliashberg theory by
introducing the effects of spatial fluctuations and higher order terms in the perturbation
theory. The aim is to develop a theory which can be used for systems with stronger
coupling, larger phonon frequencies and reduced dimensionality. I begin by motivating
http://arxiv.org/abs/0704.0365v1
the need for a more sophisticated theory from the experimental viewpoint. I also discuss
limiting cases of the Holstein model, and how the large phonon frequency limit of that
model implies that the conventional theories of superconductivity are incomplete. I then
introduce the approximations needed to develop a more sophisticated theory. Finally
I present some results from the new approximation, and discuss them in relation to
Cuprate superconductors, and also with regard to conventional theories and the exact
Hohenberg theorem [7].
MOTIVATION
When the high-temperature cuprate superconductors were discovered in 1986 [3], the
possibility that phonons could be attributed to the microscopic mechanism was quickly
discounted by many people. In part, this was due to the absence of an isotope effect at
optimal doping, and also an assumption that phonon-mediated superconductivity could
not occur above 30K. The mechanism for high-TC superconductors remains highly
controversial, and many different hypotheses are suggested (some examples are spin
fluctuations [8] and exotic phonon mechanisms such as bipolarons [9]). An increasing
body of evidence shows that phonons as well as Coulomb repulsion have an effect on the
physics of the cuprate materials. I shall give a brief review of the current experimental
situation in this section, and argue that (1) Electron-phonon interactions need to be
treated on an equal footing to Coulomb repulsion if the Cuprates are to be understood,
and (2) In order to treat the phonons in the Cuprates, extensions to the current theories
of electron-phonon interactions and phonon-mediated superconductivity are required.
There are several experiments demonstrating strong electron-phonon coupling in the
cuprates. The most compelling is the existence of a strong isotope effect on exchanging
O16 for O18 [10]. There are also some more recent experiments which demonstrate
the effects of electron-phonon interactions in a transparent manner. Figure 1 shows
schematic representations of electron and phonon dispersions in the cuprates. Panel
(a) details the main features of the electronic dispersion measured by Angle-Resolved
Photo-Emmission Spectroscopy (ARPES) in the [11] direction [11]. At energies close to
the Fermi-surface, there are coherent excitations with a long lifetime. As εk = |ω0 − εF |
is approached, the gradient of the dispersion changes at a sharp kink. The phonon is
of the transverse optic variety, and its frequency (ω0) is of the order of 100meV. It
suffices here to mention that this is very large. The ratio of the gradients above and
below the kink is related to the dimensionless coupling constant (λ = g2/tω0), and it is
found that λ can take values of up to 2 [11]. Panel (b) shows a schematic representation
of some neutron scattering results measuring the phonon dispersion [12, 13]. Above
the transition temperature, this looks like the solid line, but as the system moves from
normal to superconducting state, the spectral weight in the circled area vanishes. This
indicates that the superconductivity (bound pairs of electrons) affects the phonons, and
is additional evidence for a strong electron-phonon coupling.
A frequent misconception about the cuprates is that electron-phonon terms in the
Hamiltonian can be neglected on the basis that they are small. To demonstrate that this
is not the case, figure 2 shows approximate energy scales in the cuprates. The largest
energy by far is the Coulomb repulsion (or Hubbard U ) which weighs in at some 10eV.
(a) (b)
INCOHERENT
COHERENT
WEIGHT LOST ON
TRANSITION
kk F 0
PHONON DISPERSIONELECTRON DISPERSION
FIGURE 1. Schematics showing the effect of electron-phonon interactions on the electron and phonon
dispersions in the cuprates. Both panels describe measurements along the [11] direction. Panel (a) shows
a schematic representation of the electronic dispersion measured by Angle-Resolved Photo-Emmission
Spectroscopy (ARPES) [11]. At energies close to the Fermi-surface, there are coherent excitations with
a long lifetime. As εk = |ω0 − εF | is approached, the gradient of the dispersion changes and a kink is
introduced. The phonon is of the transverse optic variety, and its frequency (ω0) is ∼ 75meV. The ratio
of the gradients above and below the kink is related to the coupling constant [11]. Panel (b) shows a
schematic representation of some neutron scattering results measuring the phonon dispersion [12, 13].
Above the transition temperature, this looks like the solid line, but as the system moves from the normal
to the superconducting state, the spectral weight in the shaded area vanishes. This indicates that the
superconducting state affects the phonons, and is further evidence for strong electron-phonon coupling.
Next is the intersite hopping integral t, which is of the order of 1eV. Using a simple
2nd order perturbation theory at strong coupling, an effective exchange interaction is
generated [14], with J = t2/U of the order of 100meV. This J often used to argue for a
spin-fluctuation theory of high-TC superconductivity that neglects phonons. The problem
with this viewpoint is immediately clear if one reviews the experimental data. First, the
energies of the phonons are also approximately 100meV, so they cannot be treated as
a small energy scale. Second, a dimensionless coupling constant of order unity implies
dimensionfull coupling g with similar magnitude. Thus with three very close energy
scales, it is important that the contributions from both phonon and Coulomb mechanisms
are treated on equal footing in a theory for the cuprates. Unfortunately, as I discuss in the
next section, current theories of electron-phonon interactions are not capable of handling
the large phonon energies and coupling constants in the cuprates. The remainder of this
paper focuses on how the theory can be extended to describe this regime.
MODEL AND LIMITS
A generic model of electron-phonon interactions includes the motion of the electrons
Hel, the motion of the ions (or phonons) Hph and the interaction between the electrons
and the phonons (which may be absorbed or emitted) which is denoted Hel−ph. In this
Energy: 10meV 100meV 1eV 10eV
T J t U
FIGURE 2. Schematic showing the energy scales in the cuprates. The largest energy by far is the
Coulomb repulsion (or Hubbard U) of order 10eV. The intersite hopping integral t, is ∼1eV. Using a
simple 2nd order perturbation theory, an effective exchange interaction is generated, with J = t2/U of
the order of 100meV. This J is then used to argue for the spin-fluctuation theory of high TC. However,
the energies of the phonons are also approximately 100meV and the dimensionful coupling g has around
the same value. Thus with 3 similar energy scales, it is important that the contributions from both spin-
fluctuations and phonon mechanisms are treated on equal footing.
way, H = Hel +Hel−ph +Hph is the total Hamiltonian.
Hel = ∑
kck ≈− ∑
<i j>σ
tc†iσ c jσ (1)
Hel−ph =−∑
k−qck(b
q +b−q)≈−∑
niσ gri (2)
Hph = ∑
b†kbk +
Mω20 r
The first term in the Hamiltonian is the general form for free electrons, i.e. the total
energy is the sum of the kinetic energies of all occupied states. In a special case,
which is known as the Holstein Hamiltonian, the electrons in a tight binding model
may hop between nearest-neighbour sites only, and εk = −2t ∑Di=1 cos(ki), where t is
the overlap integral. In the generic form of the electron-phonon interaction, an electron
may be scattered by absorbing a phonon with momentum −q or emitting a phonon with
momentum q. An additional approximation uses a momentum independent electron-
phonon coupling, g, and in that case the Fourier transform shows that the second
term connects the local ion displacement, ri to the local electron density. Finally, the
free phonon term may be simplified by using the Einstein approximation ωk ≈ ω0
and Fourier transforming, the bare phonon Hamiltonian is shown to be a series of
independent simple harmonic oscillators at each site index. The creation of electrons
and phonons is represented by c† and b† respectively, pi is the ion momentum and M the
ion mass. By choosing t = 0.25, a bandwidth of W = 2 is chosen. A small interplanar
hopping of t⊥ = 0.01 is included to remove the logarithmic singularity in the 2D density
of states at ε = 0.
Figure 3 shows the parameter space of the Holstein model. For very large phonon
frequency, the effective interaction is instantaneous, and a Lang–Firsov transformation
[15] results in an attractive Hubbard model (which is one of the standard models for
correlated electron systems) [16]. Alternatively, taking the limit of very small phonon
frequency, a fast moving electron cannot ‘see’ the nuclei move in the time it takes to
EXTENDED THEORY
CUPRATE?
FIGURE 3. Parameter space of the Holstein model. For very large phonon frequency, the effective
interaction is instantaneous, and a Lang–Firsov transformation results in an attractive Hubbard model.
Alternatively, taking the limit of very small phonon frequency, a fast moving electron cannot ‘see’ the
phonons move, and the problem maps to a static disorder problem (similar to the Falikov–Kimball model
[19]). This makes the phonon problem extremely hard, and little is known about the middle of the
parameter space. The range of the Eliashberg theory is shown in the bottom left corner. The expected
position of the cuprates is shown as the single diamond. The expected validity of an extended theory
including all 2nd order Feynman diagrams is also shown.
traverse many sites, so the problem maps to a static disorder problem (which is essen-
tially uncorrelated). One may therefore think of the phonon frequency as possessing the
ability to “tune” the effect of correlations, and one therefore obtains a second motivation
for the study of electron-phonon systems of trying to understand electronic correlations
[17]. The correlation tuning makes the phonon problem extremely hard, and little is
known about the intermediate regime of the parameter space. The range of the Eliash-
berg theory is shown in the bottom left corner. Contrary to Migdal’s assumption, the
theory cannot extend beyond intermediate coupling since renormalisation of the effec-
tive mass reduces εF invalidating the condition (Migdal’s theorem) ω0 ≪ εF [18, 9]. The
approximate position of the phonon parameters in the cuprates is shown as the single di-
amond. It is essential to correct the theory for weak to intermediate coupling at larger
phonon frequencies. The extension is clear by looking at the large phonon frequency
limit. The Hubbard limit requires that all 2nd order processes in U are included in the
self-energy, or the incorrect weak coupling limit is found. An extended theory including
all 2nd order Feynman diagrams is required to understand the weak coupling limit, from
small to large phonon frequency.
FIGURE 4. Series of Feynman diagrams used in the current approximation. Σ is the electron and Π the
phonon self-energy. Series (a) is the Migdal-Eliashberg approximation and (b) the vertex corrected series.
EXTENDING THE ELIASHBERG THEORY
Extending the Eliashberg theory involves inserting the lowest order vertex corrections
into the electron and phonon self energies. In the Eliashberg theory, emitted phonons
are reabsorbed in a last-out-first-in order. Vertex corrections essentially allow this order
to be changed once. Such contributions are shown diagrammatically in figure 4. All
the diagrams must be included in the calculation, or electron number would not be
conserved. Momentum dependence is included in the approximation, which is essential
in low-dimensions. The inclusion of vertex corrections leads to double 2-fold integration
over the Brillouin zone in combination with a double sum over matsubara frequencies,
which is time consuming for the numerics. In order to reduce the number of points in k-
space while maintaining the thermodynamic limit, the dynamical cluster approximation
is applied [20]. Additionally, superconducting states can be considered by using the
Nambu formalism. The full details of the implementation of the extended approximation
can be found in references [21] and [22].
Using a maximum entropy technique, it is possible to compute the spectral function
from the Matsubara axis Green function. Figure 5 shows the spectral function of the
Holstein model calculated using the extended Migdal–Eliashberg theory. The results are
qualitatively similar to ARPES measurements of the cuprates. In particular the change
between incoherent and coherent particles occurs at the phonon frequency (shown as the
dashed line), associated with a kink in the [11] direction. It is noted here that the effect
of the phonon self-energy is a softening of the phonon mode. In the standard ME theory
in 2D, the mode at the (π ,π) point is completely softened, leading to a fatal instability
of the theory. However, the vertex corrections act against this softening, and relieve the
instability. In such a way, it is clear that a vertex corrected Eliashberg theory is essential
for the study of quasi-2D materials [21].
One can also compute properties in the superconducting state. One such property is
the momentum-dependent pairing density, ns(k) = T ∑n F(iωn,k), where F(iωn,k) is
the anomalous Green function associated with the pairing of electrons with momentum
k and −k. It is possible to transform the momentum dependent order parameter to
determine the magnitude of individual spherical harmonics. Figure 6 shows such a
decomposition. A cluster size of NC = 64 is used, with U = 0.6 and ω0 = 0.4. Note
ω0=0.2, U=0.3, DCA(VC)
-2 -1
A(k,ω)
FIGURE 5. Spectral function of the Holstein model in the extended Migdal–Eliashberg theory. The
results are qualitatively similar to ARPES measurements of the cuprates. In particular the change between
incoherent and coherent particles occurs at the phonon frequency, associated with a kink in the [11]
direction. ©Institute of physics publishing 2003 [21].
how higher order harmonics develop as the filling is increased. In particular, it can be
seen that no single harmonic (such as the s-wave symmetry) is sufficient to describe the
order parameter. Some of the higher order terms come about due to increased pairing at
momentum k = (π/2,π/2), in particular, pairs with angular momentum.
Finally, by varying the temperature and chemical potential, the phase diagram can
be computed. Figure 7 shows phase diagrams of the Holstein model for the different
approximations. U = 0.6 and ω0 = 0.4.The top diagram shows the result from the
Eliashberg approximation (dynamical mean-field theory NC = 1). On the bottom the
results from the current approximation with NC = 4 are shown. The superconducting
order is suppressed close to half filling. Assuming a form for the density of states in 2D
(with small interplane hopping) of D(ε) = (1− t log((ε2+ t2⊥)/16t
2))/tπ2 (for |ε|< 4t)
[23], which matches the full density of states with reasonable accuracy. From this the
BCS result may be calculated using the expression
TC(n) = 2ω0 exp(−1/|U |D(µ(n)))/π , (4)
with the chemical potential taken from the self-consistent solution for a given n. This
result also drops off monotonically. Results in the dilute limit are in good agreement with
the BCS result (line with points). Close to half-filling, the DMFT result is significantly
smaller than the BCS result (which predicts TC(n = 1)> 0.07). The difference in results
-0.025
-0.02
-0.015
-0.01
-0.005
0.005
0.01
1 1.1 1.2 1.3 1.4 1.5 1.6 1.7
s, m=0
d, m=0
g, m=0
g, m=4,-4
FIGURE 6. Decomposition of the order parameter into spherical harmonics. A cluster size, NC = 64 is
used, with U = 0.6 and ω0 = 0.4. Note how higher order harmonics develop as the filling is increased. In
particular, the g harmonics can be almost as strong as the s harmonics at n = 1.45. ©Institute of physics
publishing 2005 [22].
between the two mean-field theories at half-filling is due to the self-consistency in the
DMFT. When vertex corrections and spatial fluctuations are included, the dilute limit
is relatively unchanged. However at half-filling, there is a huge drop in the transition
temperature. The suppression at half-filling is a manifestation of Hohenberg’s theorem,
which implies that there may be no superconducting order in 2D. Here I have computed
for quasi-2D, so it is interesting that in real materials with low dimensional character the
maximum in superconductivity is shifted away from half-filling.
CONCLUDING REMARKS
I end the paper with a warning for constructing theories of high-temperature supercon-
ductivity using electron-phonon interactions alone, while neglecting the Coulomb re-
pulsion. If one takes the phase diagrams from the previous section, and assigns similar
energy scales to those in the cuprates, it is possible to obtain a temperature in Kelvins
for the maximum in the phase diagram at n = 1.2. This comes out as around 172K - one
could say approximately the TC in the cuprates.
So why isn’t this the solution for the cuprates? Cuprates are very tightly bound ma-
terials, which is why the “Fermi energy” is low, and the ratio ω0/εF is large enough
to justify extending Eliashberg theory. The problem is that a small Fermi energy also
means the the Hubbard U is a comparatively large quantity. On a simple mean-field
level, one can include the Coulomb repulsion in the theory of superconductivity. For ex-
ample, the Eliashberg equations can be extended to include an effective electron-electron
interaction (otherwise known as the Coulomb pseudopotential µC). The effect of this is
to modify λ → λ − µC. Substitution into equation 4 means that the transition temper-
0.02
0.04
0.06
0.08
U=0.6, ω0=0.4, Nc=4, VC
0.01
0.02
0.03
0.04
1 1.2 1.4 1.6 1.8
0.02
0.04
0.06
0.08
0.02
0.04
0.06
0.08
0.01
0.02
0.03
0.04
1 1.2 1.4 1.6 1.8
0.02
0.04
0.06
0.08
U=0.6, ω0=0.4, Nc=1
FIGURE 7. Phase diagrams of the Holstein model. U = 0.6 and ω0 = 0.4. The top diagram shows the
result from the Eliashberg approximation (dynamical mean-field theory NC = 1). Also shown is the BCS
result (line with points). On the bottom the results from the current approximation with NC = 4 are shown.
The superconducting order is suppressed close to half filling in the vertex corrected theory. ©Institute of
physics publishing 2005 [22].
ature is considerably reduced, or that superconductivity of the BCS type is completely
destroyed. Any phonon-based mechanism for the cuprates must address this point and
be compatible with the electron-electron interaction. Alternatively (and this is a warn-
ing against the other extreme) on the basis of the similarity of energy scales, any spin-
fluctuation mechanism (which is essentially Coulombic) must also treat the phonons (or
at least be compatible with them) to be plausible.
ACKNOWLEDGMENTS
I sincerely thank the organising committee of the course for their generous financial
support. Aspects of this research were carried out under the MPIPKS guest scien-
tist program, and as a visitor at the University of Leicester. I thank A.S.Alexandrov,
J.L.Beeby, E.M.L.Chung, N.d’Ambrumenil, J.K.Freericks, M.Jarrell, P.E.Kornilovitch,
J.H.Samson and M.Yethiraj for stimulating discussions, both about this work and the
problems of electron-phonon interactions and superconductivity in general. I acknowl-
edge support at Loughborough University under EPSRC grant no. EP/C518365/1.
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|
0704.0366 | Generalized Nariai Solutions for Yang-type Monopoles | Generalized Nariai Solutions for Yang-type
Monopoles
Pablo Diaz∗, Antonio Segui†
Departamento de Fisica Teorica, Universidad de Zaragoza, 50009-Zaragoza, Spain.
October 28, 2018
Abstract
A detailed study of the geometries that emerge by a gravitating generalized
Yang monopole in even dimensions is carried out. In particular, those which
present black hole and cosmological horizons. This two-horizon system is ther-
mally unstable. The process of thermalization will drive both horizons to coalesce.
This limit is what is profusely studied in this paper. It is shown that eventhough
coordinate distance shrinks to zero, physical distance does not. So, there is some
remaining space which geometry has been computed and identified as a generalized
Nariai solution. The thermal properties of this new spacetime are then calculated.
Topics, as the elliptical relation between radii of spheres in the geometry or a dis-
cussion about whether a mass-type term should be present in the line element or
not, are also included.
Keywords: Yang monopole, Nariai geometry, Horizon, Black hole.
∗[email protected]
†[email protected]
http://arxiv.org/abs/0704.0366v2
Contents
1 Introduction 2
2 The gravitational coupling. Some geometrical features 3
3 The horizon coalecence geometry 4
3.1 Case m = 0 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5
3.2 Case m 6= 0 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 7
4 Conclusions 11
A Proof of the finite nonzero physical distance 12
B Horizon coalescence as a flow on the line 13
1 Introduction
Monopoles have been subject of deep study and controversy all over the last century. This
is so because, although no experimental evidence of their existence has been found, many
theoretical issues make them almost unavoidable. They already appeared as solutions of
Maxwell equations as long as the null B-divergence condition was relaxed, that is, ∇·B 6=
0. It was Dirac [1] in the early thirties who first proposed the theoretical possibility of
creating an experiment to actually produce a “fake” monopole, in a way that its fakeness,
say, the Dirac string, was undetectable. As a consequence, the product of the electric
and the magnetic charges was quantized. Many years later, in 1959, the quantization
requirement was confirmed by the celebrated Aharonov-Bohm experiment [2].
Since 1954, owing to the papers by Yang and Mills [3] and by Utiyama [4], gauge
theories of a group of symmetry larger than U(1), in particular non abelian symmetry
groups SU(2) and SU(3) (which eventually would conform the Standard Model of par-
ticle physics) where gradually developped. In 1969, Lubkin [5] realized that monopoles
can be classified by the homotopy group of the gauge symmetry group of the theory,
so that the magnetic charge is replaced by the topological charge of the field configura-
tion. In the case of the Dirac monopole, the homotopy group π1 of U(1) is exactly Z.
However, it was not till 1975 that Yang [6] generalized the abelian monopole to the case
of an SU(2)-invariant gauge theory in six dimensions, see also [7]. Modern approaches
use the formalism of fiber bundles for a suitable description of monopoles. It generalizes
the traditional classification in terms of the homotopic group of the gauge theory. In
this way, magnetic monopoles are identified with the different instanton configurations
which come up basically as non trivial maps of the gauge group, usually SU(N), onto Sd,
where d is the spatial dimension. That is, magnetic monopoles are all those non trivial
principal bundles with group structure SU(N) that can be realized on the hypersurface
Sd. The classification coincides, as said before, with the different classes of homotopy
groups. The genalization of Yang monopoles to an arbitrary even dimension was carried
out in [8]. Using slightly different methods similar analysis have recently been done [9].
The reader can find good reviews on the subject in [10], [11] and the references therein.
As every existing object in nature, monopoles couple to gravity via their energy-
momentum tensor. The resulting geometry is obtained by solving the Yang-Mills-
Einstein equations, which get greatly simplified by imposing spherical symmetry (as
expected from a magnetic monopole field configuration). This geometry is fully specified
by choosing a point in the space of parameters {µ,m,Λ, k}, the meaning of which will
be explained in detail later on. For a given range of parameters, it is easy to prove
that the geometry presents both a cosmological and an event horizon. A full analogy
with the Schwarzschild-de Sitter solution reveals that, in these cases, the geometry is dy-
namically driven through the parameter space into a thermally stable point where both
horizons coalesce [12], the final line element being the analogue of Nariai’s spacetime in
four dimensions.
This paper is organized as follows: the next section sets a general framework and
fixes the notation used later. The main body of the article concerns the analysis of the
coalescence solutions. This is achieved in two subsections corresponding to the massless
and massive cases respectively. An explicit computation of the resulting geometry is
carried out in each case. A final section includes some conclusions and comments. Two
appendixes have been added to the article. They are topics which lie somehow out of the
main line of the paper, either for being technical aspects of a computation (Appendix
A) or for presenting a new idea the exposition of which would need a new section, as in
Appendix B. the absence of B, in turn, would not have prevented the reader from a full
understanding of the paper.
2 The gravitational coupling. Some geometrical fea-
tures
The gravitational effects of these monopoles have been recently studied [9]. It was done,
as usual, by minimally coupling the Yang-Mills energy-momentum tensor to gravity.
Variations of the Einstein-Hilbert action
− det g
(R− 2Λ)− 1
Tr|F |2
with respect to the metric tensor leads to
Gmn = 8πGTmn − gmnΛ, (2)
where
Tmn = γ
tr(F pm Fnp)−
gmntr(FpqF
is the energy momentum tensor of the YM strength field. The traces are taken in the
colour index and γ is the YM coupling constant. Finding general solutions for (2) is
a highly complicated problem. However, imposing spherical symmetry simplifies the
task enormously. According to this, the ansatz will be a spatially spherically symmetric
(2k + 2)-dimensional metric whose line element reads
ds2 = −∆dt2 +∆−1dr2 + r2dΩ22k. (4)
The last equation is consistent with (2) and (3) when [9]
∆(r) = 1− 2Gm
r2k−1
, (5)
where R =
k(2k+1)
is the de Sitter radius, µ2 is proportional to 1
and measures
the magnetic charge of the monopole, m comes up as a constant of integration with
dimensions of mass and G is the Newton constant in 2k + 2 spacetime dimension. At
first sight, (4) with (5) look like a Schwarzschild-de Sitter geometry in 2k + 1 spatial
dimensions with an extra term, the one involving µ, which seems to be independent of
the dimension of spacetime. It seems reasonable to think of this term as a contribution
of the magnetic monopole. This simple image, even if not exact1, is helpful and, unless
we face the vanishing limits, it may be kept in mind in the following.
The next step (and the next temptation) is to analyze how the causal structure
of this spacetime depends on given values of the parameters. The main body of this
work concerns a deep analysis of the solution in the case when parameters µ,Λ, m and
k allow the existence of two horizons. Then, inspired by the Schwarzschild-de Sitter
unstable solution, it is claimed that the system gets dynamically driven to a value of the
parameters where both horizons coalesce. Eventhough coordinate distance shrinks to
zero, physical distance does not. A generalized Nariai geometry “between” the horizons
is then explicitly obtained. The Nariai line element [13] is a nonsingular solution of the
Einstein’s vacuum equations with a positive cosmological constant, Rµν = Λgµν . It was
first found by Kasner [14] and its electrically charged generalization dates of 1959 [15].
However, the important fact that it emerges as an extremal limit of Schwarzschild-de
Sitter black holes was not noticed until 1983 [12].
Nariai spacetime in four dimensions is the direct product dS2 × S2, dS2 being no
more than the hyperbolic version of S2 as we change t → iτ . In 2k + 2 dimensions, the
solution gets generalized to dS2 × S2k. Again, it is the direct product of two constant
curvature spaces and admits a 3 + k(2k+ 1) group of isometries SO(2, 1)× SO(2k+ 1).
The space is homogeneous since the group acts transitively and is locally static, given
that a global dS-type spacetime cannot be described by merely one static coordinate
chart. In four dimensions, radii of curvature of the two product spaces are equal if the
black hole is neutral, and different in the charged case. If the black hole is electrically
charged, the respective radii a and b are different and related by the equation
a−2 + b−2 = 2Λ (6)
as shown in [15]. This relation will be generalized in the magnetic case, the object of
our study. A short but instructive recent work on the four dimensional geometry can be
found in [16].
3 The horizon coalecence geometry
Studying the horizons of a geometry like (4) is equivalent to searching the divergencies
of grr for finite values of the coordinates. This leads us to analyze the zeroes of function
∆(r), where the horizons will be located. For a certain range of values of {µ,Λ, k,m}
there will be two horizons. Finding this region in the parameter space will be the first
1The resulting geometry is, of course, not just the sum of terms of different geometries, but it casually
coincides. Differences are bound to exist on the limit of vanishing of a given contribution. For instance,
let us suppose that, given a set of parameters, say {m,µ,Λ, k }, we can switch off µ (by neglecting it
with respect to the others). The resulting geometry is topologically different to the one obtained by not
assuming any monopole at all at the beginning, that is, the limit does not coincide. However, in the
cases studied here, this is no more than an enough-to-be-aware-of subtlety.
task. After that, attention will be focused on the coalescence point of the horizons2. The
analysis consists of two steps, first, the parameterization of the coordinate separation
of the horizons (ǫ) and the calculation of the physical distance between them when
coalescence takes place (ǫ → 0). Then, following the strategy in [12], the computation
of the line element of the remaining geometry. This program is carried out on two cases:
m = 0 and m 6= 0, which are treated in the next subsections, respectively. The massless
case must be seen as a toy model of the massive one. This distinction is made not
merely for simplicity but also because, as will be explained, the mass parameter comes
out naturally for dynamical requirements.
3.1 Case m = 0
In the massless case, ∆(r) gets reduced to
∆(r) = 1−
. (7)
Solving ∆ = 0 is equivalent to finding the zeroes of a biquadratic equation as long as
r = 0 is not considered. We perform the change z ≡ r2 and solve a second order ordinary
equation. The horizons are found to be at
1− 4µ
z++ =
1− 4µ
. (9)
R > 2µ guarantees the existence of two positive solutions and, therefore, four solutions
for the quartic equation. Two of them, r+ = +
z+ and r++ = +
z++, correspond to the
radial coordinate of the inner (black hole) and outer (cosmological) horizon respectively.
If R = 2µ, both solutions coincide, which means that the horizons coalesce. As said
before, this does not mean that the geometry vanishes as a naive observation (given
a wrong choice of coordinates) would make one think. Physical distance between the
horizons, on the contrary, remains finite at the limit. In order to prove this, let us
compute it. For fixed time and angular coordinates, the physical distance is
D(µ,R) =
∫ r++
[−r4 +R2r2 − µ2R2]1/2
∫ z++
− µ2R2
z − R2
)2]1/2dz (10)
The requirement R > 2µ implies R
− µ2R2 > 0 so the above integral is exactly solved
as an cos−1-type. The result is
D(R) =
πR. (11)
Surprisingly, the physical distance does not depend on µ. It means that, given a cos-
mological constant Λ, one could “switch on” the monopole and go on till the horizons
coalesce but the distance would remain unchanged. However, because of quantization
2Coalescence as seen in Schwarzschild coordinates.
requirements, monopole charge µ cannot be tuned, but needs to have, instead, a fixed
value upto a sign. On the other hand, the cosmological constant, Λ, should be chosen
when writing the lagrangian. It means that changing its value does not drive us from
one model to another but implies an essential change in the theory [17]. Therefore, we
are not free to adjust any parameter arbitrarily as done with the mass of the black hole
in the Schwarzschild-de Sitter case. Then, eventhough physical reasons would lead the
horizons to coalesce, the absence of any free parameter in our model makes it impossible.
In the next section, m will come to our help as a free parameter for the model.
Despite the last remark, one could wonder about the kind of geometry that re-
mains when the horizons coalesce. This task, even if seems just a curious exercise now,
will be useful for the next section. Applying a technique similar to the one Gingspar
and Perry [12] used to study the geometry of Nariai’s solution, we proceed by, first,
parametrizing the separation of horizons as
R = 2µ(1 + ǫ2), (12)
in a way that coalescence corresponds to taking ǫ = 0. Then, we define a “wise” change
of coordinates
χ = cos−1
(r2 − r20)
τ = ǫ
, (13)
where A =
1− 4µ2
and r20 =
, and the angular coordinates remain unchanged. The
new coordinates (13) might seem randomly chosen at first sight. However, there are
some reasons that justify such a functional dependence. For instance, χ is nothing but
the physical distance between r+ and r. The timelike coordinate t is multiplied by i in
order to work in the Euclidean region3 and by ǫ because ∆/ǫ2 is expected to have a finite
limit when ǫ → 0. Now, we apply (12) and (13) and expand ∆(r(χ))dτ 2, ∆−1(r(χ))dχ2
and r2(χ) up to first order in ǫ. The line element (4) reads
ds2 = µ2dχ2 + µ2 sin2(χ)
1 + ǫ
2 cos(χ)
dτ 2 +
+ 2µ2
2 cos(χ)ǫ
dΩ22k. (14)
We take limit ǫ → 0 to obtain
ds2 = µ2
dχ2 + sin2(χ)dτ 2
+ 2µ2dΩ22k. (15)
As seen in (15), the 2k-sphere decouples from the rest. The resulting geometry is S2×S2k
for k ≥ 2. Notice the parallelism between this geometry and Nariai’s solution, which is
S2 × S2. The “classical” relation between radii (6) gets also generalized to
a−2 + b−2 = C0Λ, (16)
where C0 =
k(2k+1)
. The geometry (15) can be viewed as a “degenerate” black hole, in
which the two horizons have the same (maximum) size and are in thermal equilibrium.
This could be interpreted by an observer as a bath of radiation coming from both horizons
3τ will be periodic at both horizons, although different in each case. Equality will hold at the
coalescence point, when thermal stability is reached.
at a precise temperature [19]. The temperature can be calculated by means of surface
gravity κ, as computed in the new coordinates (13)
k(2k + 1)
. (17)
The entropy can also be computed as a quarter of the sum of the two horizons [18], so
k(2k + 1)
, (18)
where ω2k is the area of the 2k-dimension unit sphere.
3.2 Case m 6= 0
In the massive case we recover the full expression (5) for ∆. Since the singular point
r = 0 is not to be considered, we better analyze the function r2k−1∆(r)
∆̃ ≡ r2k−1∆ = −
r2k+1
+ r2k−1 − µ2r2k−3 − 2Gm. (19)
It is known that a polynomial equation with powers equal to or higher than five is not
generally solvable in a symbolical way. This happens for k ≥ 2. So, the purpose of doing
a study for the massive case analogous to that achieved in the first section is ruined.
Nevertheless, some information can be extracted from (19). We should first remember
the sign of the parameters: R2 > 0 (de Sitter), µ2 > 0 for k ≥ 2, and m will be free in
principle. Derivating (19) and equating to zero leads to a biquadratic equation of the
(2k + 1)r4 + (2k − 1)r2 − (2k − 3)µ2 = 0, (20)
which, as long as
Λµ2 ≤ k
(2k − 1)2
2k − 3
, (21)
has two positive (and two negative) roots, rmin and rmax ≡ rc. In terms of the cosmo-
logical constant
r2c ≡
k(2k − 1)
1− 4(2k − 3)Λµ
k(2k − 1)2
, (22)
rmin is obtained from (22) by swapping the sign of the square root. A quick look at (19)
shows that the smallest root is a minimum and the largest is a maximum of function ∆̃.
Now, let us plug rc into (19):
1. If m > 0, then (see fig.1)
a) ∆̃(rc) ≥ 0 implies that there are two event horizons, the black hole and the
cosmological horizon. The inequality gets saturated at the coalescence point.
b) ∆̃(rc) < 0 means that no horizon is found.
2. If m < 0, then (see fig.2)
a) ∆̃(rmin) < 0 together with ∆̃(rc) < 0 implies that there is just one Cauchy
horizon.
b) ∆̃(rmin) < 0 together with ∆̃(rc) > 0 assures the existence of a Cauchy horizon
and both black hole and cosmological horizon.
c) ∆̃(rmin) > 0 leaves us with the cosmological horizon only.
The case we will study is ∆̃(rmin) < 0 and ∆̃(rc) > 0 which, independently of the sign
of m, assures4 the existence of black hole and cosmological horizons. This corresponds
to values of m within range (see fig.3)
1.a 1.b
Figure 1: Case m > 0. The curve represents function ∆̃(r). Figure 1.a has two roots which
correspond to the black hole (r+) and cosmological horizon (r++) respectively. Figure 1.b
shows the absence of horizons.
2.a 2.b 2.c
Figure 2: Case m < 0. This time ∆̃(r) permits the existence of one (Cauchy) horizon as
in Figure 2.a, three horizons (Cauchy, black hole and cosmological) as in 2.b, or just the
cosmological horizon as shown in 2.c.
m− < m < m+, (23)
where
Gmc ≡ Gm+ =
1 + 2k
r2k−3c (r
c − 2µ2). (24)
The value of Gm− is obtained by replacing rc → rmin. In terms of Λ and µ we get
Gm± =
(2Λ)−k+1/2
1 + 2k
− k + 2k2 ±
k2(1− 2k)2 − 4Λµ2(2k − 3)k
]k−3/2
− k + 2k2 − 4Λµ2 ±
k2(1− 2k)2 − 4Λµ2(2k − 3)k
. (25)
4The value of m can be negative. That is because m should not be thought of as an entity with
physical meaning but as a geometrical parameter. Short calculation in (25) shows that m gets negative
values for Λµ2 ≥ k
(1 + 2k).
Hr,m+L
Hr,m-L
Figure 3: This figure shows the range of “masses” which are consistent with the existence of
both black hole and cosmological horizons. The curve ∆̃(r) “moves down” in the process of
coalescence.
The crucial point is that both horizons coalesce when rc is a root of (19) which
happens at m = mc(k,Λµ
2,Λ). Two relations have been imposed so far:
de∆(r;m)
|rc= 0,
that is, (20), which defines rc, and ∆(rc;mc) = 0 which leads to mc. In order for m to
be real, the bound which must be impossed on Λµ2 coincides with (21) which, in turn, is
nothing but the condition for the existence of two horizons. So, if a given a value for Λµ2
is low enough to produce two horizons, there always exists a real value of m which makes
them coalesce. Again, as in the Schwarzschild-de Sitter example, the system is unstable
and the equilibrium point is reached at m = mc. Unlike the massless case, plugging m
gives us enough room for maneuvre to drive the system to equilibrium.
At this point, we would like to remark that the procedure of horizon coalescence, as
studied in detail below, may be seen as a flow in a line which undergoes a Pitchfork
bifurcation at the coalescence point. Parameter m, moved by thermal instability, drives
the system to the critical situation. For concreteness see Appendix B.
Let us focus on the near coalescence point. This can be parameterized by
r = rc + δr = rc(1 + ǫ cosχ) (26)
m = mc − δm = mc(1 + bǫ2).
Parameterization of r also involves a change of coordinates r → χ and should be taken
as imposed at the moment although it will be justified later. The horizons will be
symmetrically located at: r+ = rc(1−ǫ) and r++ = rc(1+ǫ) which correspond to χ+ = π
and χ++ = 2π, respectively
5. The value of b as well as the absence of a linear term in ǫ of
the parameterization of m may be explained as follows. Near the coalescence point one
should Taylor expand ∆ around rc and have in mind that, for
≪ 1, ∆ is aproximately
parabolic, so that second order expansion is enough. By definition ∆(r+) = ∆(r++) = 0
5For a small enough ǫ, it is expected that the parabolic approach holds and, then, both horizons are
symmetrically located with respect to rc.
and ∆ reaches a maximum at rc. So,
0 = ∆(r++) = ∆(rc, m) + ∆
′(rc, m)(rcǫ) +
∆′′(rc, m)(rcǫ)
r2k−1c
∆′′(rc;mc)r
2, (27)
which means that
∆′′(rc;mc)r
. (28)
Calculating the physical distance near the coalescence point would, again, imply
solving the integration
D(ǫ) =
∫ r++
∆1/2(r)
, (29)
where r++ = r+ + 2rcǫ. Although the exact result is not computed, an explicit proof of
its finite nonzero value is given in Appendix A. The procedure of calculating the physical
distance also brings us some light on which is the change of coordinates that should be
made in order to understand the resulting geometry. It turns out to be
χ = cos−1
(r − rc)
t, (30)
where
k − 2k2 + 2Λr2c
is a dimensionless factor.
The coalescence of horizons takes place at ǫ = 0. In order to study the geometry at
the limit we proceed by calculating −∆dt2, ∆−1dr2 and r2 in the new coordinates (30)
and expand in ǫ around ǫ = 0. The new line element gets determined by taking the zero
order of the expansion. The relations for r and m in (26) are in accordance with (30),
where b takes the value of (28), by virtue of the parabolic approach. From (30), it is
straightforward to see that r2 takes a constant value r2c . Surprisingly, as in the massless
and Schwarzschild-de Sitter cases, the geometry splits in two disconnected parts which
lead to a product manifold S2 × S2k. The line element reads
ds2 = Br2c
dχ2 + sin2(χ)dτ 2
+ r2cdΩ
2k, (32)
where χ ∈ [π, 2π] and τ is periodic6.
As seen in (32), S2 has radius a2 = Br2c , and S2k has radius b2 = r2c . Now, the
generalized Bertotti relation (6) is
a−2 + b−2 =
2(1− k)
= CΛ, (33)
where C(k,Λµ2) is obtained by inserting (22) in (33). Note that C
k,Λµ2 = k(2k +
= C0, and then (33) turns into (16), that is, into the massless case. This is no
6τ is periodic on both horizon surfaces all over the process in order to avoid the conical singularity
at the horizons. At the coalescence point, however, both periods equal.
surprising since Λµ2 = k(2k + 1)/4 is the condition for coalescence in the massless case
(equivalent to R = 2µ), and, at the same time, it makes mc = 0. So, the massive
geometry is a consistent extension of the massless one. Now, fixing Λ does not determine
uniquely the geometry. Another dimensionless variable Λµ2 is required.
As in the last section, the geometry (32) can be viewed as a “degenerate” black hole,
in which the two horizons have the same (maximum) size and are in thermal equilibrium.
In the present case the temperature is given in terms of the surface gravity κ by
. (34)
In Planck units,the entropy associated with this solution may be calculated (given that
it is not extreme7) by means of the total area of the horizons as
c . (35)
4 Conclusions
The spherically symmetric solution of gravity due to a magnetic monopole in arbitrary
dimension has been studied, in particular, when the set of parameters {Λ, µ,m, k} allows
the existence of two horizons. In these cases, thermal instabilities drive a process of
horizon coalescence. Even though coordinate separation between the horizons shrinks
to zero, it has been proven in both the massless and the massive case that the physical
distance does not. The geometry of the remaining space between the horizons has been
calculated in both cases. They turned out to be Nariai-type solutions, that is, the product
of a 2-sphere and a 2k-sphere for a (2k + 2)-dimensional spacetime. In each solution,
the radii of the spheres are not independent. They are related by an elliptical equation
which should be understood as the generalization of the relation found by Bertotti. The
unique generalized equation involving these radii for both the massless and the massive
case has been given. After computing the line element in each case, the thermodynamical
properties (Hawking temperature and entropy) due to the existence of horizons have been
calculated.
The Yang monopole corresponds to the six dimensional case, where k = 2. The
geometry obtained after coalescence is S2 × S4 as can be explicitly read in (32). This
case is especially interesting since it may be described in String Theory (a realization
of the Yang monopole in Heterotic String Theory has recently been done [21] as well as
another complementary picture in Type-IIA String Theory [20]). In the same context,
it looks possible to find results (18) and (35) for the entropy by application of some
attractor mechanism [22, 23]. We believe that this would be an interesting topic to be
addressed in future research.
7A charge black hole is said to be extreme when it has the minimum mass. Then, as it cannot release
any energy without losing charge, it is supposed not to emit, and its associated Hawking temperature is
0. The black hole we are dealing with in this paper is extreme in the sense of carrying the “maximum
mass” allowed by the cosmological constant Λ. Obviously, the temperature will not be zero.
A Proof of the finite nonzero physical distance
Computing the physical distance is equivalent to performing the integration
∫ r++
∆1/2(r)
, (36)
where, for small ǫ, r++ = r+ + 2rcǫ. Divergencies might appear at the points where
∆ → 0. The case we have been considering all along section (3.2) concerns the existence
of two horizons which coalesce, that is, two single roots r+ and r++ of ∆ which join to
form a double one. Function ∆̃ can always be expressed as ∆̃ = (r − r+)(r++ − r)g(r),
where g(r) is a polynomial function of powers of degree 2k − 1 and no zeroes within the
range [r+, r++] are to be found by construction. Explicitly, equation (36) is
D(ǫ) =
∫ r++2rcǫ
(r − r+)1/2(r+ + 2rcǫ− r)1/2
rk−1/2
g1/2(r)
︸ ︷︷ ︸
. (37)
Now, h(r) is a continuous divergenceless strictly positive function in the compact [r+, r++],
which means that it will reach a positive maximum and minimum for certain r′s. Let us
call hmax and hmin the values of the function h in these points
8. Then
∫ r++2rcǫ
(r − r+)1/2(r+ + 2rcǫ− r)1/2
≤ D(ǫ) ≤
≤ hmax
∫ r++2rcǫ
(r − r+)1/2(r+ + 2rcǫ− r)1/2
. (38)
The integration can be performed:
∫ r++2rcǫ
(r − r+)1/2(r+ + 2rcǫ− r)1/2
= π. (39)
D(ǫ → 0) = πh(rc), (40)
where the value of rc is given in (22).
Integrations of form (39) are solved exactly by a cos−1 type function, and a nonzero
finite result is obtained. It is remarkable that the same can be said for any ∆ we would
choose, as long as no more than two single roots were to join to form a double one. The
key point is that (39), which could be problematic, is independent of ǫ and therefore the
distance is finite in the limit, when ǫ → 0. So, eventhough (39) was neither exactly the
physical distance in the massive case nor in Schwarzschild-de Sitter solution (however, it
was in the massless case as we have already seen in the first section), it is closely related
to it. This fact gives us a hint or, at least, justifies the change of coordinates we were
performing once and again to study the geometry at the limit ǫ → 0.
8These, in principle, depend on ǫ but coincide when ǫ → 0: hmin = hmax ≡ h(rc).
B Horizon coalescence as a flow on the line
The main phenomenon that concerns this paper, as said before, can be described in terms
of the dynamics of a vector field on the line. The coalescence point, in this picture, is no
more than a supercritical Pitchfork bifurcation. Let us remember some general features
of the dynamics of a one-dimensional flow. The equation of a general vector field on the
line can be expressed as:
ẋ = f(x, α) (41)
where f is any real function with real support, the dot means differentiation with respect
to t and α is a parameter of the model. Fixed points of (41) require ẋ = 0, which must
be obtained by finding the roots of f , that is
f(x∗, α) = 0. (42)
Equation (42) is solved by an n-collection of fixed points x∗i for a given value of α. Let
us suppose that f has three roots if α = α0. Fixed points come closer as α moves and
get “condensed” in a “fat” fixed point (bifurcation point) at α = αc. A paradigmatic
example of a Pitchfork bifurcation is shown by function
f(x) = x(α− x2). (43)
One question arises naturally now about the role the horizons play in this picture. Let
us claim that horizons are fixed points and the role of α is played by m. We will justify
this identification by constructing the vector flow.
Constructing a flow in a manifold (in our case it will be a line) is equivalent to giving
a family of curves r̄(t) which covers the manifold or part of it. Each of the curves gets
specified by the initial condition, say, r̄(t = 0). Now, let us consider geodesic motions.
Without loss of generality, the angular coordinates of our geometry will be frozen, θ and
φ are constants, and only radial curves r(t) are to be regarded. Static coordinate system
will serve us to describe the movement for any r ∈ (r+, r++). Let us invoke intuition at
this point. If r(0) is near the cosmological or the black hole horizon it is clear that a
test particle will move out of the region by approaching each horizon respectively. Then,
there is a point r = rg where the test particle will not “feel” any force and, consequently,
it will not move9. This is the first (unstable) fixed point.
Let us move the origin by defining r′ = r− rg, after this, primes will be dropped out
to simplify notation. The flow at each point will be determined by the physical velocity
ṙ(r) (as measured by an observer placed at r = 0) that a test particle would adquire at
r if it is dropped with ṙ = 0 at around r = 0 (as close as possible). It is not hard to see
that the velocity of the test particle, as seen by the static geodesic observer, is bound to
be zero at both horizons. So, horizons are fixed points. Now, our system can be treated
as a vector flow ṙ(r) which covers the region between the horizons. The vector flow has
three fixed points: {r+, r++, 0} where the first two are stable. As m runs towards mc,
the system shrinks into a Pitchfork bifurcation. Near the bifurcation point the flow can
be approximated by
ṙ = βr(r − r+)(r++ − r), (44)
9rg in our geometry, plays the role the asymptotic infinity does in Schwarzschild solution,that is, the
point where the time-like Killing vector should be normalized in order to define the horizon temperature.
Note that rg ≡ rc at the coalescence point, that is, when ǫ = 0.
where β is a positive constant which depends on µ, k and Λ. On the one hand, in the
coordinate system {χ, τ}, and using (30), we have
−→ ṙ =
. (45)
On the other hand, equation (44), expressed in the new coordinate system, reads ṙ =
ǫ3βr3c cosχ sin
2 χ, and so
= −iǫβr3cB cosχ sin2 χ. (46)
As expected, in the new coordinate system, every point converts into a fixed point as
horizons coalesce (ǫ → 0). Since the flux lines were identified with geodesics of test
particles, this can be understood as the abscence of forces at the end of the process.
Acknowledgment
We thank P. K. Townsend and Adil Belhaj for helpful discussions and Jean Nuyts for
critical reading of the manuscript. This work has been supported by MCYT ( Spain)
under grant FPA 2003-02948.
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http://arxiv.org/abs/hep-th/0609055
http://arxiv.org/abs/hep-th/0703255
http://arxiv.org/abs/hep-th/0607193
Introduction
The gravitational coupling. Some geometrical features
The horizon coalecence geometry
Case m=0
Case m=0
Conclusions
Proof of the finite nonzero physical distance
Horizon coalescence as a flow on the line
|
0704.0367 | Instanton representation of Plebanski gravity. Consistency of the
initital value constraints under time evolution | arXiv:0704.0367v6 [gr-qc] 18 Mar 2011
Instanton representation of Plebanski gravity.
Consistency of the initital value constraints under
time evolution
Eyo Eyo Ita III
October 24, 2018
Department of Applied Mathematics and Theoretical Physics
Centre for Mathematical Sciences, University of Cambridge, Wilberforce Road
Cambridge CB3 0WA, United Kingdom
[email protected]
Abstract
The instanton representation of Plebanski gravity provides as equa-
tions of motion a Hodge self-duality condition and a set of ‘general-
ized’ Maxwell’s equations, subject to gravitational degrees of freedom
encoded in the initial value constraints of general relativity. The main
result of the present paper will be to prove that this constraint sur-
face is preserved under time evolution. We carry this out not using
the usual Dirac procedure, but rather the Lagrangian equations of mo-
tion themsleves. Finally, we provide a comparison with the Ashtekar
formulation to place these results into overall context.
http://arxiv.org/abs/0704.0367v6
1 Introduction
In [1] a new formulation of general relativity was presented, named the
instanton representation of Plebanski gravity. The basic dynamical variables
are an SO(3, C) gauge connection Aaµ and a matrix Ψae taking its values in
two copies of SO(3, C).1 The consequences of the associated action IInst
were determined via its equations of motion, which hinge crucially on weak
equalities implied by the the initial value constraints. For these consequences
to be self-consistent, the constraint surface must be preserved for all time
by the evolution equations. The present paper will demonstrate that this
is indeed the case. We will not use the usual Hamiltonian formulation for
totally constrained systems [2], since we will not make use of any canonical
structure implied by IInst. Rather, we will deduce the time evolution of the
dynamical variables directly from the equations of motion of IInst.
Sections 2 and 3 of this paper present the instanton representation ac-
tion and derive the time evolution of the basic variables. Sections 4, 5 and 6
demonstrate that the nondynamical equations, referred to as the diffeomor-
phism, Gauss’ law and Hamiltonian constraints, evolve into combinations of
the same constraint set. The result is that the time derivatives of these con-
straints are weakly equal to zero with no additional constraints generated on
the system. While we do not use the usual Dirac method in this paper, the
result is still that the instanton representation is in a sense Dirac consistent.
We will make this inference clearer by comparison with the Ashtekar vari-
ables in the discussion section. On a final note, the terms ‘diffeomorphism’
and ‘Gauss’ law’ constraints are used loosely in this paper, in that we have
not specified what transformations of the basic variables these constraints
generate. The use of these terms is mainly for notational purposes, due to
their counterparts which appear in the Ashtekar variables.
2 Instanton representation of Plebanski gravity
The starting action for the instanton representation of Plebanski gravity is
given by [1]
IInst =
d3xΨaeB
F a0i + ǫkjmB
−iN(detB)1/2
Λ+ trΨ−1
, (1)
1Index labelling conventions for this paper are that symbols a, b, . . . from the begini-
ing of the Latin alphabet denote internal SO(3, C) indices while those from the middle
i, j, k, . . . denote spatial indices. Both of these sets of indices take takes 1, 2 and 3. The
Greek symbols µ, ν, . . . refer to spacetime indices which take values 0, 1, 2, 3.
where Nµ = (N,N i) are the lapse function and shift vector from metric
general relativity, and Λ is the cosmological constant. The basic fields are
Ψae and A
i , and we action (1) is defined only on configurations restricted
to (detB) 6= 0 and (detΨ) 6= 0.2 In the Dirac procedure one refers to Nµ
as nondynamical fields, since their velocities do not appear in the action.
While the velocity Ψ̇ae also does not appear, we will distinguish this field
from Nµ since the action (1), unlike for the latter, is nonlinear in Ψae.
The equation of motion for the shift vector N i, the analogue of the
Hamilton equation for its conjugate momentum Π ~N , is given by
δIInst
= ǫmjkB
eΨae = (detB)(B
−1)diψd ∼ 0, (2)
where ψd = ǫdaeΨae is the antisymmetric part of Ψae. This is equivalent to
the diffeomorphism constraint Hi owing to the nondegeneracy of B
a, and
we will often use Hi and ψd interchangeably in this paper. The equation of
motion for the lapse function N , the analogue of the Hamilton equation for
its conjugte momentum ΠN , is given by
δIInst
= (detB)1/2
Λ+ trΨ−1
= 0. (3)
Nondegeneracy of Ψae and the magnetic field B
e implies that on-shell, the
following relation must be satisfied
Λ + trΨ−1 = 0, (4)
which we will similarly take as synonymous with the Hamiltonian constraint.
The equation of motion for Ψae is
δIInst
= BkeF
0k + ǫkjmB
m + iN
detΨ(Ψ−1Ψ−1)ea ∼ 0, (5)
up to a term proportional to (4) which we have set weakly equal to zero.
One could attempt to define a momentum conjugate to Ψae, for which (5)
would be the associated Hamilton’s equation of motion. But since Ψae forms
part of the canonical structure of (1), then our interpretation is that this is
not technically correct.3
The equation of motion for the connection Aaµ is given by
δIInst
∼ ǫµσνρDσ(ΨaeF eνρ)−
4ǫmjkN
mBkeΨ[de]
+N(B−1)dj
Λ+ trΨ−1
, (6)
2The latter case limits the application of our results to spacetimes of Petroc Types I,
D and O (See e.g. [3] and [4].
3This is because (5) contains a velocity Ȧa
within F a
and will therefore be regarded
as an evolution equation rather than a constraint. This is in stark contrast with (2) and
(3), which are genuine constraint equations due to the absence of any velocities.
where we have defined
ea(x, y) ≡
δAai (x)
Bje(y) = ǫ
−δae∂k + fedaAdk
δ(3)(x, y); D
ea ≡ 0. (7)
The terms in large round brackets in (6) vanish weakly, since they are pro-
portional to the constraints (2) and (4) and their spatial derivatives. For
the purposes of this paper we will regard (6) as synonymous with
ǫµσνρDσ(ΨaeF
νρ) ∼ 0. (8)
In an abuse of notation, we will treat (5) and (8) as strong equalities in
this paper. This will be justified once we have completed the demonstration
that the constraint surface defined collectively by (2), (3) and the Gauss’
constraint from (8) is indeed preserved under time evolution. As a note
prior to proceeding we will often make the identification
N(detB)1/2
detΨ ≡
−g (9)
as a shorthand notation, to avoid cluttering many of the derivations which
follow in this paper.
2.1 Internal consistency of the equations of motion
Prior to embarking upon the issue of consistency of time evolution of the
initial value constraints, we will check for internal consistency of IInst, which
entails probing of the physical content implied by (8) and (5). First, equation
(8) can be decomposed into its spatial and temporal parts as
Di(ΨbfB
f ) = 0; D0(ΨbfB
f ) = ǫ
ijkDj(ΨbfF
0k). (10)
The first equation of (10) is the Gauss’ law constraint of a SO(3) Yang–
Mills theory, when one makes the identification of ΨbfB
f ∼ Eib with the
Yang–Mills electric field. The Maxwell equations for U(1) gauge theory
with sources (ρ, ~J), in units where c = 1, are given by
~∇ · ~B = 0; Ḃ = −~∇× ~E = 0; ~∇ · ~E = ρ; ~̇E = − ~J + ~∇× ~B. (11)
Equations (10) can be seen as a generalization of the first two equations of
(11) to SO(3) nonabelian gauge theory in flat space when one: (i) identifies
with the SO(3) generalization of the electric field ~E, and (ii) one
chooses Ψae = kδae for some numerical constant k.
When ρ = 0 and ~J = 0, then one has the vacuum theory and equations
(11) are invariant under the transformation
( ~E, ~B) −→ (− ~B, ~E). (12)
Then the second pair of equations of (11) become implied by the first pair.
This is the condition that the Abelian curvature Fµν , where F0i = Ei and
ǫijkFjk = Bi, is Hodge self-dual with respect to the metric of a conformally
flat spacetime. But equations (10) for more general Ψae encode gravitational
degrees of freedom, which as shown in [1] generalizes the concept of self-
duality to more general spacetimes solving the Einstein equations. Let us
first attempt to derive the analogue for (10) of the second pair of (11) in the
vacuum case. Acting on the first equation of (10) with D0 yields
D0Di(ΨbfB
f ) = DiD0(ΨbfB
f ) + [D0,Di](ΨbfB
f ) = 0. (13)
Substituting the second equation of (10) into the first term on the right hand
side of (13) and using the definition of temporal curvature as the commutator
of covariant derivatives on the second term we have
ijkDj(ΨbfF
0k)) + fbcdF
0iΨdfB
f = fbcd
0k +B
Ψdf = 0 (14)
where we have also used the spatial part of the commutator ǫijkDiDjva =
fabcB
b vc. Note that the term in brackets in (14) is symmetric in f and c,
and also forms the symmetric part of the left hand side of (5)
0i + i
−g(Ψ−1Ψ−1)fb + ǫijkBifB
k = 0, (15)
re-written here for completeness. To make progress from (14), we will sub-
stitute (15) into (14). This causes the last term of (15) to drop out due to
antisymmetry, which leaves us with
−gfbcd
Ψdf (Ψ
−1Ψ−1)fc +Ψdf (Ψ
−1Ψ−1)fc
= −2i
−gfbcdΨ−1dc . (16)
The equations are consistent only if (16) vanishes, which is the requirement
that Ψae = Ψea be symmetric. This of course is the requirement that the
diffeomorphism constraint (2) be satisfied. So the analogue of the second
pair of (11) in the vacuum case must be encoded in the requirement that
Ψae = Ψea be symmetric.
3 The time evolution equations
We must now verify that the initial value constraints are preserved under
time evolution defined by the equations of motion (5) and (6). These equa-
tions are respectively the Hodge duality condition
0k + i
−g(Ψ−1Ψ−1)fb + ǫijkN iBjbB
f = 0, (17)
and one of the Bianchi identity-like equations
ǫijkDj(ΨaeF
ok) = D0(ΨaeB
e). (18)
Since the initial value constraints were used to obtain the second line of (17)
from (1), then we must verify that these constraints are preserved under
time evolution as a requirement of consistency. Using F b0i = Ȧ
i −DiAb0 and
defining
−g(B−1)fi (Ψ
−1Ψ−1)fb + ǫmnkN
mBnb ≡ iHbk, (19)
Then equation (17) can be written as a time evolution equation for the
connection, which is not the same as a constraint equation as noted earlier
F b0i = −iHbi −→ Ȧbi = DiAb0 − iHbi . (20)
From equation (20) we can obtain the following equation governing time
evolution equation for the magnetic field
Ḃie = ǫ
ijkDjȦ
k = ǫ
ijkDj
0 − iHek
= febcB
0 − iǫijkDjHek = −δ~θB
e − iǫijkDjHek, (21)
which will be useful. On the first term on the right hand side of (21) we
have used the definition of the curvature as the commutator of covariant
derivatives. The notation δ~θ in (21) suggests that that B
e transforms as
a SO(3, C) vector under gauge transformations parametrized by θb ≡ Ab0.4
Since we have not specified anything about the canonical structure of IInst,
then δ~θ as used in (21) and in (24) should at this stage simply be regarded
as a definition useful for shorthand notation.
We will now apply the Liebnitz rule in conjunction with the definition
of the temporal covariant derivatives to (18) to determine the equation gov-
erning the time evolution of Ψae. This is given by
D0(ΨaeB
e) = B
eΨ̇ae +ΨaeḂ
e + fabcA
0(ΨceB
e) = ǫ
ijkDj(ΨaeF
0k). (22)
Substituting (21) and (20) into the left and right hand sides of (22), we have
BieΨ̇ae +Ψae
febcB
0 − iǫijkDjHek
+ fabcA
0(ΨceB
e) = −iǫijkDj(ΨaeHek).(23)
In what follows, it will be convenient to use the following transformation
properties for Ψae as A
i under SO(3, C) gauge transformations
δ~θΨae =
fabcΨce + febcΨac
Ab0; δ~θA
i = −DiAa0; δ~θB
e = −febcBibAc0. (24)
4We will make the identification with SO(3, C) gauge transformations later in this
paper when we bring in the relation of IInst with the Ashtekar variables.
Then using (24), the time evolution equations for the phase space variables
ΩInst can be written in the following compact form
Ȧbi = −δ~θA
i − iHbi ; Ψ̇ae = −δ~θΨae − iǫ
ijk(B−1)ei (DjΨaf )H
k . (25)
We have found evolution equations for Ψae and A
i from the covariant equa-
tions of Aaµ and the Hodge-duality condition We have obtained these without
using Poisson brackets, and by assuming that the Hamiltonian and diffeo-
morphism constraints are satisfied. Therefore the first order of business is
then to check for the preservation of the initial value constraints under the
time evolution generated by (25). This means that we must check that
the time evolution of the diffeomorphism, Gauss’ law and Hamiltonian con-
straints are combinations of terms proportional to the same constraints and
their spatial derivatives, and terms which vanish when the constraints hold.5
These constraints are given by
we{Ψae} = 0; (detB)(B−1)diψd = 0; (detB)1/2
Λ+ trΨ−1
= 0(26)
where (detB) 6= 0 and (detΨ) 6= 0. We will occasionally make the identifi-
cation
N(detB)1/2(detΨ)1/2 ≡
−g (27)
for a shorthand notation. Additionally, the following definitions are provided
for the vector fields appearing in the Gauss’ constraint
we = B
eDi; ve = B
e∂i (28)
where Di is the SO(3, C) covariant derivative with respect to the connection
Aai . Equations (26) are the equations of motion for the auxilliary fields A
N i and N .
4 Consistency of the diffeomorphism constraint un-
der time evolution
The diffeomorphism constraint is directly proportional to ψd = ǫdaeΨae, the
antisymmetric part of Ψae. So to establish the consistency condition for
this constraint, it suffices to show that the antisymmetric part of the second
equation of (25) weakly vanishes. This is given by
ǫdaeΨ̇ae = −δ~θ(ǫdaeΨae)− iǫdaeǫ
ijk(B−1)ei (DjΨaf )H
k , (29)
5This includes any nonlinear function of linear order or higher in the constraints, a
situation which involves the diffeomorphism constraint.
which splits into two terms. Using (24), one finds that the first term of (29)
is given by
−ǫdaeδ~θΨae = −ǫdae
fabcΨce +Ψacfebc
δebδdc − δecδbd
Ψce +
δdbδac − δdcδab
Ψdb − δbdtrΨ + δdbtrΨ−Ψbd
Ab0 = 2Ψ[bd]A
0 = −ǫdbhAb0ψh, (30)
which is proportional to the diffeomorphism constraint. The second term of
(30) has two contributions due to H
k as defined in (19). The first contribu-
tion reduces to
−iǫdaeǫijk(B−1)ei (DjΨaf )(H(1))
= −iǫdaeǫijk(B−1)ei (DjΨaf )
−g(B−1)gk(Ψ
−1Ψ−1)gf
= iǫdae(detB)
−1ǫegh(Ψ−1Ψ−1)gfB
hDjΨaf
= i(detB)−1(Ψ−1Ψ−1)gf
a − δgaδhd
vv{Ψaf}
= i(detB)−1(Ψ−1Ψ−1)gf
dva{Ψaf} − vd{Ψgf}
= i(detB)−1
(Ψ−1Ψ−1)dfGf + vd{Λ+ trΨ−1}
. (31)
The first term on the final right hand side of (31) is the Gauss’ constraint
and the second term is the derivative of a term direction proportional to
the Hamiltonian constraint.6 The second contribution to the second term
of (29) is given by
ǫdaeǫ
ijk(B−1)ei (DjΨaf )(H(2))
k = ǫdaeǫ
ijk(B−1)ei (DjΨaf )ǫmnkN
= ǫdac
n − δinδjm
(B−1)ei (DjΨaf )N
= ǫdaeN
i(B−1)eivf{Ψaf} −N jDj(ǫdaeΨae) = ǫdaeN i(B−1)efGa −N jDjψd.(32)
The result is that the time evolution of the diffeomorphism constraint is
directly proportional to
ψ̇d =
i(detB)−1(Ψ−1Ψ−1)da + ǫdaeN
i(B−1)ei
Ab0ǫbdh − δdhN jDj
ψh + i(detB)
vd{(−g)−1/2H}, (33)
which is a linear combination of terms proportional to the constraints (26)
and their spatial derivatives. The result is that the diffeomorphism con-
straint Hi = 0 is consistent with respect to the Hamiltonian evolution gen-
erated by the equations (25). So it remains to verify consistency of Gauss’
law and the Hamiltonian constraints Ga and H.
6We have added in a term Λ, which can be regarded as a constant of integration with
respect to the spatial derivatives from vd.
5 Consistency of the Gauss’ constraint under time
evolution
Having verified the consistency of the diffeomorphism constraint under time
evolution, we now move on to the Gauss’ constraint. Application of the
Liebnitz rule to the first equation of (26) yields
Ġa = Ḃ
eDiΨae +B
eDiΨ̇ae +B
fabfΨfe + febgΨag
Ȧai . (34)
Upon substituion of (21) and (25) into (34), we have
Ġa =
−δ~θB
e − iǫijkDjHek
DiΨae +B
−δ~θΨae − iǫ
ijk(B−1)ei (DjΨaf )H
fabfΨfe + febgΨag
−δ~θA
i − iHbi
.(35)
Using the Liebniz rule to combine the δ~θ terms of (35), we have
Ġa = −δ~θGa − iǫ
k)DiΨae +B
e Dm((B
−1)ei (DjΨaf )H
fabfΨfe + febgΨag
i . (36)
The requirement of consistency is that we must show that the right hand
side of (36) vanishes weakly. First, we will show that the third term on the
right hand side of (36) vanishes up to terms of linear order and higher in
the diffeomorphism constraint. This term, up to an insignificant numerical
factor, has two contributions. The first contribution is
fabfΨfe + febgΨag
Bie(H(1))
fabfΨfe + febgΨag
(Ψ−1Ψ−1)eb
fabf (Ψ
−1)fb + febg(Ψ
−1Ψ−1)ebΨag
∼ δ(1)a (~ψ) ∼ 0, (37)
which is directly proportional to a nonlinear function of first order in ψd
which is proportional to the diffeomorphism constraint. The second contri-
bution to the third term on the right hand side of (36) is
fabfΨfe + febgΨag
Bie(H(2))
fabfΨfe + febgΨag
ǫkmnN
kBme B
fabfΨfe + febgΨag
(detB)Nk(B−1)dkǫdeb
= (detB)Nk(B−1)dk
δfdδae − δfeδad
Ψfe + 2δdgΨag
= (detB)Nk(B−1)dk
Ψda − δadtrΨ + 2Ψad
≡ δ(2)a ( ~N) (38)
which does not vanish, and neither is it expressible as a constraint. For
the Gauss’ law constraint to be consistent under time evolution, a necessary
condition is that this δ
a ( ~N) term must be exactly cancelled by another
term arising from the variation.
Let us expand the terms in square brackets in (36). This is given, using
the Liebniz rule on the second term, by
ǫijk(DjH
k)(DiΨae) + ǫ
ijkBme Dm((B
−1)ei (DjΨae)H
= ǫijk(DjH
k)(DiΨae)− ǫijkBme (B−1)en(DmBng )(B−1)
i (DjΨaf )H
+ǫmjk(DmDjΨaf )H
+ ǫmjk(DjΨaf )(DmH
). (39)
The first and last terms on the right hand side of (39) cancel, which can be
seen by relabelling of indices. Upon application of the definition of curvature
as the commutator of covariant derivatives to the third term, then (39)
reduces to
−ǫijk(DnBng )(B−1)
i (DjΨaf )H
fabcΨcf + ffbcΨac
. (40)
The first term of (40) vanishes on account of the Bianchi identity and the
second term contains two contributions which we must evaluate. The first
contribution is given by
(H(2))
fabcΨcf + ffbcΨac
= (detB)Nk(B−1)dkǫdbf
fabcΨcf + ffbcΨac
= (detB)Nk(B−1)dk
δdaδfc − δdcδfa
Ψcf − 2δdcΨac
= (detB)Nk(B−1)dk
δdatrΨ−Ψda − 2Ψad
= −δ(2)a ( ~N ),(41)
with δ
a ( ~N) as given in (37). So putting the results of (39), (40) and (41)
into (36), we have
Ġa = −δ~θGa + δ
~N) + δ(1)a (
~ψ) + δ(1)a (
~ψ)− δ(2)a ( ~N) = −δ~θGa + 2δ
(1)(~ψ).(42)
The velocity of the Gauss’ law constraint is a linear combination of the
Gauss’ constraint with terms of the diffeomorphism constraint of linear or-
der and higher. Hence the time evolution of the Gauss’ law constraint is con-
sistent in the sense that we have defined, since δ(1)(~ψ) vanishes for ψd = 0.
6 Consistency of the Hamiltonian constraint un-
der time evolution
The time derivative of the Hamiltonian constraint, the third equation of
(26), is given by
((detB)1/2(detΨ)1/2
(Λ + trΨ−1) +
(Λ + trΨ−1) (43)
which has split up into two terms. The first term is directly proportional
to the Hamiltonian constraint, therefore it is already consistent. We will
nevertheless expand it using (21) and (25)
(B−1)di Ḃ
d + (Ψ
−1)aeΨ̇ae
(detB)1/2(detΨ)1/2(Λ + trΨ−1)
(B−1)di
−δ~θB
d − iǫijkDjHdk
+(Ψ−1)ae
−δ~θΨae − iǫ
ijk(B−1)ei (DjΨaf )H
H. (44)
We will be content to compute the δ~θ terms of (44). These are
(B−1)di δ~θB
d = (B
−1)di fdbfB
0 = δdbfdbfA
0 = 0 (45)
on account of antisymmetry of the structure constants, and
(Ψ−1)eaδ~θΨae = (Ψ
−1)ea
fabfΨfe + febgΨag
= 0, (46)
also due to antisymmetry of the structure constants. We have shown that
the first term on the right hand side of (43) is consistent with respect to
time evolution. To verify consistency of the Hamiltonian constraint under
time evolution, it remains to show that the second term is weakly equal to
zero. It suffices to show this just for the second term, in brackets, of (43)
(Λ + trΨ−1) = −(Ψ−1Ψ−1)feΨ̇ef
= (Ψ−1Ψ−1)ef
δ~θΨae − iǫ
ijk(B−1)ei (DjΨaf )H
, (47)
where we have used (25). Equation (47) has split up into two terms, of
which the first term is
(Ψ−1Ψ−1)eaδ~θΨae = (Ψ
−1Ψ−1)ea
fAbfΨfe + febgΨag
fabf (Ψ
−1)fa + febg(Ψ
−1)eg
Ab0 = m(
~ψ) ∼ 0 (48)
which vanishes weakly since it is a nonlinear function of at least linear order
in ψd. The second term of (47) splits into two terms which we must evaluate.
The first contribution is proportional to
(Ψ−1Ψ−1)eaǫijk(B−1)ei (DjΨaf )(H(1))
−g(Ψ−1Ψ−1)eaǫijk(B−1)ei (DjΨaf )(B−1)dk(Ψ−1Ψ−1)df
−g(Ψ−1Ψ−1)ea(Ψ−1Ψ−1)df (detB)−1ǫedgBjgDjΨaf
−g(detB)−1ǫedg(Ψ−1Ψ−1)ea(Ψ−1Ψ−1)dfvg{Ψaf} ≡ v{~ψ} (49)
for some vector field v. We have used the fact that the term in (49) quartic
in Ψ−1 in antisymmetric in a and f due to the epsilon symbol. Hence Ψaf
as acted upon by vg can only appear in an antisymmetric combination, and
is therefore proportional to the diffeomorphism constraint ψd whose spatial
derivatives weakly vanish. Hence (49) presents a consistent contribution to
the time evolution of H, which leaves remaining the second contribution to
the second term of (47). This term is proportional to
(Ψ−1Ψ−1)eaǫijk(B−1)ei (DjΨaf )(H(2))
= (Ψ−1Ψ−1)eaǫijk(B−1)ei (DjΨaf )ǫmnkN
n − δinδjm
(B−1)eiB
−1Ψ−1)ea(DjΨaf )
N i(B−1)eiB
f − δefN
(Ψ−1Ψ−1)ea(DjΨaf )
= (−g)−1/2N iHai vf{Ψaf} − (Ψ−1Ψ−1)fa(N jDjΨaf )
= (−1)−1/2N iHai Ga −N jDj(Λ + trΨ−1). (50)
The first term on the final right hand side of (50) is proportional to the
Gauss’ law constraint, and the second term is proportional to the derivative
of the Hamiltonian constraint. To obtain this second term we have added
in Λ as a constant of differentiation with respect to ∂j . Substituting (48),
(49) and (50) into (47), then we have
Ḣ =∼ Ô(~ψ) + (−g)−1/2N iHai Ga + T̂ ((−g)−1/2H), (51)
where Ô and T̂ are operators consisting of spatial derivatives acting to the
right and c numbers. The time derivative of the Hamiltonian constraint
is a linear combination of the Gauss’ law and Hamiltonian constraints and
its spatial derivatives, plus terms of linear order and higher in the diffeo-
morphism constraint and its spatial derivatives. Hence the Hamiltonian
constraint is consistent under time evolution.
7 Recapitulation
The final equations governing the time evolution of the initial value con-
straints are given weakly by
ψ̇d =
i(detB)−1(Ψ−1Ψ−1)da + ǫdaeN
i(B−1)ei
Ab0ǫbdh − δdhN jDj
ψh + i(detB)
vd{Λ + trΨ−1};
Ġa = −fabcAb0Gc + δ(1)a (~ψ);
ǫijk(B−1)di (DjH
k ) + ǫ
ijk(B−1)ei (Ψ
−1)ae(DjΨaf )H
−N j∂j
(Λ + trΨ−1)
+(−g)−1/2N iHai Ga −
−g(detB)−1ǫedg(Ψ−2Ψ−1)ea(Ψ−1Ψ−1)dfvg{ǫafhψh}+m(~ψ).(52)
Equations (52) show that all constraints derivable from the the action (1)
are preserved under time evolution, since their time derivatives yield linear
combinations of the same set of constraints and their spatial derivatives.
There are no additional constraints generated which implies that the action
(1) is consistent in the Dirac sense. On the other hand, we have not defined
the canonical structure of (52) or any Poisson brackets.
Equations (52) can be written schematically in the following form
~̇H ∼ ~H + ~G+H; ~̇G ∼ ~G+Φ( ~H); Ḣ ∼ H + ~G+Φ( ~H), (53)
where Φ is some nonlinear function of the diffeomorphism constraint ~H,
which is of at least first order in ~H. In the Hamiltonian formulation of a
theory, one identifies time derivatives of a variable f with via ḟ = {f,H}
the Poisson brackets of the variable with the Hamiltonian H . So while we
have not specified Poisson brackets, equation (53) implies the existence of
Poisson brackets associated to some Hamiltonian HInst for the action (1),
{ ~H,HInst} ∼ ~H + ~G+H; {~G,HInst} ∼ ~G+Φ( ~H);
{H,HInst} ∼ H +Φ( ~H) + ~G. (54)
So the main result of this paper has been to demonstrate that the instanton
representation of Plebanski gravity forms a consistent system, in the sense
that the constraint surface is preserved under time evolution. As a direction
of future research we will compute the algebra of constraints for (1) directly
from its canonical structure. Nevertheless it will be useful for the present
paper to think of equations (52) in the Dirac context, mainly for compari-
son with other formulations of general relativity. This will bring us to the
Ashtekar variables.
8 Discussion: Relation of the instanton represen-
tation to the Ashtekar variables
We will now provide the rationale for not following the Dirac procedure for
constrained systems [2] with respect to (1), by comparison with the Ashekar
formulation of GR. The action for the instanton representation (1) can be
written in the following 3+1 decomposed form
IInst =
0we{Ψae} − ǫijkN iBjaBkeΨae
−iN(detB)1/2(detΨ)1/2
Λ+ trΨ−1
, (55)
which regards Ψae and A
i as phase space variables. But the phase space of
(55) is noncanonical since its symplectic two form
δθInst = δ
d3xΨaeB
d3xBieδΨae ∧ δAai +
d3xΨaeǫ
ijkDj(δA
k) ∧ δAai , (56)
is not closed owing to the presence of the second term on the right hand
side. The initial stages of the Dirac procedure applied to (55) state that the
momentum conjugate to Aai yields the primary constraint
Πia =
δIInst
δȦai
= ΨaeB
e. (57)
Then making the identification σ̃ia = Π
a and upon substitution into (57)
and into (55), one obtains the action
IAsh =
σ̃iaȦ
0Ga −N iHi −
, (58)
which is the action for the Ashtekar complex formalism of general relativity
[5], [6], with σ̃ia being the densitized triad. This is a totally constrained sys-
tem with (Aa0, N
i, N), respectively the SO(3, C) rotation angle Aa0, the shift
vector N i and the densitized lapse function N = N(detσ̃)−1/2 as auxilliary
fields. The constraints in (58) smearing the auxilliary fields are the Gauss’
law, vector and Hamiltonian constraints
Ga = Diσ̃
a; Hi = ǫijkσ̃
a ; H = ǫijkǫ
abcσ̃iaσ̃
σ̃kc +B
. (59)
From (58) one reads off the symplectic two form ΩAsh given by
ΩAsh =
d3xδσ̃ia ∧ δAai = δ
d3xσ̃iaδA
= δθAsh, (60)
which is the exact functional variation of the canonical one form θAsh.
The actions (55) and (58) are transformable into each other only under
the condition (detB) 6= 0 and (detΨ) 6= 0. In (58) it is clear that σ̃ia and Aai
form a canonically conjugate pair, which suggests that (55) is a noncanonical
version of (58). The constraints algebra for (59) is
{ ~H[ ~N ], ~H [ ~M ]} = Hk
N i∂kMi −M i∂kNi
{ ~H[N ], Ga[θa]} = Ga[N i∂iθa];
{Ga[θa], Gb[λb]} = Ga
fabcθ
{H(N), ~H [ ~N ]} = H[N i∂iN
{H(N ), Ga(θa)} = 0;[
H(N),H(M )
= Hi[
N∂jM −M∂jN
H ij], (61)
which is first class due to closure of the algebra, and is therefore consistent
in the Dirac sense. Let us consider (61) for each constraint with the total
Hamiltonian HAsh and compare with (54). This is given schematically by
{ ~H,HAsh} ∼ ~H + ~G+H; {~G,HAsh} ∼ ~G+ ~H;
{H,HAsh} ∼ H + ~H. (62)
Comparison of (62) with (54) shows an essentially similar structure for the
top two lines involving ~H and ~G.7 But there is a marked dissimilarity with
respect to the Hamiltonian constraint H. Note that there is a Gauss’ law
constraint appearing in the right hand side of the last line of (54) whereas
there is no such constraint on the corresponding right hand side of (62).
This means that while the Hamiltonian constraint is gauge-invariant under
SO(3, C) gauge-transformations as implied by (61) and (62), this is not
the case in (54). This means that the action (1), which as shown in [1]
describes general relativity for Petrov Types I, D and O, has a different
role for the Gauss’ law and Hamiltonian constraints than the action (58),
which also describes general relativity. Therefore IInst and IAsh at some
level correspond to genuinely different descriptions of GR, a feature which
would have been missed had we applied the step-by-step Dirac procedure.
9 Appendix: Commutation relations for IInst
We will now infer the Poisson brackets for (55) by inference from the corre-
sponding canonical Ashtekar Poisson brackets
{Aai (x), σ̃
b (y)} = δ
(3)(x,y) (63)
along with the vanishing brackets
{Aai (x), Abj(y)} = {σ̃ia(x), σ̃
(x)} = 0. (64)
To find the analogue of (63) and (64) for (55), we will use the tranformation
equation
σ̃ia = ΨaeB
e, (65)
which corresponds to a noncanonical transformation. Substitution of (65)
into (63) yields
{Aai (x),Ψbf (y)Bif (y)} = δ
(3)(x,y)
{Aai (x),Ψbf (y)}B
(y) + Ψbf (x){Aai (x), B
(y)}. (66)
The second term on the right hand side of (66) vanishes on account of the
first relations of (64), and upon multiplying (66) by the inverse magnetic
field (B−1)ei , assumed to be nondegenerate, we obtain
{Aai (x),Ψbf (y)} = δab (B−1(y))
(3)(x,y). (67)
7The linearly versus nonlinearly of the diffeomorphism constraints on the right hand
side is just a minor difference.
This gives us the Poisson brackets {A,A} ∼ 0 and {A,Ψ} ∼ B−1, which
leaves remaining the brackets {Ψ,Ψ}. To obtain these, we substitute (65)
into the second equation of (64), yielding
{σ̃ia(x), σ̃bj(y)} = {Ψae(x)Bie(x),Ψbf (y)B
f (y)}
= Ψae(x){Bie(x),Ψbf (y)}B
(y) + {Ψae(x),Ψbf (y)}Bie(x)B
+Ψbf (x)Ψae(x){Bie(x), B
f (y)}+Ψbf (y){Ψae(x), B
f (y)}B
e(x) = 0. (68)
Noting that the third term vanishes on account of the first equation of (64),
equation (68) reduces to
{Ψae(x),Ψbf (y)}Bie(x)B
f (y)
+Ψae(x){Bie(x),Ψbf (y)}B
(y)−Ψbf (y){Bjf (y),Ψae(x)}B
e(x) = 0. (69)
The bottom two terms of (69) can be computed using (67)
{Bie(x),Ψbf (y)} = ǫimnDxm{Aen(x),Ψbf (y)} = ǫimnDxm(δeb (B−1(y))fnδ(3)(x,y)).(70)
Substituting (70) into (69) and cancelling a pair of magnetic fields, then we
have that
{Ψae(x),Ψbf (y)}Bie(x)B
f (y) = ǫ
Ψae(x)D
m +Ψba(y)D
δ(3)(x,y). (71)
Left and right multiplying (71) by the inverse of the magnetic fields, we have
{Ψae(x),Ψbf (y)} = ǫijm
(B−1(y))
mΨab(x)(B
−1(x))ei
+(B−1(x))eiD
mΨba(y)(B
−1(y))
δ(3)(x,y). (72)
One sees that the internal components of Ψae have nontrivial commutation
relations with themselves.
References
[1] Eyo Ita ‘Instanton representation of Plebanski gravity. Gravitational
instantons from the classical formulation.’ arXiv: gr-qc/0703057
[2] Paul Dirac ‘Lectures on quantum mechanics’ Yeshiva University Press,
New York, 1964
[3] Hans Stephani, Dietrich Kramer, Maclcolm MacCallum, Cornelius
Hoenselaers, and Eduard Herlt ‘Exact Solutions of Einstein’s Field
Equations’ Cambridge University Press
[4] R. Penrose and W. Rindler ‘Spinors and space-time’ Cambridge Mono-
graphs in Mathematical Physics
[5] Ahbay Ashtekar ‘New Hamiltonian formulation of general relativity’
Phys. Rev. D36(1987)1587
[6] Ahbay Ashtekar ‘New variables for classical and quantum gravity’ Phys.
Rev. Lett. Volume 57, number 18 (1986)
|
0704.0368 | Metal-insulator transition in the low-dimensional organic conductor
(TMTSF)2FSO3 probed by infrared microspectroscopy | EPJ manuscript No.
(will be inserted by the editor)
Metal-insulator transition in the low-dimensional organic
conductor (TMTSF)2FSO3 probed by infrared microspectroscopy
A. Pashkin1a, K. Thirunavukkuarasu1, Y.-L. Mathis2, W. Kang3, and C. A. Kuntscher1b
1 Experimentalphysik II, Universität Augsburg, 86159 Augsburg, Germany
2 Institute for Synchrotron Radiation, Forschungszentrum Karlsruhe, P.O. Box 3640, 76021 Karlsruhe, Germany
3 Department of Physics, Ewha Womans University, Seoul 120-750, Korea
Received: October 28, 2018
Abstract. We present measurements of the infrared response of the quasi-one-dimensional organic conduc-
tor (TMTSF)2FSO3 along (E‖a) and perpendicular (E‖b
′) to the stacking axis as a function of temper-
ature. Above the metal-insulator transition related to the anion ordering the optical conductivity spectra
show a Drude-like response. Below the transition an energy gap of about 1500 cm−1 (185 meV) opens,
leading to the corresponding charge transfer band in the optical conductivity spectra. The analysis of the
infrared-active vibrations gives evidence for the long-range crystal structure modulation below the transi-
tion temperature and for the short-range order fluctuations of the lattice modulation above the transition
temperature. Also we report about a new infrared mode at around 710 cm−1 with a peculiar temperature
behavior, which has so far not been observed in any other (TMTSF)2X salt showing a metal-insulator
transition. A qualitative model based on the coupling between the TMTSF molecule vibration and the
reorientation of electrical dipole moment of the FSO3 anion is proposed, in order to explain the anomalous
behavior of the new mode.
PACS. 71.30.+h Metal-insulator transitions and other electronic transitions – 74.70.Kn Organic super-
conductors
1 Introduction
The organic Bechgaard salts (TMTSF)2X consist of
stacks of planar TMTSF (tetramethyltetraselenafulva-
lene) molecules separated by anions (X = PF6, AsF6,
ClO4, Br, etc.). The charge transport in these systems
is restricted to the direction along the molecular stacks,
making the Bechgaard salts prime examples of one-
dimensional metals. However, on cooling down most of
them undergo a metal-insulator transition which prevents
the onset of a superconducting state [1]. In Bechgaard
salts with noncentrosymmetric anions such as ReO4, BF4
or FSO3 the metal-insulator transition is related to the
anion ordering [2]. It was furthermore demonstrated that
in some cases the metal-insulator transition can be sup-
pressed by the application of external pressure, leading to
a superconducting ground state [3].
The case of the anions X=FSO3 in this class of ma-
terials is particularly interesting, since these anions are
noncentrosymmetric and in addition possess a permanent
electrical dipole moment. The first study of the basic prop-
erties of (TMTSF)2FSO3 has been reported by Wudl et
al. in 1982 [4]. Further studies have shown that this com-
a email: [email protected]
b email: [email protected]
pound has the highest superconducting transition temper-
ature (2.5 K at 8.5 kbar) among the Bechgaard salts. It
was proposed that this is due to the interaction of the
conducting electrons with the FSO3 anion dipoles [5]. A
recent detailed study [6] revealed a very rich pressure-
temperature phase diagram of (TMTSF)2FSO3 with a va-
riety of different phases, which have not been completely
identified up to now. Furthermore, by magnetoresistance
measurements a two-dimensional electronic behavior was
found in (TMTSF)2FSO3 under a pressure of around 6.2
kbar [7].
The interaction of the FSO3 anions with each other via
long-range Coulomb forces and with the centrosymmetric
surrounding formed by the TMTSF cations tends to order
the anions below a certain temperature. The first-order
structural phase transition related to this anion ordering
occurs at around TMI=89 K in (TMTSF)2FSO3 at ambi-
ent pressure. The change of the crystal structure modifies
the electronic band structure: The effective half-filled con-
ducting band splits into one filled and one empty band
separated by an energy gap, leading to a sharp metal-
insulator transition [5]. The structural analysis suggested
a modulation of the crystal structure with wavevector q =
(1/2, 1/2, 1/2) below the phase transition, which implies
an antiferroelectric state [8,2]. The ordering of the FSO3
anions modulates the lattice resulting in a new unit cell
http://arxiv.org/abs/0704.0368v1
2 A. Pashkin et al.: Metal-insulator transition in (TMTSF)2FSO3 probed by infrared spectroscopy
of size 2a × 2b × 2c. Thus, there are eight formula units
of (TMTSF)2FSO3 per unit cell in the low temperature
phase. Correspondingly, one can expect a splitting of each
vibrational mode into up to eight components [9].
The ratio of the energy gap to the transition tem-
perature in (TMTSF)2FSO3 is ∼ 12.5 [4], which is ap-
preciably higher than the value 3.5 predicted by the
mean-field theory for the Peierls transition. Therefore,
the metal-insulator in the Bechgaard salts with non-
centrosymmetric anions was attributed to a special type of
Peierls instability which originates from the anion-electron
coupling [10].
In this work we present the results of a temperature-
dependent polarized infrared reflectivity study of
(TMTSF)2FSO3 single crystals in the far- and mid-
infrared frequency range, in order to characterize the
change of electronic and vibrational properties during
the metal-insulator transition at TMI=89 K. This is the
first infrared spectroscopic investigation of the compound
(TMTSF)2FSO3. Our results allow a direct determination
of the charge gap in the insulating state. Furthermore, we
determined and analyzed the behavior of the vibrational
modes during the metal-insulator transition, which can
clarify details of the dipolar ordering.
2 Experimental
(TMTSF)2FSO3 single crystals were grown by standard
electrochemical techniques from TMTSF molecules and
tetrabutylammonium-FSO3. The studied samples have a
needle-like shape, with a size of approximately 2 × 0.2 ×
0.1 mm3. The samples were mounted on a cold-finger Cry-
oVac Konti-Mikro cryostat. The actual measuring temper-
ature was controlled by a sensor attached in direct vicin-
ity of the sample. The measurements were performed at
the infrared beamline of the synchrotron radiation source
ANKA. The polarized infrared reflectivity was measured
in the range 150 - 10000 cm−1 using a Bruker IRscope
II microscope attached to a IFS66v/S spectrometer. The
frequency resolution was 1 cm−1 for all measured spec-
tra. Optically transparent TPX and KBr cryostat win-
dows were used for the measurements in the far- and mid-
infrared frequency range, respectively.
3 Results and discussion
3.1 Electronic properties
The reflectivity spectra of (TMTSF)2FSO3 above and be-
low the metal-insulator transition temperature of 89 K for
both polarizationsE‖a and E‖b′ (along and perpendicular
to the stacking axis, respectively) are shown in Fig. 1. The
reflectivity data in the spectral region at around 450 cm−1
are affected by the absorption features of the far-infrared
TPX cryostat window and are therefore not shown.
At 290 K the reflectivity of the sample along the
stacking axis E‖a demonstrates a typical Drude behavior
E || a
290 K
45 K
8000250
E || b'
Frequency (cm
290 K
45 K
Energy (meV)
E || a
Fig. 1. Reflectivity spectra of (TMTSF)2FSO3 above and be-
low the metal-insulator transition for E‖a and E‖b′.
(growth up to 1 when frequency tends to zero). In con-
trast, at 45 K , i.e., below TMI , the reflectivity is almost
frequency independent below 1000 cm−1, which is typical
for an insulating state.
The interference fringes observed below 400 cm−1 in
the spectra for both polarizations are due to the partial
transparency of the sample in the insulating phase. Per-
pendicular to the stacking axis (E‖b′), the optical reflec-
tivity and conductivity is much lower than along the a
axis. Nevertheless, the observed changes during the metal-
insulator transition are similar to those of the E‖a direc-
tion. These results demonstrate the opening of an energy
gap at the Fermi level for both studied directions.
The dramatic effect of the temperature decrease on the
electronic properties of (TMTSF)2FSO3 are more directly
seen in the optical conductivity spectra. The E‖a optical
conductivity σ1(ω) of (TMTSF)2FSO3 in the insulating
(at 45 K) and conducting phase (at 290 K) obtained by
means of Kramers-Kronig analysis is shown in Fig. 2. The
dominating feature of the spectrum at 45 K is a strong
charge transfer band due to electronic transitions across
the gap. The arrow shows the band gap (1500 cm−1) ob-
tained from the published temperature-dependent dc re-
sistivity measurements [4]. Obviously, the agreement of
this value with the onset of the optical interband tran-
A. Pashkin et al.: Metal-insulator transition in (TMTSF)2FSO3 probed by infrared spectroscopy 3
1000 10000
100 1000
2D = 1500 cm
Frequency (cm
290 K
45 K
E || a
Energy (meV)
Fig. 2. E‖a optical conductivity spectra of (TMTSF)2FSO3
above and below the metal-insulator transition at TMI= 89 K.
Hatched area depicts the Drude model fit of the high temper-
ature optical conductivity.
sition is very good. On the other hand, the optical con-
ductivity at room temperature is mostly dominated by
the Drude response of the free carriers. The corresponding
fit using the Drude model is shown as the hatched area
in Fig. 2. Obviously, the Drude model provides a good
description of the measured room-temperature spectrum
excluding the electron-molecular vibration (emv) antires-
onance modes. The plasma frequency ωp = 8660 cm
and the scattering rate Γ ≃ 1450 cm−1 obtained from the
fit agree well with the Drude model parameters reported
for other TMTSF salts [11]. The obtained value of the dc
conductivity, σdc ≃ 860 (Ωcm)
−1, is in reasonable agree-
ment with the dc and microwave conductivity values of
1600 and 300 (Ωcm)−1, respectively, reported by Wudl et
al. [4].
3.2 Vibrational modes
The TMTSF molecule with the point group symmetry
D2h has in total 72 local vibrational modes classified ac-
cording to the following representations [12]
ΓD2h = (12ag + 11b3g + 11b1u + 11b2u)
+(6b1g + 7b2g + 7au + 7b3u), (1)
where the vibrations in the first brackets are polarized in
the molecular plane (perpendicular to the stacking a axis)
and the vibrations in the second brackets are polarized out
of the plane (along the stacking axis a). The symmetric
(gerade) vibrations are Raman active and the asymmet-
ric (ungerade) vibrations are infrared active excluding the
au silent modes. Some of the totally symmetric ag Ra-
man modes are expected to appear in the infrared spectra
for E‖a due to efficient emv coupling in the modulated
stacking structure [11,13].
Table 1. The eigenfrequencies and assignment of some vi-
brational modes observed in (TMTSF)2FSO3 for E‖a at 45 K
below TMI . All numbers are in cm
45 K calculated frequency1 assignment
580 571 ν3(a1) FSO3
728 702 ν51(b2u)
902, 911, 915, 916 ν8(ag)
917, 924, 932
1020, 1031, 1036 1060 ν7(ag)
1067, 1072 1060 ν7(ag)
1362, 14502 1469 ν4(ag)
1354, 1364, 1369 1369 ν6(ag)
1373, 1379, 1385
1550, 1584, 1606 1596 ν3(ag)
1847, 1854, 1863 1863 ν3(ag) + ν11(ag)
The tetrahedral FSO3 anion has C3v point group sym-
metry which gives in total nine vibrational modes
ΓC3v = 3a1(z, x
2 + y2, z2) + 3e(x, y, x2 − y2, xy, yz, xz),
where e species correspond to the doublets. Thus, in the
infrared spectra one expects six modes, with the 3a1 and
3e modes being polarized along and perpendicular to the
polar axis of the anion, respectively.
In this section we want to concentrate on the changes
in the infrared phonon spectra for both polarizations
across the metal-insulator transition. For E‖a several ag
vibrations of the TMTSF molecules become infrared ac-
tive in the insulating phase. This is due to the effec-
tive emv coupling of these vibrations to the on-chain
charge transfer band in the structure modulated due to
the anion ordering. The list of the new modes observed
below the transition together with their tentative as-
signment is given in Table 1. Most of them are emv
coupled ag modes polarized in the molecular plane or
their combination as a triplet at around 1850 cm−1. The
ν4(ag) mode involving the central C=C bond stretching is
known to have especially strong emv coupling and there-
fore it appears as a strong antiresonance mode in the
optical conductivity spectrum. It should be pointed out
that the observed appearance of ag modes for E‖a in
the ordered phase is typical only for (TMTSF)2X com-
pounds with non-centrosymmetric anions. In comparison,
(TMTTF)2X salts possess a stronger stack dimerization,
resulting in the emv coupling of the ag modes already
in the disordered phase, and therefore the anion ordering
transition causes only a frequency shift and an intensity
change of the emv coupled modes [15].
4 A. Pashkin et al.: Metal-insulator transition in (TMTSF)2FSO3 probed by infrared spectroscopy
1264 1280 1296
560 570 580 590 600 1140 1150 1160
1358 1365 1372
Frequency (cm
Frequency (cm
Frequency (cm
Frequency (cm
Fig. 3. Reflectivity spectra (shifted for clarity) of some
phonons which experience changes during the metal-insulator
transition at 89 K: (a) vibration polarized along the a axis;
(b)-(d) vibrations polarized along the b′ axis.
The ν3(a1) vibrational mode of the FSO3 anion at
580 cm−1 is observed for the whole studied temperature
range. However, the lineshape of this mode in the metallic
phase above TMI is inverted with respect to the insulating
phase [see Fig. 3(a)], since the background dielectric con-
stant is negative as expected for highly conducting metals
at low frequencies. Such a change is a clear evidence for
the suppression of the Drude conductivity in the insulat-
ing phase of (TMTSF)2FSO3.
The mode at 728 cm−1 observed for temperatures
below TMI is particularly interesting, since its intensity
gradually increases on temperature decrease (see Fig. 4).
A similar behavior is found for the polarization perpen-
dicular to the stacks, E‖b′. Moreover, above the transi-
tion temperature a strong asymmetric mode is seen at
710 cm−1. This mode shifts to lower frequencies and gets
stronger with increasing temperature. This mode has not
been observed in any other earlier study of the Bechgaard
salts. Therefore, it would be natural to assign it to a vibra-
tion of FSO3 anion. However, such an assignment would be
in contradiction to the experimental observations, since:
680 700 720 740 760
E || b'
E || a
200 K
200 K
Frequency (cm
Fig. 4. Reflectivity spectra (shifted for clarity) of the vibration
at around 710 cm−1 at different temperatures for E||a and
E||b′.
(i) the ν5(e) and ν2(a1) vibrations of FSO3 located close
to the observed mode have frequencies which are by more
than 100 cm−1 higher or lower [14]; (ii) the intensity of
the anion vibration should not vanish at the order-disorder
transition point. Thus, one has to attribute the modes at
around 710 and 728 cm−1 to vibrations of the TMTSF
molecules. We suggest that both modes originate from the
ν51(b2u) in-plane vibration of the TMTSF molecule. Ac-
cording to the normal-coordinate analysis [12,16] its fre-
quency for a free TMTSF0.5+ cation is 702 cm−1.
The corresponding atomic movements involve stretch-
ing of the Se-C side bond and rocking of the adjacent
methyl group. For the b2u vibration the inversion symme-
try of the molecule is not preserved, causing its infrared-
activity for the polarization perpendicular to the stacks.
However, it is known that in (TMTSF)2X salts the dipole
moment corresponding to the ν51(b2u) vibration is very
small, and therefore this mode can hardly be detected
even for E‖b′ where it should have the strongest inten-
sity [16]. Nevertheless, in (TMTSF)2FSO3 this mode is
particularly strong even at room temperature. This find-
ing can be explained by the electrical dipole of the FSO3
anion pointing towards the Se-F bond. Similar to other
A. Pashkin et al.: Metal-insulator transition in (TMTSF)2FSO3 probed by infrared spectroscopy 5
Fig. 5. Schematic illustration of the ν51(b2u) vibration coupled
to the reorientation of the FSO3 electrical dipole moment. The
projection of the crystal structure on the b− c plane is shown.
Only the Se (large open circles) and C (small filled circles)
atoms of the TMTSF molecules are presented, together with
the displacements of the Se atoms. The grey filled circles be-
tween molecules denote the positions of the FSO3 anions, the
bold arrows show the two possible orientations of the anion
dipole moment (p1 and p2). Because of the symmetry prop-
erties of the b2u vibration the reorientation of the electrical
dipole moment leads to a change of polarization ∆p in the
perpendicular direction for any orientation of the anion dipole
moment.
non-centrosymmetric anions (ReO4, ClO4 etc.) the FSO3
anion has two possible symmetrically equivalent orienta-
tions for which the dipole moment points towards the Se
atoms of the neighboring TMTSF molecules. This situa-
tion is sketched in Fig. 5, where p1 and p2 are two possible
orientations of the FSO3 electrical dipole moment. During
the vibration the dipole moment of the anion follows the
position of the Se atom. Due to the symmetry properties
of the b2u vibration the nearest Se atoms on both sides
of the anion move in the same direction. Thus, for both
possible orientations of the dipole the b2u vibration results
in a change of the average polarization along the direction
of ∆p (Fig. 5).
The described coupling mechanism between the b2u
vibration and the dipole moment of the anion in
(TMTSF)2FSO3 should lead to a strong enhancement of
the infrared strength of the ν51(b2u) vibration for E‖b
since ∆p has the largest projection along this direction.
On the other hand, ∆p is perpendicular to the stacking
axis and the b2u mode should not appear for E‖a. This
is indeed observed in our experiment above the transition
temperature. Below the transition the long-range order of
the anion sublattice builds up. Then the anion dipole mo-
ment orientation is determined by the modulation of the
whole lattice and it is not dependent on the movement of
neighboring TMTSF molecules, i.e., the ν51(b2u) vibration
is decoupled from the FSO3 anions. Therefore, its inten-
sity should drop abruptly below TMI , in agreement with
our observations (see Fig. 4). Moreover, the observed de-
crease of the intensity of the coupled b2u mode at around
Table 2. The eigenfrequency, width (given in bracket), and
assignment of some vibrational modes observed for E‖b′ at
selected temperatures. All numbers are in cm−1.
95 K 80 K 45 K assignment
580 (1.3) 580 (0.9) 580 (0.8) ν3(a1) FSO3
710 (7.1) 728 728 (2.0) ν51(b2u)
1154 (5.0) 1150 (3.4) 1150 (3.3) ν48(b2u)
1157 (3.8) 1158 (2.5)
1280 (5.1) 1276 (4.2) 1276 (2.1) ν4(e) FSO3
1288 (16) 1286 (4.4) 1286 (1.7)
1363 (3.2) 1361 (2.8) 1361 (1.3) ν47(b2u)
1366 (4.3) 1367 (2.4) 1367 (1.8)
710 cm−1 at 95 K compared to higher temperatures can be
explained by taking into account short-range order fluc-
tuations above the transition, evidence for which is also
given below. Indeed, in the large enough dynamical regions
where the anions are ordered, the coupling is suppressed
and therefore the strength of the ν51(b2u) should decrease.
On cooling down below TMI a vibration appears again
at somewhat higher frequency (728 cm−1) for E‖b′ and
its strength gradually increases with decreasing temper-
ature. We suggest that this is the same ν51(b2u) vibra-
tion described above. Since it is decoupled from the anion
sublattice, its frequency is expected to increase abruptly
below the transition. The increase in strength for both
polarizations should be obviously related to the tempera-
ture dependence of the order parameter (i.e., the degree
of lattice modulation). One of the possible mechanisms
can be the emv coupling of the ν51(b2u) vibration to the
charge transfer bands along a and b′ directions. However,
the detailed picture of this emv coupling is not clear, since
the symmetry of b2u mode does not allow such kind of
coupling. One can speculate that the electric field of the
FSO3 dipoles in the ordered phase distorts the TMTSF
molecules making them non-centrosymmetric. Then the
emv coupling may become allowed for the b2u(ν51) mode.
Noticeable changes in the phonon mode spectra across
the metal-insulator transition are observed for E‖b′. The
list of the parameters of these modes at temperatures
above and below TMI is given in Table 2. An obvious split-
ting into two components is seen for the ν48(b2u) mode at
1154 cm−1 [see Figure 3(b)]. In addition, the damping of
the split components directly below TMI is lower than the
damping of the single component directly above the tran-
sition (see Table 2). This difference is probably related
to the precursor short-range order fluctuations above the
transition, which can induce a small splitting already in
the disordered phase. An evidence for such fluctuations
was found in x-ray diffuse scattering experiments [8,2].
This effect is even more clearly seen in the splitting of
two other modes: the doublet ν4(e) vibration of the FSO3
anion at 1280 cm−1 [Fig. 3(c)] and the ν47(b2u) mode at
around 1365 cm−1 [Fig. 3(d)]. For each of these modes
6 A. Pashkin et al.: Metal-insulator transition in (TMTSF)2FSO3 probed by infrared spectroscopy
above TMI one can resolve two weakly split components.
However, below TMI the splitting abruptly increases and
the damping decreases (see Table 2) indicating the onset
of long-range order. Since the described effect is observed
not only for the FSO3 anion vibration but also for two vi-
brations of the TMTSF cation, we can conclude that the
short-range order fluctuations involve the modulation of
the whole (TMTSF)2FSO3 lattice and not only the anion
sublattice.
4 Conclusion
We have performed an infrared spectroscopic study of
the metal-insulator transition in (TMTSF)2FSO3. The ob-
tained optical conductivity spectra for E‖a show a Drude-
like conductivity above the anion ordering temperature
and a charge transfer band formed below the transition.
The onset of this band is in agreement with the energy
gap value of 1500 cm−1 obtained from transport measure-
ments [4].
The analysis of the infrared-active vibrations leads to
the following conclusions: (i) the crystal structure modula-
tion below the metal-insulator transition leads to a strong
emv coupling of several ag vibrations which therefore be-
come infrared-active; (ii) short-range order fluctuations of
the FSO3 anions and the corresponding lattice modulation
exist above the transition temperature, as it is seen from
the splitting of some infrared-active modes for E‖b′; (iii)
a new infrared-active mode located at around 710 cm−1
with a peculiar temperature behavior is detected and as-
signed to the coupling between the b2u TMTSF molecule
vibration and the electrical dipole moment of the FSO3
anion. The latter feature has not been observed in any
other (TMTSF)2X salt showing a metal-insulator transi-
tion. This points out the important role of the electrical
dipole moment of the anion on the structural and dynam-
ical properties of the (TMTSF)2FSO3 salt.
5 Acknowledgements
We acknowledge the ANKA Angströmquelle Karlsruhe for
the provision of beamtime and thank M. Süpfle, D. Moss,
and B. Gasharova for technical assistance at the ANKA
IR beamline. The financial support of the DFG (Emmy
Noether-program) is acknowledged.
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Introduction
Experimental
Results and discussion
Conclusion
Acknowledgements
|
0704.0369 | The effect of Topcolor Assisted Technicolor, and other models, on
Neutrino Oscillation | arXiv:0704.0369v1 [hep-ph] 3 Apr 2007
August 9, 2021 18:21 WSPC - Proceedings Trim Size: 9in x 6in SCGT06-takeuchi
OCHA-PP-270, YITP-07-09, VPI-IPNAS-07-02
The effect of Topcolor Assisted Technicolor, and other models,
on Neutrino Oscillation
Minako Honda1, Yee Kao2, Naotoshi Okamura3, Alexey Pronin2, and Tatsu Takeuchi2∗
1Physics Department, Ochanomizu Women’s University, Tokyo 112-8610, Japan
2Physics Department, Virginia Tech, Blacksburg VA 24061, USA
3Yukawa Institute for Theoretical Physics, Kyoto University, Kyoto 606-8502, Japan
New physics beyond the Standard Model can lead to extra matter effects on
neutrino oscillation if the new interactions distinguish among the three flavors
of neutrino. In Ref. 1, we argued that a long-baseline neutrino oscillation ex-
periment in which the Fermilab-NUMI beam in its high-energy mode2 is aimed
at the planned Hyper-Kamiokande detector3 would be capable of constraining
the size of those extra matter effects, provided the vacuum value of sin2 2θ23
is not too close to one. In this talk, we discuss how such a constraint would
translate into limits on the coupling constants and masses of new particles in
models such as topcolor assisted technicolor.4
1. Introduction
When considering matter effects on neutrino oscillation, it is customary to
consider only the W -exchange interaction of the νe with the electrons in
matter. However, if new interactions beyond the Standard Model (SM) that
distinguish among the three generations of neutrinos exist, they can lead
to extra matter effects via radiative corrections to the Zνν vertex which
effectively violate neutral current universality, or via the direct exchange of
new particles between the neutrinos and matter particles.
For instance, topcolor assisted technicolor4 treats the third generation
differently from the first two and the Z ′ in this class of models couples
more strongly to the ντ than to the νe or νµ. In Extended Technicolor
(ETC) Models, such as that of Appelquist, Piai, and Shrock,5 the neutral
technimesons, which mix with the Z, couple to different generation fermions
differently, distinguishing among νe, νµ, and ντ . The diagonal ETC gauge
bosons also couple to the different generations differently, as well as the
∗Presenting Author
http://arxiv.org/abs/0704.0369v1
August 9, 2021 18:21 WSPC - Proceedings Trim Size: 9in x 6in SCGT06-takeuchi
large variety of leptoquark states in the model. Flavor distinguishing matter
effects from diagonal ETC and leptoquarks are induced by ETC gauge
boson mixing.
The effective Hamiltonian that governs neutrino oscillations in the pres-
ence of neutral-current lepton universality violation, or new physics that
couples to the different generations differently, is given by1
H = Ũ
λ1 0 0
0 λ2 0
0 0 λ3
Ũ † = U
0 0 0
0 δm221 0
0 0 δm231
U † +
a 0 0
0 0 0
0 0 0
be 0 0
0 bµ 0
0 0 bτ
where U is the MNS matrix,
a = 2EVCC , VCC =
2GFNe = Ne
, (2)
is the usual matter effect due to W -exchange between νe and the electrons,
and be, bµ, bτ are the extra matter effects which we assume to be non-equal.
We define the parameter ξ as
bτ − bµ
= ξ . (3)
Then, the effective Hamiltonian can be rewritten as
H = Ũ
λ1 0 0
0 λ2 0
0 0 λ3
Ũ † = U
0 0 0
0 δm221 0
0 0 δm231
U †+a
1 0 0
0 −ξ/2 0
0 0 +ξ/2
, (4)
where we have absorbed the extra b-terms in the (1, 1) element into a.
The extra ξ-dependent contribution in Eq. (4) can manifest itself when
a > |δm231| (i.e. E & 10GeV for typical matter densities in the Earth) in
the νµ and ν̄µ survival probabilities as
P (νµ → νµ) ≈ 1− sin2
2θ23 −
δm231
P (ν̄µ → ν̄µ) ≈ 1− sin2
2θ23 +
δm231
, (5)
where
∆ ≈ ∆31c213 −∆21c212 , ∆ij =
δm2ij
L , cij = cos θij , (6)
and the CP violating phase δ has been set to zero. As is evident from
these expressions, the small shift due to ξ will be invisible if the value of
sin2 2θ23 is too close to one. However, if the value of sin
2 2θ23 is as low
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as sin2 2θ23 = 0.92 (the current 90% lower bound), and if ξ is as large
as ξ = 0.025 (the central value of from CHARM/CHARM II6), then the
shift in the survival probability at the first oscillation dip can be as large as
∼ 40%. If the Fermilab-NUMI beam in its high-energy mode2 were aimed at
a declination angle of 46◦ toward the planned Hyper-Kamiokande detector3
in Kamioka, Japan (baseline 9120 km), such a shift would be visible after
just one year of data taking, assuming a Mega-ton fiducial volume and
100% efficiency. The absence of any shift after 5 years of data taking would
constrain ξ to1
|ξ| ≤ ξ0 ≡ 0.005 , (7)
at the 99% confidence level.
In the following, we look at how this potential limit on ξ would translate
into constraints on the Z ′ in topcolor assisted technicolor, and various types
of leptoquarks. A more comprehensive analysis will be presented in Ref. 7.
2. Topcolor Assisted Technicolor
Though there are several different versions of topcolor assisted technicolor,4
we consider here the simplest in which the quarks and leptons transform
under the gauge group
SU(3)s × SU(3)w × U(1)s × U(1)w × SU(2)L (8)
with coupling constants g3s, g3w, g1s, g1w, and g, respectively. It is assumed
that g3s ≫ g3w and g1s ≫ g1w. SU(2)L is the usual weak-isospin gauge
group of the SM. The first and second generation fermions are assumed to
be charged only under SU(3)w×SU(2)L×U(1)w, while the third generation
fermions are assumed to be charged only under SU(3)s × SU(2)L ×U(1)s.
The U(1) charges for both cases are set equal to the SM hypercharge. At
scale Λ ∼ 1 TeV, technicolor, which is included in the model to generate the
W and Z masses, is assumed to become strong and generate a condensate
(of something which is left unspecified) which breaks the two SU(3)’s and
the two U(1)’s to their diagonal subgroups:
SU(3)s × SU(3)w → SU(3)c , U(1)s × U(1)w → U(1)Y , (9)
which we identify with the usual SM color and hypercharge groups. The
massless unbroken U(1) gauge boson Bµ and the massive broken U(1) gauge
boson Z ′µ are related to the original U(1)s × U(1)w gauge fields Ysµ and
Ywµ by
Z ′µ = Ysµ cos θ1 − Ywµ sin θ1
August 9, 2021 18:21 WSPC - Proceedings Trim Size: 9in x 6in SCGT06-takeuchi
Bµ = Ysµ sin θ1 + Ywµ cos θ1 (10)
where
tan θ1 =
. (11)
The currents to which the Bµ and Z
µ couple to are:
1sYsµ+g1wJ
1wYwµ = g
′ (cot θ1J
1s − tan θ1J
1s + J
1w)Bµ ,
where
. (13)
The current J
1s + J
1w is the SM hypercharge current, and g
′ is the
SM hypercharge coupling constant.
The exchange of the Z ′ leads to the current-current interaction
(cot θ1J1s − tan θ1J1w) (cot θ1J1s − tan θ1J1w) , (14)
the J1sJ1s part of which does not contribute to neutrino oscillations on the
Earth, while the J1wJ1w part is suppressed relative to the J1wJ1s part by
a factor of tan2 θ1 ≪ 1. Therefore, we only need to consider the J1sJ1w
interaction which only affects the propagation of ντ . The effective potential
felt by ντ due to this interaction is
Vντ =
, (15)
and the effective ξ is
ξTT =
Vντ − Vνµ
(g′/MZ′)
(g/MW )2
tan2 θW
sin2 θW
The limit |ξTT | ≤ ξ0 = 0.005 then translates into:
MZ′ ≥ MZ
sin2 θW
≈ 440GeV . (17)
Unfortunately, this potential limit from the measurement of ξ is weaker
than what is already available from precision electroweak data,8 and from
direct searches for pp̄ → Z ′ → τ+τ− at CDF.9,10
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3. Generation Non-diagonal Leptoquarks
The interactions of leptoquarks with ordinary matter can be described in
a model-independent fashion by an effective low-energy Lagrangian as dis-
cussed in Ref. 11. Assuming the fermionic content of the SM, the most gen-
eral dimensionless SU(3)C × SU(2)L ×U(1)Y invariant couplings of scalar
and vector leptoquarks satisfying baryon and lepton number conservation
are given by:
L = LF=2 + LF=0 , (18)
where
LF=2 =
ejL − dciLνjL) + g
γµejL) + g
γµejR)
V +2µ
γµνjL) + g
γµejR)
V −2µ
γµejL)Ṽ
2µ + (u
γµνjL)Ṽ
ejL)S
3 − (uciLejL + dciLνjL)S
νjL)S
, (19)
LF=0 =
(uiRejL) + h
(uiLejR)
(uiRνjL)− hij2R(diLejR)
(diRejL)S̃
2 + (diRνjL)S̃
(uiLγ
µνjL + diLγ
µejL) + h
(diRγ
µejR)
V 01µ
2(uiLγ
µejL)V
3µ + (uiLγ
µνjL − diLγµejL)V 03µ +
2(diLγ
µνjL)V
Here, the scalar and vector leptoquark fields are denoted by S and V , their
subscripts indicating the dimension of their SU(2)L representation, and the
superscripts indicating the sign of the weak-isospin of each component. We
allow for generation non-diagonal couplings with the indices i and j indicat-
ing the quark and lepton generation numbers, respectively. The subscript L
or R on the coupling constants indicate the chirality of the lepton involved
in the interaction. For simplicity, color indices have been suppressed. The
leptoquarks S1, ~S3, V2, Ṽ2 carry fermion number F = 3B + L = −2, while
the leptoquarks S2, S̃2, V1, ~V3 have F = 0. The interactions that affect neu-
trino oscillation are those with (ij) = (12) or (13).
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Table 1. Constraints on the leptoquark couplings with all the leptoquark
masses set to 300 GeV. To obtain the bounds for a different leptoquark mass
MLQ, simply rescale these numbers with the factor (MLQ/300 GeV)
LQ CLQ δλ
upper bound from |ξ| ≤ ξ0
S1 +3 |g
|2 − |g13
|2 0.01
~S3 +9 |g
|2 − |g13
|2 0.003
S2 −3 |h
|2 − |h13
|2 0.01
S̃2 −3 |h̃
|2 − |h̃13
|2 0.01
V2 +6 |g
|2 − |g13
|2 0.005
Ṽ2 +6 |g̃
|2 − |g̃13
|2 0.005
V1 −6 |h
|2 − |h13
|2 0.005
~V3 −18 |h
|2 − |h13
|2 0.002
It is straightforward to calculate the effective potentials due to the ex-
change of these leptoquarks, as well as the effective values of ξ.7 Assuming a
common mass for leptoquarks in the same SU(2)L weak-isospin multiplet,
the effective ξ due to the exchange of any particular type of leptoquark can
be written in the form
ξLQ = CLQ
δλ2LQ/M
g2/M2
δλ2LQ
. (21)
Here, CLQ is a constant prefactor, and δλ
LQ represents
δλ2LQ = |λ12LQ|2 − |λ13LQ|2 , (22)
where λ
is a generic coupling constant. The values of CLQ and δλ
for the different types of leptoquark are listed in Table 1. The constraint
|ξLQ| ≤ ξ0 translates into:
MLQ ≥
|CLQ||δλ2LQ|
2GF ξ0
|CLQ||δλ2LQ| × (1700GeV) . (23)
Alternatively, one can fix the leptoquark mass and obtain upper bounds on
the leptoquark couplings:
∣δλ2LQ
2GF ξ0
M2LQ ≈
300GeV
. (24)
The values when MLQ = 300GeV are listed in the rightmost column of
Table 1. Thought it is often stated that generation non-diagonal couplings
of leptoquarks are strongly constrained by the absence of flavor chang-
ing neutral currents, it is only the products of the (ij) = (12) and (13)
couplings with other couplings that are constrained.12 The limits on the
August 9, 2021 18:21 WSPC - Proceedings Trim Size: 9in x 6in SCGT06-takeuchi
individual couplings can be improved considerably. The current leptoquark
mass bounds from direct searches at the Tevatron, LEP, and HERA are in
the 200∼300 GeV range assuming generation diagonal couplings set equal
4πα. At the LHC, leptoquarks, if they exist, can be expected to be
pair-produced copiously through gluon-gluon fusion. The expected sensi-
tivity is up to about 1.5 TeV.13 Depending on the value assumed for δλ2LQ,
the bound from Eq. (23) can be competitive.
Acknowledgments
We would like to thank Drs. Andrew Akeroyd, Mayumi Aoki, Masafumi
Kurachi, Robert Shrock, and Hiroaki Sugiyama for helpful discussions. This
research was supported in part by the U.S. Department of Energy, grant
DE–FG05–92ER40709, Task A (Kao, Pronin, and Takeuchi).
References
1. M. Honda, N. Okamura, and T. Takeuchi, arXiv:hep-ph/0603268.
2. NUMI Technical Design Handbook,
available at http://www-numi.fnal.gov/numwork/tdh/tdh index.html
3. Y. Itow et al., arXiv:hep-ex/0106019;
updated version available at http://neutrino.kek.jp/jhfnu/.
4. C. T. Hill, Phys. Lett. B 345, 483 (1995); G. Buchalla, G. Burdman,
C. T. Hill, and D. Kominis, Phys. Rev. D 53, 5185 (1996).
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6. J. Dorenbosch et al. [CHARM Collaboration], Phys. Lett. B 180, 303 (1986);
P. Vilain et al. [CHARM-II Collaboration], Phys. Lett. B 320, 203 (1994).
7. M. Honda, Y. Kao, N. Okamura, A. Pronin, and T. Takeuchi, in preparation.
8. R. S. Chivukula and J. Terning, Phys. Lett. B 385, 209 (1996);
W. Loinaz and T. Takeuchi, Phys. Rev. D 60, 015005 (1999).
9. D. Acosta et al. [CDF Collaboration], Phys. Rev. Lett. 95, 131801 (2005).
10. W. M. Yao et al. [Particle Data Group], J. Phys. G 33 (2006) 1.
11. W. Buchmüller, R. Rückl and D. Wyler, Phys. Lett. B 191, 442 (1987);
M. Tanabashi, p.412 of Ref. 10.
12. S. Davidson, D. C. Bailey and B. A. Campbell, Z. Phys. C 61, 613 (1994);
M. Leurer, Phys. Rev. D 49, 333 (1994); M. Leurer, Phys. Rev. D 50, 536
(1994); E. Gabrielli, Phys. Rev. D 62, 055009 (2000).
13. ATLAS detector and physics performance. Technical design report. Vol. 2,
CERN-LHCC-99-15, ATLAS-TDR-15; CMS physics : Technical Design Re-
port v.2 : Physics performance, CERN-LHCC-2006-021, CMS-TDR-008-2;
V. A. Mitsou, N. C. Benekos, I. Panagoulias and T. D. Papadopoulou, Czech.
J. Phys. 55, B659 (2005).
|
0704.0370 | Shaped angular dependence of the spin transfer torque and microwave
generation without magnetic field | The magnetization of a ferromagnetic body can be manipulated by transfer of spin angular momentum from a spin-polarized curren
Shaped angular dependence of the spin transfer torque and microwave generation without
magnetic field
O. Boulle1, V. Cros1, J. Grollier1, L. G. Pereira1,*, C. Deranlot1, F. Petroff1, G. Faini2, J. Barnaś3,
A. Fert1
1 Unité Mixte de Physique CNRS/Thales and Université Paris Sud XI, Route départementale 128,
91767 Palaiseau, France
2 Laboratoire de Photonique et de Nanostructures LPN-CNRS, Route de Nozay, 91460
Marcoussis, France
3 Department of Physics, Adam Mickiewicz University, Umultowska 85, 61-614 Poznań, Poland
Abstract: The generation of oscillations in the microwave frequency range is one of the most
important applications expected from spintronics devices exploiting the spin transfer
phenomenon. We report transport and microwave power measurements on specially designed
nanopillars for which a non-standard angular dependence of the spin transfer torque (wavy
variation) is predicted by theoretical models. We observe a new kind of current-induced
dynamics that is characterized by large angle precessions in the absence of any applied field, as
this is also predicted by simulation with such a wavy angular dependence of the torque. This type
of non-standard nanopillars can represent an interesting way for the implementation of spin
transfer oscillators since they are able to generate microwave oscillations without applied
magnetic field. We also emphasize the theoretical implications of our results on the angular
dependence of the torque.
The magnetization of a ferromagnetic body can be manipulated by transfer of spin angular
momentum from a spin-polarized current. This is the concept of spin transfer introduced by
Slonczewski [1] and Berger [2] in 1996. In most experiments, a spin-polarized current is injected
from a spin polarizer into a “free” magnetic element, for example in pillar-shaped magnetic
trilayers [3-6]. The phenomenon of spin transfer has a great potential for applications. It can be
used either to switch a magnetic configuration (the configuration of a magnetic memory for
example) [3-5] or to generate magnetic precessions and voltage oscillations in the microwave
frequency range[6-7]. In the most usual situations, such oscillations are observed in the presence
of a magnetic field.
From a fundamental point of view, spin transfer effects raise two different types of
problems [8]. First the spin transfer torque acting on a magnetic element is related to the
transverse spin polarisation of the current (transverse meaning perpendicular to the magnetization
axis of the element) and can be derived from spin-dependent transport equations [8-17]. On the
other hand, the description of the magnetic excitations generated by the spin transfer torque raises
problems of non-linear dynamics [8,18-20]. For example, in the simple limit where the excitation
is supposed to be a uniform precession of the magnetization (macrospin approximation), this
precession can be determined by introducing the spin transfer torque into a Landau-Lifshitz-
Gilbert (LLG) equation for the motion of the magnetic moment. However, the determination of
the spin transfer torque and the description of the magnetization dynamics cannot be regarded as
independent problems. In standard trilayered structures with in-plane magnetizations and with the
usual angular dependence, a switching regime is found at zero and low magnetic field and the
precession regime with generation of voltage oscillations is mainly observed above some
threshold field [8]. We will show that a new behavior, characterized by large angle precessions in
the absence of any magnetic field, can be obtained in specially designed structures presenting a
non-standard dependence of the spin transfer torque as a function of the angle between the fixed
magnetization of the polarizer and the magnetization of the free layer. This non-standard angular
dependence of the torque, that we call “wavy”, is obtained by choosing materials with different
spin diffusion lengths for the “fixed” and “free” magnetic layers, which changes the distribution
of the spin currents and spin accumulations in the structure.
The observation of spin transfer oscillations at zero field in structures with a “wavy” angular
dependence of the torque can represent a new way to obtain spin transfer oscillators operating
without any applied field, an other possible way being the use of exchange interactions or
anisotropy to generate local effective fields or non-collinear equilibrium configurations [21]. In
addition, the observation of a wavy angular dependence of the torque represents a valuable test of
the theory and shows that realistic predictions of the spin transfer torque and its angular
dependence in a given structure are now possible. As we will see, in the models we consider here
[15-16], the torque is calculated from parameters which, for most of them, can be derived from
former CPP-GMR experiments [22-23].
The usual behaviour observed in pillars with in-plane magnetizations along an anisotropy axis
corresponds to the standard angular dependence of the inset of Fig.1a, in which the torque starts
from zero at ϕ = 0 (P equilibrium state with parallel magnetizations of the fixed and free
magnetic layers) and keeps the same sign till it comes back to zero at ϕ = π (AP antiparallel
state). At zero field and starting from a P state for example (Fig.1b), a negative current (electrons
going from the free to the fixed layer in our convention) will destabilize the P state and stabilize
the AP state, i.e. can switch the system from P to AP. In the presence of a large enough applied
field favouring the P configuration, the torque cannot stabilize the AP state and leaves the system
in an intermediate precession state. This is what we call the standard behaviour with irreversible
switching at low field and precession at high field, as illustrated by Fig.2a (remark: in some low
field experiments however, the irreversible switching is preceded by precessions in a very narrow
current range just below the switching current).
The non-standard behaviour with precession at zero and/or low field presented in this article
is related to the existence of a wavy angular dependence of the torque acting on the free magnetic
layer. This oscillatory angular dependence, with an inversion of the torque between ϕ = 0 and ϕ =
π, is shown in Fig. 1a. We present the results of calculations in the models of Fert et al [15] and
Barnaś et al [16-17] for a Py(8)/Cu(10)/Co(8) pillar. With respect to standard structures like
Co/Cu/Co or Py/Cu/Py, the difference we have introduced is a large asymmetry between the spin
diffusion lengths (SDL) in the magnetic layers, with a long SDL in Co ≈ 38 nm (at room
temperature) and a short SDL in Py ≈ 4 nm [22-23]. The smaller spin asymmetry of the resistivity
in Co could also affect the angular dependence but we have checked by additional calculations
that the wavy variation comes primarily from the shorter SDL in the Py free layer and not from
the different spin asymmetry coefficients, as this has been mistakenly written in Ref.[24]. The
solid curves in Fig.1a correspond to the calculation in the model of Barnaś et al [16]. A wavy
angular dependence is also predicted by the model of Fert et al [15] which gives the terms of first
order in ϕ and (π-ϕ) in the vicinity of the colinear P or AP states (the solid straight lines at the
left and right edges of the graph in Fig.1). Due to the inversion at small values of ϕ, a negative
current (Fig. 1c) now stabilizes not only the AP state but also the P one and should be ‘inactive’.
This can be a solution, for example, to reduce the spin-transfer-induced noise that is detrimental
to read heads. In contrast, an appropriate positive current can destabilize both the P and AP states,
leading to a precessional solution the motion equation, even at zero field.
To validate these predictions, we have performed transport and microwave power
measurements at room temperature on Py(8)/Cu/Co(8) elliptical nanopillars of approximate
dimensions 100x155 nm². Only the top Py layer (free layer) and the Cu spacer are etched
through. The unetched Co layer (“fixed” layer) lies directly on the Ta/Cu bottom electrode. Very
similar results have been obtained on Py(8)/Cu/Co(4)/IrMn nanopillars in which the extended Co
layer is exchange biased by the IrMn one. We show in Fig. 2b the GMR signal of a
Py(8)/Cu/Co(8) sample. Starting, for example, from large negative fields, the switching to an AP
state at about 40 Oe is related to the magnetization reversal of the free layer (Py) to the positive
direction, as this can be found from subsequent CIMS experiments in which the current-induced
return to P is made harder by a larger positive field (consistently with a positive orientation of the
Py magnetization in the switching to the AP state). From the GMR minor cycles of the Py layer
(see Supplementary Information), we find that the coercive field of the Py layer is 90 Oe and the
dipolar field acting on it is 43 Oe.
The different behaviours observed for standard and wavy angular dependences are first
illustrated in Fig.2a and 2c. In Fig. 2a, we show the standard variation of differential resistance
(dV/dI) versus I measured on a Py(4 nm)/Cu(10 nm)/Py(15 nm) pillar: starting from a P state, a
negative current induces an irreversible switching from P to AP at low field and a reversible
variation with the characteristic peak of steady precessions at high field. In contrast, starting
again from a P magnetic configuration with magnetizations in the positive field direction but now
with a Py(8 nm)/Cu(10 nm)/Co(8 nm) pillar for which a wavy angular dependence is expected,
we detect (Fig. 2c) reversible peaks of dV/dI for positive currents and at very small fields on both
sides of Happ= 0. The peak current increases with increasing applied positive field as expected
since the P state becomes more stable. We have also performed experiments with an AP initial
state. We find that dV/dI first drops to the level of a P state at some positive current and then, at
higher current, exhibits the same characteristic precession peak we observe in measurements with
a P initial state (data not presented).
In Fig. 3, we present microwave power spectra recorded with the same P initial state and for
several values of the current. Fig.3a is for zero applied field (actually, Happ ≈ 2 Oe) and Fig.3b for
zero effective in-plane field (after subtracting the dipolar field). Coloured dots in the insets
indicate the values of the current on the corresponding dV/dI vs I curves. A peak in the
microwave power spectrum turns out approximately in some current range above the maximum
of dV/dI. The frequency f of the microwave peak increases with the current (blue shift), in
contrast with the red shift generally observed in standard pillars with in-plane magnetization.
Actually, with the standard angular dependence of the torque, the theoretical prediction is a
succession of red and blue shift regimes at increasing current but, in experiments with in-plane
applied fields, the crossover to a blue shift regime has been seldom observed [25]. In macrospin
simulations, a blue shift in f is predicted for the regime of out-of-plane (OP) precessions and is
also associated with a decrease of f with increasing in-plane field. As shown in Fig.3 c, we
observe this decrease of f with Happ. In Fig.4a, we present the current-field diagram of the
microwave power. Microwave signals are emitted only in the top left corner of the diagram, i.e.
at low field and in a zone which is also a region of increased resistance (Fig.4b). No excitation is
observed at higher field.
We can therefore put forward two main results from our microwave power data: i) Pillars in
which a wavy angular dependence of the spin transfer torque is expected, generate microwave
oscillations, but, in contrast with the standard behavior, when there are excited by positive
currents and at zero field; ii) These microwave oscillations present a blue shift of their frequency
with current, a behaviour generally associated with out-of-plane precessions.
We first want to exclude that the effects described in the preceding paragraphs could arise
from other origins than the wavy angular dependence of the STT. Could they arise from
excitations of the Co “fixed” layer? We can first argue that the same behaviour is also observed
when the 4nm thick and extended Co layer is pinned by an IrMn layer and that an excitation of a
thin Co layer in the presence of such a strong pinning is quite improbable. We can also point out
that, for un-pinned continuous magnetic layers, the switching current densities obtained by Chen
et al. [26] are about one order of magnitude larger than ours. In addition, whereas a reduction of
the thickness of the Co layer to 4 nm for the same 8nm Py thickness should make the excitation
of Co easier (smaller current), our experimental results are in the opposite direction.
The sample of Fig.2-3 exhibits the relatively simple behaviour predicted for a wavy angular
dependence of the torque in a macrospin picture, i.e. precessions at zero and low field in positive
current. However, in a series of five similar samples (with or without pinning by IrMn), we have
also observed additional features in transport measurements. For example, in some samples and
with an initial P state, we see not only peaks in dV/dI in positive current at zero or low field but
also partial or total switchings in negative current. These excitations can be ascribed to a non-
uniform distribution of the magnetization [27]. For a part of the sample, the angle ϕ between the
magnetizations of the two layers is above ϕc, the angle of torque inversion, and can be excited by
a negative current. However, we emphasize that these additional excitations observed in transport
measurements are never associated with peaks in the emitted power in the Gigahertz range. All
the samples share the same main features with microwave emission only at low field in positive
current.
We now present the theoretical implications of our experimental results and first comment briefly
on the origin of the wavy angular dependence of the spin transfer torque in our samples. The
physics governing this angular dependence can be discussed simply by considering that, in all the
models [8,13-17] based on interfacial absorption of the transverse spin component and boundary
conditions of the mixing conductance type (the language can be different in different
formalisms), the spin transfer torque is proportional to the transverse component of the spin
accumulation in the spacer layer. The key point is that the spin accumulation in a nonmagnetic
conductor is directly related to the gradient of the spin current along the current axis z,
dzjdm m /)(−∝ [28]. In configurations close to the P state of a standard pillar, with a thick
fixed layer and a thin free layer made of the same material, the spin polarization of the current in
the spacer decreases from the fixed layer to the free layer. This corresponds to a given sign of
the spin accumulation. But an opposite sign is expected if, in the same configuration, the spin
polarization of the current increases from the fixed layer to the free layer. This is what occurs for
our Py(8nm)/Cu(10nm)/Co(8nm) pillars in an angular range close to the P configuration, as this
can be seen from the spin accumulation calculated in the Section Methods. As shown in Fig 1a,
calculations of STT based on two different models reflect this inversion of the spin accumulation
by an inversion of the torque on the left part of the figure with respect to the standards case.
However, as shown the figure, the inversion is a little less pronounced (less steep slope) in the
model of Ref.[15] which goes beyond the simple mixing conductance approximation of Ref.[16].
For a further understanding, we have performed additional macrospin simulations of the
current-induced precessions by solving a Landau-Lifschitz Gilbert equation including a spin
transfer term using parameters compatible with the actual structure of the measured samples (see
Methods, the simulations have been performed by two of the co-authors, O.B. and J.G.,
independently of those published in Ref.[24]). The simulated current-field diagram at T = 0 K is
presented on Fig 4.d with a colour scale corresponding to the change of resistance. At high field
(Happ larger than the anisotropy field) and in the current range we have considered, the only
excitations are in-plane (IP) precessions occurring above a threshold current Ic1 and associated
with a small change of resistance (which also corresponds to a small microwave power). At low
field, the IP precessions above Ic1 (black and blue trajectories in Fig. 4c) are followed by out-of-
plane (OP) precessions (orange and red trajectories) above a second current threshold Ic2.
There is a general good agreement between the main features of the experimental and
calculated phase diagrams. In particular, the zone of OP precessions in the top left corner of the
diagram of Fig. 4d turns out to be also the zone where we measure the larger DC resistance
increase (Fig. 4b) and also detect microwave excitations (Fig.4a). Quantitatively, if one compares
the colours in Fig.4 b and c, one can see that the distribution of the resistance change in the
diagram is well reproduced and that the experimental ΔR in the OP zone is only somewhat
smaller than the calculated one (by about 20% in average). The simulations also give a
distribution of microwave power (not shown) concentrated in OP top-left zone as in the
experimental plot of Fig.4a but with a power which is about 80 times larger than the experimental
one. This could be due to several reasons. First, there are certainly technical factors, like a large
impedance mismatch in the detection circuit. Second, for the OP excitations, the limits of a
macrospin approach for a quantitative prediction [6,30], have been put forward by several
publications. Finally, for the IP precessions we could not detect in the microwave spectra, it can
be pointed out that a very small variation of GMR is expected for angles between P and an angle
similar to our ϕc in structures with our type of torque angular dependence [29]. This has also led
us probably to overestimate the resistance change and the microwave power, since our
calculation is based on a standard angular dependence of the GMR as sin2(ϕ/2).
A confirmation that the zone of maximum resistance and microwave excitations in the top-
left corner (positive currents and low fields) of the diagrams in Fig.4a-b can be identified with the
zone of Out-of-Plane precessions in the calculated diagram (Fig. 4c) comes from the current and
field dependence of the frequency. As shown in the inset of Fig. 4c, the simulations predict that a
decrease of the frequency at increasing current for IP precessions is followed by an increase at
the transition to OP precessions. This is in agreement with the frequency blue shift of the
microwave excitations detected in the same zone of the phase diagram. The simulations also
predict correctly the red shift for the variation with the field. Our simulations therefore support
the picture of a non-standard behavior induced by a wavy angular dependence of the STT torque
and characterized by out-of-plane precessions excited by positive current at zero and low field.
During the submission process, we learned that oscillations of vortex structures in thick Py
layers excited by STT have been observed at relatively low field [31,32]. However this leads to
oscillations at relatively low frequency, below 1 GHz for layers in our aspect ratio [33], and the
oscillations above 3 GHz we observe cannot be explained by this mechanism.
Leaning on recent theoretical models of spin transfer torque, our experimental results
should help designing more efficient spin transfer oscillators operating in a very small or even
without an applied magnetic field. This is a necessary step (among others) on the implementation
of these new spintronics-based oscillators in a microwave receiver system for telecommunication
applications.
Methods :
The multilayers are grown by sputtering onto oxidized Si substrates. Two types of stacks
were deposited : structure 1 = Au(20 nm)/Cu(5 nm)/Py(8 nm)/Cu(10 nm)/Co(8 nm)/Ta(10
nm)/Cu(80 nm)/Ta(10 nm) and structure 2 = Au(25 nm)/Py(8 nm)/ Cu(8 nm)/Co(4 nm)/IrMn(15
nm)/Ru(15 nm)/Cu(35 nm). Py stands for Permalloy. The results we present, are on a nanopillar
with structure 1, but very similar results are observed with structure 2 when the fixed Co layer is
pinned with an IrMn layer. This indicates that, even without an IrMn pinning layer, the
magnetization of the extended Co layer is similarly fixed.
For the nanofabrication process, we first defined (by e-beam lithography, evaporation
deposition and lift-off) a Ti(15 nm)/Au (55 nm) elliptical mask on the magnetic multilayer. Then,
the magnetic pillar is etched by ion milling with a real-time monitoring by mass spectroscopy
down to the Cu/Co interface. The bottom electrode is defined by optical lithography and ion
milling. The next step is a planarization of the pillar with a Su-8 resist layer that is also used to
electrically isolate the bottom and the top electrode. The Su-8 layer on the top of the pillar is
removed by reactive ion etching. Finally, the top Ti/Au electrode is defined by optical
lithography, evaporation deposition and lift-off.
We measured both the dc resistance and the differential resistance dV/dI using an additional
20µA ac current modulated at 5kHz. For the frequency-domain measurements, we applied a dc
current on the sample through a bias-T. The high frequency voltage signal is then amplified (68
dB) and analysed on a commercial spectrum analyzer. The power spectra we show are extracted
from the spectrum analyser (we do not correct them from a calibration done for quantities like the
frequency-dependent amplifier gain, the attenuation in the transmission lines, and the impedance
mismatches). They are only obtained by subtracting a reference spectra measured at Idc = 0 in the
same magnetic field conditions. Note that the measured emitted power is therefore only a fraction
of the actual emitted power from the pillars. Both transport and frequency measurements have
been performed at room temperature and with in-plane magnetic field.
The torques of Fig.1a have been calculated by introducing in the models of Refs.[16] and
[17] parameters mostly derived from CPP-GMR experimental data [22-23]. For respectively Au,
Py, Cu, Co and Ta, these parameters are: bulk resistivity ρ (μΩ.cm) = 2, 15, 2.9, 24, 170; bulk
spin asymmetry coefficient β = 0, 0.76, 0, 0.46, 0; spin diffusion length lsf (nm) = 35, 4, 350, 38,
10. For the interface parameters, respectively Au/Cu, Cu/Py, Cu/Co, Co/Ta, Ta/Cu, the
parameters are : interfacial resistance rb (fΩ.m²) = 0.17, 0.5, 0.51, 0.5, 0.5; interfacial spin
asymmetry coefficient γ : 0, 0.7, 0.77, 0.7, 0; interfacial spin memory loss coefficient δ : 0.13,
0.25, 0.25, 0.25, 0.1. Note that the values of the Co and Ta resistivity have been measured on thin
film we had grown in the same conditions. The unknown value of lsf in Ta has been estimated by
fitting the calculated and experimental variations of resistance ΔR. We have also used the same
parameters in routine programmes [34] developed for the CPP-GMR to calculate the spin
accumulation in the spacer layer for our structure and in a standard structure Py (15 nm)/Cu (10
nm)/Py(2 nm), respectively – 2.2 and + 2.0 in arbitrary units and check the change of sign at the
origin of the wavy angular dependence.
For the simulations of the magnetization dynamics, we have solved a Landau Lifschitz
Gilbert equation including a spin transfer term of the form M1x(M1xM2) with the angular
dependence shown in Fig.1a. The calculations are performed at zero temperature. The saturation
magnetization µ0Ms = 0.87 T has been derived from ferromagnetic resonance experiments
performed on a Cu(6nm)/Py(7nm)/Cu(6nm) layer at room temperature. The other parameters are
the anisotropy field Han= 0.009 T, the gyromagnetic factor γ0 = 2.21 105(s.A/m)-1, α = 0.011. The
area of the pillars is about 1.38 104 nm², as derived by scanning electron microscopy.
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Metal/Metal Interfaces: an Experimentalist's Critical Review. cond-mat/0610085 (2006).
[24] Gmitra, M. and Barnaś, J., Current-driven destabilization of both collinear configurations in
asymmetric spin valves. Phys. Rev. Lett. 96, 207205 (2006).
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currents. Phys. Rev. B 72, 064430 (2005).
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film, Appl. Phys. Lett. 84, 380-382 (2004)
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macrospin concept. Phys. Rev. Lett. 96, 217202 (2006).
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Phys. Rev. B 48, 7099-7113 (1993).
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Phys. Rev. B 71, 100401 (2005).
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Rev. B 72, 024455 (2005)
[31] Pribiag, V. S., Krivorotov, I. N., Fuchs, G. D., Braganca, P. M., Ozatay, O., Sankey J. C.,
Ralph, D. C., Buhrman R. A., Magnetic vortex oscillator by dc spin-polarized current, cond-
mat/0702253 (2007).
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Oscillations in Spin Transfer Nanocontacts, , cond-mat/0702416 (2007).
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[34] H. Jaffrès, http://www.trt.thalesgroup.com/ump-cnrs-thales
Correspondence and requests for materials should be adressed to V. C. The authors declare they
have no competing financial interests
Acknowledgements
The authors thank M. Gmitra for the calculations of Fig 1b based on the model of Ref[16]. We
would like also to acknowledge H. Hurdequint for FMR measurements, L. Vila for assistance in
fabrication, O. Copie and B. Marcilhac for assistance in transport and frequency measurements
and M.R. Pufall for discussions. This work was partly supported by the french National Agency
of Research ANR through the PNANO program (MAGICO PNANO-05-044-02) and the EU
through the Marie Curie Training network SPINSWITCH (MRTN-CT-2006-035327). J. B.
acknowledges support by funds from the Polish Ministry of Science and Higher Education as a
research project (2006-2009).
*Present address : Instituto de Física, UFRGS, 91501-970 Porto Alegre, RS, Brazil
Figure captions
Figure 1 Angular dependence of the spin transfer torque for a standard and a ‘wavy’
angular dependence. a, Variation of the spin transfer torque on the free Py layer of a
Au(infinite)/Cu(5 nm)/Py(8 nm)/Cu(10 nm)/Co(8 nm)/Ta(10 nm)/Cu(infinite) multilayer as a
function of the angle ϕ between the magnetizations of the free Py and fixed Co layers for positive
and negative currents. The solid curves are calculated in the model of Barnaś et al [17], the solid
straight lines represent the slopes of the torque variation as the angle tends to 0 and π and have
been derived from the small angle expression of Fert et al [16]. The parameters used in the
calculations and mainly derived from CPP-GMR data are listed in the Section Methods. Inset :
typical variation of the spin transfer torque as a function of the angle between the magnetizations
of the free and fixed layers for a standard trilayer structure (case of Co/Cu/Co from Ref.[10]). b-
c, Sketches showing schematically the direction (blue arrow) of the spin transfer torque on the
free layer for configurations close to the P and AP configurations of the free layer (m) and fixed
layer (M) magnetizations for a standard (b) and a wavy (c) angular dependence of the torque.
Figure 2 Transport measurements on nanopillars with standard or “wavy” angular
dependence of the spin transfer torque. a, Differential resistance vs current measured for a
nanopillar with a standard structure Py(15 nm)/Cu(10 nm)/Py(4 nm) at “low field” (H = 6 Oe)
and “high field” (H = 133 Oe). In the latter case (precession), the applied field is larger than the
coercive field equal to H = 133 Oe. Curves are offset for clarity. b-c : Transport data for a Co(8
nm)/Cu(10 nm)/Py(8 nm). nanopillar. b, Resistance vs field at low current (I = 200 µA). c,
Differential resistance vs current for different applied fields around zero. These fields correspond
to the coloured symbols in b.
Figure 3. Microwave power spectra for the Co(8 nm)/Cu(10 nm)/Py(8 nm) nanopillar of
Fig.2b-c. a, Microwave power spectra for an applied field close to zero (Happl = 2 Oe) at different
currents corresponding to the coloured symbols in the inset. Inset: dV/dI vs I for Happl = 2 Oe. b,
Microwave spectra for different applied currents corresponding to the symbols in inset for an
effective (applied + dipolar) field of about zero (Happ = 43 Oe). Inset in b : dV/dI vs I for Happ =
43 Oe. c, Microwave spectra for I = 9 mA at different positive applied fields. Spectra are offset
for clarity.
Figure 4 Experimental and simulated spin-transfer-induced high frequency dynamics for a
Co(8nm)/Cu10nm)/Py(8nm) nanopillar. a, Experimental integrated power between 0.1 to 8
GHz in colour scale as a function of field and current. b, Normalized experimental resistance in
colour scale as a function of field and current (a reference curve has been subtracted to the
experimental R vs I curves to remove the changes in resistance due to Joule heating). c-d :
Simulated dynamics of the magnetization in a macrospin approach c, Results of macrospin
numerical calculations of LLG equation as a function of current and field at T = 0K. The black
line indicates the onset of current-induced precession. Inset in c : Variation of the calculated
frequency as a function the current for Happ = 0 Oe. d Magnetization trajectories for Happ=0 (black
arrow in c) at several increasing applied currents.
π/2 π
Standard angular
dependence of the spin
transfer torque
I > 0
P stable
AP unstable
(H = 0)
Wavy angular
dependence of the spin
transfer torque
I < 0
P unstable
AP stable
(H = 0)
I < 0
P stable
AP stable
(H = 0)
I > 0
P unstable
AP unstable
(H = 0)
Fig. 1 Boulle et al.
-10 -5 0 5 10
Current (mA)
- 19 Oe
- 3 Oe
9 Oe
22 Oe
-100 -50 0 50 100
12.95
12.97
12.99
Magnetic field (Oe)
-6 -4 -2 0 2 4 6
Current (mA)
Fig. 2
Boulle et al.
9,5 mA
8,5 mA
7,5 mA
6,5 mA
4 6 8 10
12.25
12.39
I (mA)
I (mA)
1 2 3 4
36 Oe
24 Oe
- 4 Oe
Frequency (GHz)
-19 Oe
Happ ≈ 0 (2 Oe) a
Heff ≈ 0
(Happ=43 Oe)
I = 9 mA
11 mA
10 mA
Fig. 3
Boulle et al.
9,5 mA
8,5 mA
7,5 mA
6,5 mA
4 6 8 10
12.25
12.39
I (mA)
I (mA)
1 2 3 4
36 Oe
24 Oe
- 4 Oe
Frequency (GHz)
-19 Oe
Happ ≈ 0 (2 Oe) a
Heff ≈ 0
(Happ=43 Oe)
I = 9 mA
11 mA
10 mA
Fig. 3
Boulle et al.
4 6 8 10
12.25
12.39
I (mA)
I (mA)
1 2 3 4
36 Oe
24 Oe
- 4 Oe
Frequency (GHz)
-19 Oe
Happ ≈ 0 (2 Oe) a
Heff ≈ 0
(Happ=43 Oe)
I = 9 mA
11 mA
10 mA
Fig. 3
Boulle et al.
0 50 100 150
Magnetic field (Oe)
Power (pW)0 2.74
0 50 100 150
Magnetic field (Oe)
9.4 mA
6.9 mA
5.7 mA
3.7 mA
0 50 100 150
4 6 8
Current (mA)
=0 Oe
Magnetic field (Oe)
plane
In plane
Parallel state
ΔRdc(mΩ)
ΔRdc(mΩ)
Fig 4.
Boulle et al.
Article File #1
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Figure 1
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|
0704.0371 | Dark energy interacting with neutrinos and dark matter: a
phenomenological theory | Dark energy interacting with neutrinos and dark matter: a
phenomenological theory
G. M. Kremer∗
Departamento de F́ısica, Universidade Federal do Paraná
Caixa Postal 19044, 81531-990 Curitiba, Brazil
October 29, 2018
Abstract
A model for a flat homogeneous and isotropic Universe composed of dark energy, dark
matter, neutrinos, radiation and baryons is analyzed. The fields of dark matter and neutrinos
are supposed to interact with the dark energy. The dark energy is considered to obey either
the van der Waals or the Chaplygin equations of state. The ratio between the pressure and the
energy density of the neutrinos varies with the red-shift simulating massive and non-relativistic
neutrinos at small red-shifts and non-massive relativistic neutrinos at high red-shifts. The
model can reproduce the expected red-shift behaviors of the deceleration parameter and of the
density parameters of each constituent.
The recent astronomical measurements of type-IA supernovae [1, 2, 3, 4] and the analysis of
the power spectrum of the CMBR [5, 6, 7, 8, 9] provided strong evidence for a present accelerated
expansion of the Universe [3, 10, 11, 12, 13, 14]; the nature of the responsible entity, called dark
energy, still remains unknown. Furthermore, the measurements of the rotation curves of spiral
galaxies [15] as well as other astronomical experiments suggest that the luminous matter represents
only a small amount of the massive particles of the Universe, and that the more significant amount
is related to dark matter. That offered a new setting for cosmological models with dark energy
and dark matter and in these contexts many interesting phenomenological models appear in the
literature analyzing the interaction of neutrinos [16, 17, 18] and dark matter [19, 20, 21, 22, 23, 24]
with dark energy. With respect to dark energy some exotic equations of state were proposed in the
literature and among others we quote the van der Waals [25, 26, 27, 28, 29] and the Chaplygin [30,
31, 32, 33] equations of state.
In the present work a very simple cosmological model – for a homogeneous, isotropic and flat
Universe composed by dark matter, dark energy, baryons, radiation and neutrinos – is investigated
where the dark energy is modeled either by the van der Waals or the Chaplygin equations of state
and interact with neutrinos and dark matter. Units have been chosen so that 8πG/3 = c = 1,
whereas the metric tensor has signature (+,−,−,−).
Let a homogeneous, isotropic and spatially flat Universe be characterized by the Robertson
Walker metric ds2 = dt2 − a(t)2δijdxidxj , where a(t) denotes the cosmic scale factor. The sources
of the gravitational field are related to a mixture of five constituents described by the fields of dark
energy, dark matter, baryons, neutrinos and radiation. The components of the energy-momentum
tensor of the sources is written as
(T µν) = diag(ρ,−p,−p,−p), (1)
where ρ and p denote the total energy density and pressure of the sources, respectively. In terms
of the energy densities and pressures of the constituents it follows
ρ = ρb + ρdm + ρr + ρν + ρde, p = pb + pdm + pr + pν + pde. (2)
∗[email protected]
http://arxiv.org/abs/0704.0371v1
Above the indexes (b, dm, r, ν, de) refer to the baryons, dark matter, radiation, neutrinos and dark
energy, respectively.
The conservation law of the energy-momentum tensor T µν ;ν = 0 leads to the evolution equation
for the total energy density of the sources, namely
ρ̇+ 3
(ρ+ p) = 0, (3)
where the dot refers to a differentiation with respect to time.
The baryons and radiation are considered as non-interacting fields so that the evolution equa-
tions for their energy densities read
ρ̇b + 3
ρb = 0, ρ̇r + 4
ρr = 0, (4)
once the baryons represent a pressureless fluid, i.e., pb = 0, and the radiation pressure is given in
terms of its energy density by pr = ρr/3.
According to a model proposed by Wetterich [19] the evolution equation for the energy density
of a pressureless (pdm = 0) dark matter field which interacts with a scalar field φ is given by
ρ̇dm + 3
ρdm = βρdmφ̇. (5)
Here the scalar field plays the role of the dark energy and β is a constant which couples the fields
of dark matter and dark energy.
For interacting neutrinos with dark energy it is supposed that the evolution equation of the
energy density is given by (see [17, 18])
ρ̇ν + 3
(ρν + pν) = α(ρν − 3pν)φ̇. (6)
The coefficient α is connected with the mass of the neutrinos and for more details one is referred
to [17, 18] and to the references therein. Here α will be consider a phenomenological coefficient
that couples the dark energy field with the neutrinos. Note that if pν = ρν/3, there is no coupling
between the fields of dark energy and neutrinos. Moreover, it is also important to note that the
neutrinos in the past must behave as massless particles where the relationship between the pressure
and the energy density is pν = ρν/3. Due to the coupling of the neutrinos with the scalar field
they become massive and non-relativistic. For these reasons a barotropic equation of state for the
neutrinos is proposed where the ratio between the pressure and the energy density wν = pν/ρν ,
given in terms of the red-shift z, reads
K3(1/z)
K2(1/z)
K3(1/z)
K2(1/z)
. (7)
Above K2(1/z) and K3(1/z) are modified Bessel functions of second kind. For small values of z,
wν tends to the non-relativistic limit equal to 2/3, whereas for large values of z, wν tends to the
relativistic limit equal to 1/3. It is noteworthy that for red-shifts z ≈ 10 this ratio reaches the value
wν ≈ 1/3 and the coupling between the neutrinos and the dark energy is negligible. The expression
given in (7) is motivated by the equation of the specific heat of a relativistic gas (see e.g. [34]).
The evolution equation for the energy density of the dark energy field is obtained from equations
(2) through (6), yielding
ρ̇de + 3
(ρde + pde) = −αφ̇(ρν − 3pν)− βρdmφ̇. (8)
The energy density and pressure of the dark energy are connected with the scalar field by φ̇ =√
ρde + pde. Since the purpose of this work is to develop a phenomenological theory, it is assumed
0 2 4 6 8 10
vw Ωb
Figure 1: Density parameters as functions of red-shift: van der Waals fluid (solid lines) and Chap-
lygin fluid (dashed lines).
that the dark energy field behaves either as a van der Waals or a Chaplygin fluid with an equation
of state given by [28, 29, 30, 31, 32, 33]
pvw =
8wvwρvw
3− ρvw
− 3ρ2vw, pch = −
, (9)
where wvw and A are positive free parameters in the van der Waals and Chaplygin equations of
state, respectively.
For the determination of the time evolution of the energy densities one has to close the system
of differential equations by introducing the Friedmann equation
= ρ. (10)
From now on the red-shift will be used as a variable instead of time thanks to the following
relationships
ρ(1 + z)
. (11)
Equations (4) can be easily integrated leading to the well-known dependence of the energy
densities of the baryons and radiation with the red-shift
ρr(z) = ρr(0)(1 + z)
4, ρb(z) = ρb(0)(1 + z)
3, (12)
whereas equations (5), (6) and (8) become a system of coupled differential equations for the energy
densities ρdm, ρν and ρde, namely,
(1 + z)ρ′dm − 3ρdm
(ρde + pde)/ρ
= −βρdm, (13)
(1 + z)ρ′ν − 3(ρν + pν)
(ρde + pde)/ρ
= −α(ρν − 3pν), (14)
(1 + z)ρ′de − 3(ρde + pde)
(ρde + pde)/ρ
= βρdm + α(ρν − 3pν). (15)
In the above equations the prime refers to a differentiation with respect to the red-shift.
In order to solve the coupled system of differential equations (13) – (15) one has to specify initial
values for the energy densities at z = 0. The following initial values for the density parameters
Ωi(z) = ρi(z)/ρ(z) taken from the literature (see [35] for a review) were chosen: Ωde(0) = 0.72,
0 500 1000 1500 2000 2500 3000
Figure 2: Density parameters as functions of red-shift for a van der Waals fluid as dark matter.
Ωdm(0) = 0.229916, Ωb(0) = 5 × 10−2, Ωr(0) = 5 × 10−5, Ων(0) = 3.4 × 10−5. Moreover, one has
to specify values for the coupling parameters α and β and for the parameters wvw and A which
appear in the van der Waals and Chaplygin equations of state (9). One way to fix the two last
parameters is through the use of the value of the deceleration parameter q = 1/2+ 3p/2ρ at z = 0.
Indeed, by considering q(0) = −0.55 it follows wvw = 0.33851 and A = 0.50403. For the coupling
parameters two sets of values were chosen, namely, (a) α = 5 × 10−5 and β = −5 × 10−5 for the
van der Waals equation of state and (b) α = 10−1 and β = −10−2 for the Chaplygin equation of
state. Its is also important to note that by increasing the value of the coupling parameter α (and/
or β) the transfer of energy between the dark energy and neutrinos (and/or dark matter) becomes
more efficient.
In Fig. 1 the density parameters are plotted as functions of the red-shift for values in the range
0 ≤ z ≤ 10. The straight lines refer to the case where the van der Waals equation of state is used
to describe the dark energy field whereas the dashed lines correspond to the Chaplygin equation
of state. The two density parameters that represent the dark energy field are denoted by Ωvw and
Ωch. One can infer from this figure that the dark energy density parameter tends to zero for high
red-shifts when the van der Waals equation of state is used, whereas it tends to a constant value
for the Chaplygin equation of state. While for high red-shifts the van der Waals equation of state
simulates a cosmological constant with pvw = −ρvw, the pressure of the Chaplygin fluid vanishes
indicating that it becomes another component of the dark matter field (see also the behavior of
the pressures indicated in Fig. 4). It is also important to note that the density parameters of the
baryons and of the dark matter increase more with the red-shift for the van der Waals equation of
state, since there is an accentuated decrease in the density parameter of the dark energy for this
case. Note that the density parameters of the radiation and neutrinos are very small in this range
of the red-shift and are not represented in this figure.
The behavior of the density parameters for the cases of the van der Waals and Chaplygin
equations of state are shown in Figs. 2 and 3, respectively, for red-shifts in the range 0 ≤ z ≤ 3000.
One can conclude from these figures, as expected, that the density parameters of the neutrinos
and radiation increase with the red-shift whereas those of the baryons and dark matter decrease.
Furthermore, the equality between the “matter” and “radiation” fields occurs when z ≈ 3000 for
the case where the dark matter field is modeled as a van der Waals fluid and z ≈ 4200 for the case of
a Chaplygin fluid. This can be easily understood, since in the latter case the dark energy becomes
dark matter for high red-shifts contributing for the density parameter of the “matter” field.
In Fig. 4 are plotted the deceleration parameter and the ratio between the pressure and the
energy density for both cases, the large frame corresponding to the van de Waals fluid whereas the
small frame to the Chaplygin fluid. For both cases the deceleration parameter at z = 0 is equal to
q(0) = −0.55, since this value was fixed in order to find the parameters wvw and A in the equations
of state (9). The transition from the decelerated to the accelerated phase of the Universe occurs at
zT = 0.73 and zT = 0.53 for the van der Waals and Chaplygin equations of state, respectively. It
0 500 1000 1500 2000 2500 3000
+ Ων
Figure 3: Density parameters as functions of red-shift for a Chaplygin fluid as dark matter.
-0.75 -0.5 -0.25 0 0.25 0.5 0.75 1 1.25 1.5 1.751.5 2
-0.5 0 0.5 1 1.51.5 2
Figure 4: Deceleration parameter and ratio between the pressure and the energy density as functions
of red-shift: large frame (van der Waals), small frame (Chaplygin).
0 2 4 6 8 10
with interactions
without interactions
Figure 5: Density parameters as functions of red-shift for a Chaplygin fluid with and without
interactions.
is interesting to note that while the Chaplygin equation of state simulates a cosmological constant
with pch = −ρch for negative red-shifts which implies an accelerated phase of the Universe in the
future, the van der Waals equation of state leads to a positive pressure and brings the Universe to
another decelerated phase in the future. It is noteworthy to call attention that for positive values
of the red-shift, the solution of the coupled differential equations (13) through (15) predicts that
the van der Waals fluid behaves close to a cosmological constant with pvw ≈ −ρvw. This behavior
does not lead to a new transition from a decelerated to an accelerated phase in the very early
Universe, since the energy density of the radiation field increases so that the radiation pressure
becomes larger than that of the van der Waals fluid. For high red-shifts the Universe first becomes
dominated by the baryon and dark matter fields and for higher red-shifts by the radiation field.
This model does not attempt to model the inflationary period, where the inflaton field dominates
a short rapid evolution of the Universe.
As final remarks we call attention to the fact that one expects that the coupling between dark
energy, dark matter and neutrinos should be weak so that the parameters α and β are restricted to
small values. The difference between the parameters adopted for the van der Waals and Chaplygin
equations of state is due to stability conditions of the non-linear coupled system of differential
equations (13) – (15), the van der Waals equation of state being more unstable for large values
of these parameters than the Chaplygin equation of state. In Fig. 5 we have plotted the density
parameters as functions of the red-shift for the case where a Chaplygin equation of state is used as
dark energy. One can infer from this figure that the decay of the dark energy density parameter
and the increase of the dark matter density parameter with the red-shift are more pronounced when
there exists a coupling between the fields. The density parameter of the baryons remains unchanged
since the baryons are uncoupled.
As final comment it is important to note that even without couplings between the fields of
dark energy, dark matter and neutrinos, this phenomenological model – with the equations of state
of van der Waals and Chaplyging as dark energy – can describe satisfactorily the evolution of a
Universe whose constituents are dark energy, dark matter, baryons, neutrinos and radiation.
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|
0704.0372 | Levy-Lieb constrained-search formulation as a minimization of the
correlation functional | Levy-Lieb constrained-search formulation as a minimization of
the correlation functional.
Luigi Delle Site∗
Max-Planck-Institute for Polymer Research
Ackermannweg 10, D 55021 Mainz Germany.
Abstract
The constrained-search formulation of Levy and Lieb, which formally defines the exact
Hohenberg-Kohn functional for any N -representable electron density, is here shown to be equivalent
to the minimization of the correlation functional with respect to the N − 1 conditional probability
density, where N is number of electrons of the system. The consequences and implications of such
a result are here analyzed and discussed via a practical example.
PACS numbers: 03.65. w, 71.10. w, 71.15.Mb
∗Electronic address: [email protected]
http://arxiv.org/abs/0704.0372v1
mailto:[email protected]
I. INTRODUCTION
The Hohenberg-Kohn (HK) theorem [1] has opened new perspectives to the calculations
of electronic-based properties of condensed matter [2], and, an aspect often disregarded, has
given profound new insights into the general understanding of quantum mechanics. In fact
the 3N -dimensional Schrödinger problem for the ground state of an electronic system:
HNψ(r1, ...r2) = E0ψ(r1, ...r2);HN =
i=1,N
∇2i ) +
i=1,N
v(ri) +
where v(ri) is the external potential,
is the electron-electron Coulomb term, E0
is the energy of the ground state and ψ(r1, ...rN) is the 3N -dimensional antisymmet-
ric ground state wavefunction [3], is transformed into a ”manageable” variational prob-
lem in three dimensions where the central role is played by the electron density: ρ(r) =
ψ∗(r, r2, .....rN)ψ(r, r2, .....rN)dr2....drN , where ΩN−1 is the N − 1 spatial domain.
In explicit terms the variational problems is written as:
E0 =MinρE[ρ] (2)
where
Ω ρ(r)dr = N (Ω being the spatial domain of definition) and E[ρ] = T [ρ] + Vee[ρ] +
Vext[ρ] is the energy functional composed respectively by the kinetic, electron-electron po-
tential and the external potential functional. However in its original formulation the HK
theorem and the related variational problem have got a restricted field of applicability; it
is valid only if the electron density ρ(r) is v-representable, that is if ρ(r) is the density cor-
responding to an antisymmetric wavefunction of the ground-state of an Hamiltonian of the
form of Eq.1. It follows that the correct formulation of the variational problem becomes:
E0 =MinρEv[ρ] (3)
where v refers to the v-representability of ρ(r). As discussed in Ref.[2], there are no general
conditions for a density to be v-representable and this makes the use of the HK theorem and
its associated variational principle not practical. A generalization of the HK theorem which
does not require ρ(r) to be v-representable was found, in parallel, by M.Levy [4] and E.Lieb
[5] and it is usually known as the Levy constrained-search formulation or Levy-Lieb
constrained-search formulation [6]; in this paper we adopt the latter terminology. We
also notice that recently P.Ayers [7] has further clarified this concept and developed an
axiomatic treatment of the Hohenberg-Kohn functional. In the following we briefly describe
the crucial aspects of the abovementioned approach which are relevant for the current work.
The starting point of the theory is the distinction between the ground state wavefunction, ψ,
and a wavefunction ψλ that also integrates to the ground state electron density ρ(r). Since
ψ is the ground state wavefunction, we have:
〈ψλ |HN |ψλ〉 ≥ 〈ψ |HN |ψ〉 = E0. (4)
Taking into account that Vext[ρ] is a functional of ρ only, Eq.4 can be written as:
〈ψλ |T + Vee|ψλ〉 ≥ 〈ψ |T + Vee|ψ〉 (5)
where T and Vee are respectively the kinetic and Coulomb electron-electron operator as
defined in Eq.1. The meaning of Eq.5 is that ψ is the wavefunction that minimize the
kinetic plus the electron-electron repulsion energy and integrates to ρ. It follows that the
initial variational problem of Eq.2 can be transformed in a double hierarchical minimization
procedure which formally allows for searching among all the ρ’s which are N -representable,
i.e. it can be obtained from some antisymmetric wavefunction; this is a condition which
is much weaker and more controllable than the v-representability. In explicit terms such a
formulation is written as:
E0 =Minρ
Minψλ→ρ 〈ψλ |T + Vee|ψλ〉+
v(r)ρ(r)dr
. (6)
The inner minimization is restricted to all wave functions ψλ leading to ρ(r), while the outer
minimization searches over all the ρ’s which integrate to N . The original HK formulation can
then be seen as a part of this new one once its universal functional, F [ρ] = 〈ψ |T + Vee|ψ〉
is written as:
F [ρ] =Minψ→ρ 〈ψ |T + Vee|ψ〉 . (7)
The purpose of this work is to show that F [ρ] can be determined solely by a minimiza-
tion with respect to the N − 1 conditional probability density of the electron correlation
functional. This latter will be shown to be composed by the non local Fisher information
functional [8] and the electron-electron two-particle Coulomb term. The advantage of this
representation is manifold; it further clarifies the connection of electronic properties to the
Fisher theory and shows that the knowledge of such a functional is the crucial ingredi-
ent in density functional based approaches; it also identifies the Weizsacker kinetic term,
∫ |∇ρ(r)|2
dr, as necessary component of the universal functional F [ρ] and, in practical terms,
offers an objective criterion of evaluation of ”approximate” exchange and correlation func-
tional, i.e. among two functionals, the physically better founded is the ”smaller” one. In
order to show the practical aspects of our idea we illustrate a potential application.
II. THE NEW REPRESENTATION
Before writing the functional in the conditional probability density formalism, we need
to define such a quantity. Let us consider a generic fermionic wavefunction ψ(r1, ....rN), for
simplicity we consider a real wavefunction, but the extension to a complex one can be also
done [9]; we do not consider the spin dependence explicitly, however this will not influence
the main conclusions. Then the N -particle probability density is [10, 11]:
Nψ∗(r1, ....rN )ψ(r1, ....rN) = Θ
2(r1, ...., rN) (8)
and this can formally decomposed as [10, 11]:
Θ2(r1, ...., rN) = ρ(r1)f(r2, ......., rN/r1) (9)
where ρ(r1) is the one particle probability density (normalized to N) and f(r2, ......., rN/r1)
is the N − 1 electron conditional (w.r.t. r1) probability density, i.e. the probability density
of finding an N −1 electron configuration, C(r2, ......., rN), for a given fixed value of r1. The
function f satisfies the following properties:
f(r2, ........, rN/r1)dr2.......drN = 1∀r1
(ii) f(r1, ..ri...rj−1, rj+1..., rN/rj) = 0; for i = j; ∀i, j = 1, N
(iii) f(r1, .....ri.., rj...rk−1, rk+1.., rN/rk) = 0; for i = j; ∀i, j 6= k (10)
The property (iii) of Eq.10 assures us that f reflects the fermionic character of an electronic
wavefunction. In fact it says that if any two particles are in the same ’state’ ”r” the
probability of that specific global configuration is zero. In principle, together with condition
(ii), this is a way to mimic the antisymmetric character of the fermionic wavefunction since
for fermions |ψ(r1, ...ri, ...rj, ...rN)|
2 = 0; for i = j, ∀i, j.It must be noticed that condition
(iii) is complementary to (ii). With this formalism the term 〈ψ |T + Vee|ψ〉 can be written
as (see Refs.[9, 10, 12]):
〈ψ |T + Vee|ψ〉 =
|∇ρ(r)|2
, ...., rN/r)|
, ....., rN/r)
....drN
(N − 1)
, ....., rN/r)
|r− r
....drN
dr(11)
where we have identified r1 with r and made use of the property of electron indistinguisha-
bility, thus r could be identified with any of the ri (and the same for r
identified here with
r2) without changing the results; a further consequence is that the Coulomb expression (last
term on the r.h.s.) is written as the sum of N − 1 identical terms for the generic r and r
particles. Using Eq.11 the Levy-Lieb constrained-search formulation can then be written as:
E0 =Minρ
Minf (Γ[f, ρ]) +
|∇ρ(r)|2
v(r)ρ(r)dr
where
Γ[f, ρ] =
, ...., rN/r)|
, ....., rN/r)
....drN
(N − 1)
, ....., rN/r)
|r− r
....drN
dr. (13)
In this way we have transferred the problem from from ψ to f which means that the focus
is now on Γ[f, ρ], i.e., as discussed in Ref.[12], the correlation functional.
III. A PRACTICAL EXAMPLE: THE PARAMETRIC EXPONENTIAL FORM
In our previous work [12], we have proposed an approximation for f based on a two-
particle factorization:
f = ΠNi=2hi(EH(r, ri)) = Π
(N−1)E(r)e−EH (r,ri) (14)
where
e−E(r) =
e−EH (r,ri)dri. (15)
here EH(r, ri) =
ρ(r)ρ(ri)
|r−ri|
, N is the number of particle, and ω the volume corresponding to
one particle. Such an approximation, due to its simplicity, allows us to write an analytic
expression of the Fisher functional which can be used in a straightforward way in numerical
calculations. However it does not satisfy the condition (iii) of Eq.10, and, for this reason,
in order to use it into the Levy-Lieb constrained-search scheme it must be extended. The
expression we propose here is the following:
f(r2, ...rN/r) = Πn=2,Ne
E(r)−γEH (r,rn) × Πi>j 6=1e
−βEH(ri,rj) (16)
with:
e−E(r) =
Πn=2,NΠi>j 6=1e
−γEH (r,rn)−βEH (ri,rj)dr2.....drN (17)
Here γ and β are two free parameters. As it can be easily verified this expression of f
satisfies all the requirements of Eq.10. The meaning of f as expressed in Eq.16 is that the
probability of finding a certain configuration for the N − 1 particles, having fixed particle
r1 = r, depends not only on the fixed particle and its interaction with the N − 1 other
particles as before, but also on the mutual arrangements of the N − 1 particles (it has
also to be kept in mind that using the particle indistinguishability the formalism can be
applied to any ri as a fixed particle). The parameters γ and β express how important
the N − 1 mutual interactions are with respect to the interactions with r. Being now f a
biparametric function, one can use the Levy-Lieb constrained search in our formulation and
find the optimal values for γ and β. This practical example shows two different aspects of
our formulation; basically we have shown that indeed it is possible to build a function f and
actually it can be chosen in a way that its optimal expression can be determined via the
constrained-search formulation. It must be noticed that this form of f is still rather simple
since the spins are not explicitly considered when constructing the function and thus one
cannot distinguish between the exchange and the correlation part of the electron-electron
interaction as it is done in standard Density Functional Theory; as a consequence one should
expect only an overall average description of these two terms which are here incorporated into
the global correlation. However the construction of a more complete expression of f , which
takes care of the effects of the spins, is the subject of current investigation. This emphasizes
once more the merit of the general procedure shown here, that is different expressions of f ,
with different degrees of complexity, can be proposed and their relative validity checked by
the constrained-search procedure.
IV. DISCUSSION AND CONCLUSIONS
As anticipated in the introduction, the consequences of Eqs.12,13 are rather interesting.
The Levy-Lieb variational principle can be reformulated as: The universal functional F [ρ]
is the one with the minimum correlation functional with respect to the electron conditional
probability density. This new interpretation of the HK universal functional tells us that
only an accurate description of the correlation effects, considering the Weizsacker term as
a necessary term, leads to an accurate description of the whole energy functional; such a
criterion is necessary and sufficient. It is obvious that it is necessary; without knowing
Γ[f, ρ], F [ρ] cannot be known; it is sufficient because once Γ[f, ρ] or better f(r2, ...rN/r1)
is (in principle) known than the whole energy functional is known explicitely. Clearly, the
”true” f(r2, ...rN/r1) is very difficult if not impossible to obtain [13], however it can be
sufficiently well described on the basis of mathematical requirements and physical intuition
as done for example in Ref.[12] and as shown in the previous section. From this point of
view, Eqs.12,13, can be seen as an objective criterion to design, on the basis of physical
intuition and fundamental mathematical requirements, valid energy functionals. In fact, as
done in Ref.[12] and in the previous section, one can construct well-founded expressions
for f keeping in mind the physical meaning of the electron correlation effects and the
necessary related mathematical prescriptions of Eq.10. Next one can make use of Eqs.12,13
and choose among different functional forms of f , the one giving the ”smaller” Γ. It
must be noticed that in this work we do not claim that finding a functional form of f is
easier or more rigorous than to find an exchange-correlation functional in standard Density
Functional Theory; it represents an alternative or complementary approach to the latter.
However, the approach based on f allows one to express in a more direct way, via the choice
of different forms of f , the physical principles related to the electron correlation effects and
to have an explicit form of the correlation term for the kinetic functional which is of great
advantage for kinetic functional based methods (see e.g. Refs.[14, 15]).An important aspect
linked to the statement above is that the term: 1
|∇rf(r
,....,rN/r)|
,.....,rN/r)
....drNdr, is
the well known non local Fisher information functional about which a vast literature
is available (see e.g. [10, 16, 17] and references therein); this term is very often linked to
the electron correlation functional and electronic properties(see Refs.[18, 19]), our work
further clarifies this connection, suggesting that the results known from the analysis of
the Fisher functional could be employed in this context. In conclusion we have shown an
alternative view of the Levy-Lieb constrained search approach and provided an example
which clarifies the practical advantage of our idea; in this sense the present work it is not
merely a marginal new formal contribution to a rather well-known method, but gives a new
powerful insight into the field of applicability for realistic systems.
Acknowledgments
I would like to thank Luca Ghiringhelli for a critical reading the manuscript.
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[18] R. F. Nalewajski, Advances in Quantum Chemistry 43, 119, 2003.
[19] R. F. Nalewajski, Chem. Phys. Lett. 386, 265 (2004).
Introduction
The new representation
A practical Example: The parametric exponential form of f
Discussion and Conclusions
References
|
0704.0373 | Reality of linear and angular momentum expectation values in bound
states | Reality of linear and angular momentum expectation values in bound states
Utpal Roy,∗ Suranjana Ghosh,† and T. Shreecharan‡
Physical Research Laboratory, Ahmedabad 380009, India
Kaushik Bhattacharya§
Insituto de Ciencias Nucleares, Universidad Nacional Autónoma de Mexico,
Circuito Exterior, C.U., A. Postal 70-543, C. Postal 04510, Mexico DF, Mexico
ABSTRACT
In quantummechanics textbooks the momentum operator is defined in the Cartesian coordinates and
rarely the form of the momentum operator in spherical polar coordinates is discussed. Consequently
one always generalizes the Cartesian prescription to other coordinates and falls in a trap. In this work
we introduce the difficulties one faces when the question of the momentum operator in spherical polar
coordinate comes. We have tried to point out most of the elementary quantum mechanical results,
related to the momentum operator, which has coordinate dependence. We explicitly calculate the
momentum expectation values in various bound states and show that the expectation value really
turns out to be zero, a consequence of the fact that the momentum expectation value is real. We
comment briefly on the status of the angular variables in quantum mechanics and the problems
related in interpreting them as dynamical variables. At the end, we calculate the Heisenberg’s
equation of motion for the radial component of the momentum for the Hydrogen atom.
I. INTRODUCTION
Quantum mechanics is a treasure house of peculiar and interesting things. Elementary textbooks of quantum
mechanics [1, 2, 3] generally start with the postulates which are required to define the nature of the dynamical
variables in the theory and their commutation relations. The choice of the dynamical variables is not that clear, as
the coordinates in Cartesian system are all elevated to the status of operators where as time remains a parameter.
More over in spherical polar coordinates only the radial component can be represented as an operator while the angles
still remain as a problem. The difficulty of giving different status to the spatial coordinates and time is bypassed in
quantum field theories where all the coordinates and time become parameters of the theory. But the problem with
angles still remain a puzzle which requires to be understood in future.
When we start to learn quantum mechanics, most of the time we begin with elementary calculations relating to the
particle in a one dimensional infinite well, particle in a finite potential well, linear harmonic oscillator and so on. The
main aim of these calculations is to solve the Schrödinger equation in the specific cases and find out the bound state
energies and the energy eigenfunctions in coordinate space representation. While solving these problems we overlook
the subtleties of other quantum mechanical objects as the definition of the momentum operator in various coordinates,
the reality of its expectation value, etc.. In the last one or two decades there has been a number of studies regarding
the self-adjointness of various operators [4]. The aim of these studies has been to analyze the self-adjointness of
various operators like momentum, Hamiltonian etc. and find out whether these operators are really self-adjoint in
some interval of space where the theory is defined, if not then can there be any mathematical method by which we
can make these operators to be self-adjoint in the specified intervals ?
In the present work we deal with a much elementary concept in quantum mechanics related to the reality of the
expectation values of the momentum operator, be it linear or angular. We do not analyze the self-adjointness of
the operators which requires different mathematical techniques. To test the self-adjointness of an operator we have
to see whether the operator is symmetric in a specific spatial interval and the functional domain of the operator
and its adjoint are the same. In the article we always keep in touch with the recent findings from modern research
about self-adjoint extensions but loosely we assume that the operators which we are dealing with are Hermitian. If
something is in contrary we point it out in the main text. A considerable portion of our article deals with the analysis
∗Electronic address: [email protected]
†Electronic address: [email protected]
‡Electronic address: [email protected]
§Electronic address: [email protected]
http://arxiv.org/abs/0704.0373v1
mailto:[email protected]
mailto:[email protected]
mailto:[email protected]
mailto:[email protected]
of the fact that the expectation value of the momentum operator in various bound states are zero, a result which
most of the textbooks only quote but never show. In the simpler cases the result can be shown by one or two lines
of calculation, but in nontrivial potentials as the Morse potential, the Coulomb potential the result is established by
using various properties of the special functions such as the associated Legendre and the associated Laguerre.
The presentation of various materials in our article is done in the following way. Next section deals with the
definition of the momentum operator and its properties. Section III deals with the intricacies of the definition of the
momentum operator in spherical polar coordinates and the problems we face when we try to mechanically implement
the quantization condition, which is invariably always written in Cartesian coordinates in most of the textbooks on
quantum mechanics. In section IV we explicitly calculate the momentum expectation values in various potentials and
show that in bound states we always get the expectation value of the linear momentum to be zero. Section V gives
a brief discussion on the Ehrenfest theorem when we are using it to find out the time derivative of the expectation
value of the radial component of momentum in the case of the Hydrogen atom. We end with the concluding section
which summarizes the findings in our article.
Before going into the main discussion we would like to mention about the convention. We have deliberately put a
hat over various symbols to show that they are operators in quantum mechanics. Some times this convention becomes
tricky when we are dealing with angular variables as there the status of these variables is in question. The other
symbols have their conventional meaning. As we are always using the coordinate representation sometimes we may
omit the hat over the position operator as in this representation the position operator and its eigenvalues can be
trivially interchanged.
II. DEFINITION OF THE MOMENTUM OPERATOR AND THE REALITY OF ITS EXPECTATION
VALUE
From the Poisson bracket formalism of classical mechanics we can infer:
[x̂i, p̂j] = ih̄ δi j , (2.1)
where δi j = 1 when i = j and zero for all other cases, and i, j = 1, 2, 3. In the above equation x̂i is the position
operator and p̂j is the linear momentum operator in Cartesian coordinates. From the above equation we can also find
the form of the momentum operator in position representation, which is:
p̂i = −ih̄
. (2.2)
It is interesting to note that the above expression of the momentum operator also gives us the form of the generator
of translations. This is because of the property:
[p̂x, F (x̂)] = −ih̄
dF (x̂)
, (2.3)
where F (x̂) is an arbitrary well defined function of x̂. The above equation ensures that the momentum operator
generates translations along the x direction.
Particularly in one-dimension the expression of the momentum operator becomes p̂x = −ih̄ ∂∂x . We know that the
expectation value of the momentum operator must be real. If we focus on one-dimensional systems to start with,
where the system is specified by the wave-function ψ(x, t), the expectation value of any operator Ô is defined by:
〈Ô〉 ≡
ψ∗(x, t) Ô ψ(x, t) dx , (2.4)
where ψ∗ signifies complex conjugation of ψ and the extent of the system is taken as −∞ < x <∞. From the above
equation we can write,
〈Ô〉∗ =
ψ(x, t) Ô∗ ψ(x, t)∗ dx . (2.5)
If 〈Ô〉 = 〈Ô〉∗ then the condition of the reality of the expectation value becomes:
ψ∗(x, t) Ô ψ(x, t) dx =
ψ(x, t) Ô∗ ψ(x, t)∗ dx . (2.6)
For a three-dimensional system the above condition becomes,
ψ∗(x, t) Ô ψ(x, t) d3x =
ψ(x, t) Ô∗ ψ(x, t)∗ d3x . (2.7)
Now if we take the specific case of the momentum operator in one-dimension we can explicitly show that its expectation
value is real if the extent of the system is infinite and the wave-function vanishes at infinity. The proof is as follows.
The expectation value of the momentum operator is:
ψ∗(x, t) p̂x ψ(x, t) dx = −ih̄
ψ∗(x, t)
∂ψ(x, t)
= −ih̄
ψ∗(x, t)ψ(x, t)|∞−∞ −
ψ(x, t)
∂ψ(x, t)∗
, (2.8)
If the wave functions vanish at infinity then the first term on the second line on the right-hand side of the above
equation drops and we have,
ψ∗(x, t) p̂x ψ(x, t) dx = −ih̄
ψ∗(x, t)
∂ψ(x, t)
= ih̄
ψ(x, t)
∂ψ(x, t)∗
ψ(x, t) p̂∗x ψ(x, t)
∗ dx . (2.9)
A similar proof holds for the three-dimensional case where it is assumed that the wave-function vanishes at the
boundary surface at infinity.
III. THE EXPECTATION VALUE OF THE MOMENTUM OPERATOR IN CARTESIAN AND
SPHERICAL POLAR COORDINATES
In non-relativistic version of quantum mechanics we know that if we have a particle of mass m which is present in
a time-independent potential we can separate the Schrödinger equation:
∂ψ(x, t)
∇2 + V (x)
ψ(x, t) , (3.1)
into two equations, one is the time-dependent one which gives the trivial solution e−
h̄ where E is the total energy
of the particle, and the other equation is the time-independent Schrödinger equation:
∇2u(x) + 2m
(E − V (x))u(x) = 0 , (3.2)
where u(x) is the solution of the time-independent Schrödinger equation and the complete solution of the Eq. (3.1)
ψ(x, t) = u(x)e−
h̄ . (3.3)
In the case of the free-particle, where V (x) = 0, we have u(x, t) = eik·x where E = k
and k = |k|. The free-
particle solution is an eigenfunction of the momentum operator with eigen value h̄k. Although if we try to find
out the expectation value of the momentum operator as is done in the last section we will be in trouble as these
wave-functions do not vanish at infinity, a typical property of free-particle solutions. But this problem is not related
to the Hermiticity property of the momentum operator, it is related with the de-localized nature of the free-particle
solution.
In physics many times we require to solve a problem using curvilinear coordinate systems. The choice of our
coordinate system depends upon the specific symmetry which we have at hand. Suppose we are working in spherical
polar coordinates and the solution of Eq. (3.2) can be separated into well behaved functions of r, θ and φ as,
u(x) = u(r, θ, φ) = R(r)Θ(θ)Φ(φ) . (3.4)
If we try to follow the proof of the Hermiticity of the linear momenta components, as done in the last section, in
spherical polar coordinates, then we should write:
〈p̂〉 = −ih̄
u∗(r, θ, φ)∇u(r, θ, φ) dτ ,
= −ih̄
R∗(r)Θ∗(θ)Φ∗(φ)
r sin θ
R(r)Θ(θ)Φ(φ)r2drdΩ ,
(3.5)
where in the above equation er, eθ, eφ respectively are the unit vectors along r, θ and φ, and dΩ = sin θ dθ dφ. τ is
the volume over which we integrate the expression in the above equation. From the last equation we can write:
〈p̂r〉 = −ih̄
|Θ(θ)|2|Φ(φ)|2dΩ
r2R∗(r)
dR(r)
dr , (3.6)
As Θ(θ) and Φ(φ) are normalized, the integration:
|Θ(θ)|2|Φ(φ)|2dΩ = 1 and we can proceed as in Eq. (2.9) as:
〈p̂r〉 = −ih̄
r2R∗(r)
dR(r)
= −ih̄
r2R∗(r)R(r)
2rR∗(r) + r2
dR∗(r)
R(r) dr
. (3.7)
If R(r) vanishes at infinity then the above equation reduces to,
〈p̂r〉 =
r2R(r)
dR∗(r)
+ 2ih̄
r|R(r)|2 dr ,
= 〈p̂r〉∗ + 2ih̄
r|R(r)|2 dr . (3.8)
The above equation implies that 〈p̂r〉 is not real in spherical polar coordinates. The solution of the above problem
lies in redefining p̂r as is evident from Eq. (3.8), and it was given by Dirac [5, 6]. The redefined linear momentum
operator along r can be:
p̂r ≡ −ih̄
= −ih̄
r . (3.9)
This definition of the p̂r is suitable because in this form it satisfies the commutation relation as given in Eq. (2.1)
where now the operator conjugate to r̂ is p̂r. The form of p̂r in Eq. (3.9) shows that for any arbitrary function of r as
F (r) we must still have Eq. (2.3) satisfied. This implies that the modified form of p̂r is still a generator of translations
along the r direction. Up to this point we were following what was said by Dirac regarding the status of the radial
momentum operator. Still everything is not that smooth with the redefined operator as we can see that it turns out
to be singular around r = 0, more over, although the radial momentum acts like a translation generator along r but
near r = 0 it cannot generate a translation towards the left as the interval ends there.
In this regard we can state that the issue of the reality of the radial component of the momentum in spherical polar
coordinates is a topic of modern research in theoretical physics [7, 8]. It has been shown that the operator −ih̄ ∂
not Hermitian and more over it can be shown [4] that such an operator cannot be self-adjoint in the interval [0,∞].
In some recent work [8] the author claims that there can be an unitary operator which connects −ih̄ ∂
to −ih̄1
and as the former operator does not have a self-adjoint extension in the semi-infinite interval so the latter is also not
self-adjoint in the same interval.
If we further try to find out whether 〈p̂θ〉 and 〈p̂φ〉 are real, then we will face difficulties. Working out naively if
we claim that p̂φ =
rsinθ
as suggested by the φ component of Eq. (3.5) we will notice that φ p̂φ does not have the
dimension of action. This means p̂φ or p̂θ is not conjugate to φ or θ. This is a direct representation of the special
coordinate dependence of the quantization condition. Only in Cartesian coordinates the variables conjugate to x, y
and z are px, py and pz. Taking the clue from classical mechanics we know the proper dynamical variables conjugate
to φ̂ and θ̂ are the angular momentum operators, namely L̂θ and L̂φ. In general L̂φ is given by:
L̂φ = −ih̄
, (3.10)
which can be shown to posses real expectation values by following a similar proof as is done in Eq. (2.8) and Eq. (2.9),
if we assume Φ(0) = Φ(2π). In this form it is tempting to say that we can have a relation of the form,
[φ̂, L̂φ] = ih̄ , (3.11)
which looks algebraically correct. But the difficulty in writing such an equation is in the interpretation of φ̂ which
has been elevated from an angular variable to a dynamical operator. In spherical polar coordinates both θ and φ
are compact variables and consequently have their own subtleties. Much work is being done in trying to understand
the status of angular variables and phases [9, 10], in this work we only present one example showing the difficulty of
accepting φ̂ as an operator.
From the solution of the time-independent Schrödinger equation for an isotropic potential we will always have:
Φ(φ) =
eiMφ , (3.12)
where M = 0,±1,±2, ·, ·. Now if φ̂ is an operator we can find its expectation value, and it turns out to be:
〈φ̂〉 = 1
φeiMφe−iMφ dφ ,
= π , (3.13)
and the expectation value of φ̂2 is:
〈φ̂2〉 = 1
φ2eiMφe−iMφ dφ ,
π2 . (3.14)
Consequently ∆φ =
〈φ̂2〉 − 〈φ̂〉2 = π√
. Similarly calculating L̂φ we get:
〈L̂φ〉 =
e−iMφeiMφ dφ ,
= Mh̄ , (3.15)
as expected, and 〈L̂2φ〉 = M2h̄
. This implies ∆Lφ =
〈L̂2φ〉 − 〈L̂φ〉2 = 0. So we can immediately see that the
Heisenberg uncertainty relation between φ̂ and L̂φ, ∆φ∆Lφ ≥ h̄/2 breaks down. This fact makes life difficult and we
have no means to eradicate this problem.
Taking the clue from the φ part we can propose that L̂θ is also of the form −ih̄ ∂∂θ . With this definition of L̂θ let
us try to prove its Hermitian nature as done in Eq. (3.7). Taking R(r) and Φ(φ) in Eq. (3.4) separately normalized,
we can write:
〈L̂θ〉 = −ih̄
Θ∗(θ)
dΘ(θ)
sin θdθ
= −ih̄
sin θΘ∗(θ)Θ(θ)|π0 −
cos θΘ∗(θ) + sin θ
dΘ∗(θ)
Θ(θ) dθ
sin θΘ(θ)
dΘ∗(θ)
+ ih̄
cos θ |Θ(θ)|2 dθ ,
= 〈L̂θ〉∗ + ih̄
cos θ |Θ(θ)|2 dθ . (3.16)
The above equation shows that 〈L̂θ〉 is not real. The rest is similar to the analysis following Eq. (3.8) where now we
have to redefine the angular momentum operator conjugate to θ as [11]:
L̂θ ≡ −ih̄
cot θ
. (3.17)
Unlike the φ case, Θ(θ) are not eigenfunctions of L̂θ. But the difficulties of establishing θ as an operator still persists
and in general θ is not taken to be a dynamical operator in quantum mechanics.
It is known that both θ and φ are compact variables, i.e. they have a finite extent. But there is a difference between
them. In spherical polar coordinates the range of φ and θ are not the same, 0 ≤ φ < 2π and 0 ≤ θ ≤ π. This difference
can have physical effects. As φ runs over the whole angular range so the wave-function corresponding to it Φ(φ) is
periodic in nature while due to the range of θ, Θ(θ) need not be periodic. Consequently there can be a net angular
momentum along the φ direction while there cannot be any net angular momentum along θ direction. And this can
be easily shown to be true. As the time-independent Schrödinger equation for an isotropic potential yields Φ(φ) as
given in Eq. (3.12) similarly it is known that in such a potential the form of Θ(θ) is given by:
Θ(θ) = Nθ P
M (cos θ) , (3.18)
where Nθ is a normalization constant depending on L, M and P
M (cos θ) is the associated Legendre function, which
is real. In the above equation L and M are integers where L = 0, 1, 2, 3, ·, · and M = 0,±1,±2,±3, ·, ·. The
quantum number M appearing in Eq. (3.12) and in Eq. (3.18) are the same. This becomes evident when we solve
the time-independent Schrödinger equation in spherical polar coordinates by the method of separation of variables.
A requirement of the solution is −L ≤ M ≤ L. Now we can calculate the expectation value of L̂θ using the above
wave-function and it is:
〈L̂θ〉 = −ih̄N2θ
PLM (cos θ)
dPLM (cos θ)
cot θPLM (cos θ)
sin θdθ
= −ih̄N2θ
PLM (cos θ)
dPLM (cos θ)
sin θdθ +
PLM (cos θ)P
M (cos θ) cos θ dθ
(3.19)
To evaluate the integrals on the right hand side of the above equation we can take x = cos θ and then the expectation
value becomes:
〈L̂θ〉 = −ih̄N2θ
PLM (x)
dPML (x)
(1− x2) 12 dx
PLM(x)P
1− x2
. (3.20)
The second term in the right hand side of the above equation vanishes as the integrand is an odd function in the
integration range. For the first integral we use the following recurrence relation [12]:
(x2 − 1)
dPLM (x)
=MxPLM (x) − (L+M)PLM−1(x) , (3.21)
the last integral can be written as,
〈L̂θ〉 = ih̄N2θ
x(1− x2)−
2PLM (x)P
M (x) dx
− (L+M)
(1 − x2)−
2PLM (x)P
M−1(x) dx
. (3.22)
PLM (x) = (−1)L+MPLM (−x) , (3.23)
we can see immediately that both the integrands in the right hand side of the above equation is odd and consequently
〈L̂θ〉 = 0 as expected. A similar analysis gives 〈L̂φ〉 = Mh̄. It must be noted that the form of L̂θ still permits it to
be the generator of rotations along the θ direction.
As the motion along φ is closed so there can be a net flow of angular momentum along that direction but because
the motion along θ is not so, a net momentum along θ direction will not conserve probability and consequently for
probability conservation we must have expectation value of angular momentum along such a direction to be zero.
In elementary quantum mechanics text books it is often loosely written that the solution of the time-independent
Schrödinger equation is real when we are solving it for a real potential. But this statement is not correct. The reality
of the solution also depends upon the coordinate system used. Specially for compact periodic coordinates we can
always have complex functions as solutions without breaking any laws of physics.
Before leaving the discussion on angular variables in spherical polar coordinates we want to point out one simple
thing which is interesting. In Cartesian coordinates when we deal with angular momentum we know that:
[L̂i, L̂j] = iǫijk L̂k , (3.24)
where L̂i stands for L̂x, L̂y or L̂z. For this reason there cannot be any state which can be labelled by the quantum
numbers of any two of the above angular momenta. But from the expressions of L̂φ and L̂θ we see that,
[L̂φ, L̂θ] = 0 , (3.25)
and consequently in spherical polar coordinates we can have wave-function solutions of the Schrödinger equation
which are simultaneous eigenfunctions of both L̂φ and L̂θ as P
M (θ).
For real V (x), we expect the solution of the time-independent Schrödinger equation u(x) to be real, when we
are solving the problem in Cartesian coordinates. In all these cases the expectation value of the linear momentum
operators must vanish. The reason is simple and can be understood in one-dimensional cases where with real u(x) we
directly see that the integral
∂u(x)
dx is real and so
∗(x) p̂x u(x) dx becomes imaginary as p̂x contains
i, as is evident from the first line in Eq. (2.9). So if the expectation value of the momentum operator has to be
real then the only outcome can be that for all those cases where we have a time-independent solution in a bounded
region of space, with a real potential and working in Cartesian coordinates, the expectation value of the momentum
operator must vanish. The above statement is true in curvilinear coordinates also, but in those cases the definition of
the momentum operators have to be modified. This fact becomes clear when we write the relationship between the
probability flux and the expectation value of the momentum operator. The probability flux for a particle of mass m
j(x, t) = −
[ψ∗(x, t)∇ψ(x, t) − (∇ψ∗(x, t))ψ(x, t)] ,
Im (ψ∗(x, t)∇ψ(x, t)) , (3.26)
where ‘Im’ implies the imaginary part of some quantity. Most of the elementary quantum mechanics books then
proceeds to show that:
d3x j(x, t) =
, (3.27)
which is obtained from Eq. (3.26) by integrating both sides of it over the whole volume. From Eq. (3.26) we immediately
see that if the solution of the time-independent Schrödinger equation is real we will have j(x, t) = 0 and consequently
from Eq. (3.27), 〈p̂〉 = 0. But this statement is also coordinate dependent, which is rarely said in elementary textbooks
of quantum mechanics. Eq. (3.26) evidently does not hold in spherical polar coordinates. If we take Eq. (3.4) as the
solution in a general isotropic central potential and use the general form of ∇ in spherical polar coordinates then it
can be seen that jr(r, θ, φ, t) = 0 for a real potential. But then Eq. (3.27) does not hold as here p̂r is simply the radial
component of ∇ and not as given in Eq. (3.9), and we know 〈 d
〉 is not zero. The reason why Eq. (3.26) is not suitable
in spherical polar coordinates is related to the fact that in deriving Eq. (3.26) one assumes that the probability density
of finding the quantum state within position x and x + dx at time t is |ψ(x, t)|2. But this statement is only true in
Cartesian coordinates, in spherical polar coordinates the probability density of the system to be within a region r and
r + dr, θ and θ + dθ, φ and φ+ dφ is not |ψ(r, θ, φ)|2 but |ψ(r, θ, φ)|2r2 sin θ and consequently the steps which follow
leading to Eq. (3.26) in Cartesian coordinates are not valid in spherical polar coordinates. In general Eq. (3.26) will
not be valid in any curvilinear coordinate system.
The next section contains the actual calculations of the expectation values of the momentum operator in various
cases where we have bound state solutions. In all the relevant cases discussed in this article it is seen that although
〈p̂x〉 = 0 but 〈p̂2x〉 is not zero as it is related to the Hamiltonian operator. In all the cases we must have,
〈(p̂x)s〉 = 0 , s = odd integer . (3.28)
The above equation can be guessed from the reality of the expectation value of the momentum operator.
IV. MOMENTUM EXPECTATION VALUES IN VARIOUS BOUND STATES
In this section we will calculate the momentum expectation values in various bound states with stiff or slowly
varying potentials.
A. Particle in one-dimensional stiff potential wells
1. Infinite square well potential
In this case we consider a particle to be confined in region −L
along the x-axis where the potential is specified
V (x̂) = ∞ , |x| ≥
= 0 , |x| <
. (4.1)
In this case the solution of the time-independent Schrödinger equation, Eq. (3.2), satisfies the boundary condition,
= 0 , (4.2)
and as the potential has parity symmetry about x = 0 we have two sets of solutions, the odd solutions:
u(o)n (x) =
, (4.3)
and the even solutions:
u(e)n (x) =
(2n− 1)πx
. (4.4)
In the above equations n is a positive integer. Both of these functions, u
n (x) for the odd case and u
n (x) for the
even case, are real and are not momentum eigenstates. But the momentum expectation values can be found out from
the above solutions. For the odd solutions we have:
〈p̂x〉 = −ih̄
u(o)n (x)
n (x)
= −4inπh̄
= 0 , (4.5)
as expected. Similarly for the even solutions it is also easy to show that the expectation value of the momentum
operator vanishes.
2. Finite square well potential
In this case,
V (x̂) = 0 , |x| ≥ a ,
= −V0 , |x| < a , (V0 > 0) . (4.6)
If we are not interested in the normalization constant of the bound state solution then the solution of the time-
independent Schrödinger equation in this case is:
u(x) ∼ e−κ|x| , |x| > a ,
∼ cos(kx) , |x| < a , (even parity)
∼ sin(kx) , |x| < a , (odd parity) , (4.7)
where,
2m(−|E|+ V0)
, (4.8)
2m|E|
. (4.9)
In this case the expectation value of the momentum operator is:
〈p̂x〉 ∼ −ih̄
du(x)
∼ −ih̄
e2κxdx−
e−2κxdx
sin(kx) cos(kx) dx
= 0 , (4.10)
where the first two lines of the above equation holds up to a constant arising from the normalization of the wave-
function. In deriving the last equation we have taken the odd parity solution, but the result remains unaffected if we
take the even parity solution as well.
3. Dirac-delta potential
In this case the potential is:
V (x̂) = −V0 δ(x̂) , (V0 > 0) . (4.11)
In this case there can be one bound state solution which is obtained after solving the Eq. (3.2). Demanding that the
solution u(x) satisfies the boundary conditions:
u(x = −ǫ) = u(x = +ǫ) , (4.12)
= −2mV0
u(x = 0) , (4.13)
where ǫ is an infinitesimal quantity tending to zero, we get the form of the solution which is:
u(x) =
κ eκx , x ≤ 0 , (4.14)
κ eκx , x ≥ 0 , (4.15)
where κ = mV0
and the energy of the bound state is E = −mV
The expectation value of the momentum operator in this case is:
〈p̂x〉 = −ih̄
du(x)
= −ih̄κ
e2κx dx −
e−2κx dx
= 0 . (4.16)
In this case, also from Hermiticity of the momentum operator we see that Eq. (3.28) holds true.
B. Particle in one-dimensional slowly varying potentials
1. Linear harmonic oscillator potential
In the case of the linear harmonic oscillator we have:
V (x̂) =
mω2x̂2 , (4.17)
where ω is the angular frequency of the oscillator. The solution of Eq. (3.2) in this case, using the series solution
method, yields:
un(q) = Nn e
2 Hn(q) , (4.18)
where n = 0, 1, 2, ·, · and q =
αx where α = mω
. Hn(q) are Hermite polynomials of order n and Nn is the
normalization constant given by,
π n! 2n
. (4.19)
The momentum expectation value in this case turns out to be,
〈p̂x〉 = −ih̄
un(q)
dun(q)
= −ih̄
Hn(q)
dHn(q)
q e−q
H2n(q) dq
= 0 . (4.20)
The first integral on the right side of the second line of the last equation vanishes because,
dHn(q)
= 2nHn−1(q) and
consequently the integral transforms into the orthogonality condition of the Hermite polynomials. The second integral
on the second line of the right side of the above equation vanishes because the integrand is an odd function of q.
The linear harmonic oscillator (LHO) has some very interesting properties. To unravel them we have to digress a
bit from the wave-mechanics approach which we have been following and follow the Dirac notation of bra and kets.
The Hamiltonian of the LHO in one-dimension is:
mω2x̂2 , (4.21)
which can also be written as:
Ĥ = h̄ω
â†â+
, (4.22)
where â and ↠are the annihilation and the creation operators given by:
, ↠≡
x̂− ip̂x
. (4.23)
It can be seen clearly from the above definitions that â is not an Hermitian operator. More over from the definition
of the operators we see that,
[â, â†] = 1 . (4.24)
Conventionally the number operator is defined as:
N̂ ≡ â†â , (4.25)
and its eigen-basis are the number states |n〉 such that,
N̂ |n〉 = n|n〉 . (4.26)
The Hamiltonian of the LHO can be written in terms of the number operator and consequently the number states
are energy eigenstates. In this basis the action of the annihilation and creation operators are as:
â|n〉 =
n |n− 1〉 , (4.27)
â†|n〉 =
n+ 1 |n+ 1〉 . (4.28)
From the definitions of the annihilation and creation operators we can write the momentum operator as:
p̂x = −i
(â− â†) . (4.29)
From Eq. (4.27), Eq. (4.28) and the above equation we can write the matrix elements of the momentum operator as:
〈n′|p̂x|n〉 = i
n δn′, n−1 +
n+ 1 δn′, n+1
. (4.30)
The above equation shows that the momentum operator can connect two different energy eigenstates.
In the case of LHO, except the number operator states, we can have another state which is an eigenstate of the
annihilation operator â. This state is conventionally called the coherent state and it is given as:
|α〉 = e−
|n〉 , (4.31)
where α is an arbitrary complex number. Now from Eq. (4.29) we can find the momentum expectation value of the
coherent state and it is,
〈p̂x〉α ≡ 〈α|p̂x|α〉 =
Im(α) . (4.32)
From the above equation we can see that although the expectation value of the momentum operator is zero in the
energy eigen-basis but it is not so when we compute the momentum expectation value in the coherent state basis,
which is essentially a superposition of energy eigenstates. It must be noted that the momentum expectation value is
non zero only when the parameter α has an imaginary part.
2. Pöschl-Teller potential
Among the potentials belonging to the hypergeometric class the Pöschl-Teller potentials have been the most ex-
tensively studied and used. This class of potentials consist of trigonometric as well as the hyperbolic type. The
trigonometric versions have found applications in molecular and solid state physics and the hyperbolic variants have
been used in various studies related to black hole perturbations.
In the present work we use the trigonometric, symmetric Pöschl-Teller potential given by:
V (x̂) = V0 tan
2(ax̂) , (4.33)
where V0 can be parameterized as:
λ(λ− 1) , (4.34)
with for some positive number λ > 1 and a is some scaling factor. The energy eigenvalues of the bound state solutions
En = −
h̄2a2
(n2 + 2nλ+ λ) , (4.35)
and the solution of the time-independent Schrödinger equation is,
un(x) = Nn
cos(ax)P
1/2−λ
n+λ−1/2 (sin(ax)) , (4.36)
where,
a(n+ λ)Γ(n + 2λ)
Γ(n+ 1)
, (4.37)
is the normalization constant and Pµν (x) is the associated Legendre function. At this point it is fair to point out that
Pµν (x) is not the Legendre polynomial P
M (x) appearing in Eq. (3.18), as µ and ν need not be integers as L and M .
Pµν (x) is not a polynomial but the function appearing in the right hand side of Eq. (4.36) is a polynomial.
Now as claimed in the text let us show that the momentum expectation value is indeed zero. Before we proceed
let us simplify the notation a bit by calling µ = 1/2− λ and ν = n+ λ − 1/2. Substituting z = ax we can write the
momentum expectation value as:
〈p̂x〉 = −ih̄N2n
∫ π/2
cos(z)Pµν (sin(z))
cos(z)Pµν (sin(z))
. (4.38)
Note the limits of the integration range from π/2 to −π/2 since at this value the potential becomes infinity hence we
need not consider the integration range to be the whole real line. For the sake of convenience let us make a change
of variable; letting y = sin(z) the above integral becomes:
〈p̂x〉 = −ih̄N2n
dy (1− y2)1/4Pµν (y)
(1 − y2)1/4Pµν (y)
. (4.39)
Taking the derivative inside the integral we get:
〈p̂x〉 = −ih̄N2n
(1− y2)1/2Pµν (y)
dPµν (y)
− y(1− y
2)−1/2
Pµν (y)P
ν (y)
. (4.40)
It is known that for associated Legendre functions [13],
Pµν (−x) = cos[(µ+ ν)π]Pµν (x)−
sin[(µ+ ν)π]Qµν (x) , (4.41)
where Qµν (x) is the other linearly independent solution of the associated Legendre differential equation. As in our
case µ+ ν = n so Pµν (x) will have definite parity. As P
ν (x) has definite parity so the contribution of the second term
in the above integral vanishes since the total integrand is an odd function. The first integral is similar to the one in
Eq. (3.20) and, due to the typical parity property of Pµν (x) as shown in Eq. (4.41), it also vanishes. Consequently we
have 〈p̂x〉 = 0 as expected.
3. Morse potential
Diatomic molecule is an exactly solvable system, if one neglects the molecular rotation. The most convenient model
to describe the system, is the Morse potential [14]:
V (x̂) = D(e−2βx̂ − 2e−βx̂) , (4.42)
where x = r/r0 − 1, which is the distance from the equilibrium position scaled by the equilibrium value of the
inter-nuclear distance r0. D is the depth of the potential, called dissociation energy of the molecule and β being a
parameter which controls the width of the potential.
In terms of the above scaled variable x, the time-independent Schrödinger equation becomes:
d2u(x)
+D(e−2βx − 2e−βx)u(x) = Eu(x) . (4.43)
Here µ is the reduced mass of the molecule and the corresponding bound state eigen function comes out to be:
uλn(ξ) = Ne
−ξ/2ξs/2Lsn(ξ) , (4.44)
where the variables are described as,
ξ = 2λe−y; y = βx; 0 < ξ <∞ , (4.45)
n = 0, 1, ..., [λ− 1/2] , (4.46)
which is nothing but the quantum number of the vibrational bound states. Here [ρ] denotes the largest integer smaller
than ρ, thus total number of bound states is [λ− 1/2] + 1. The parameters,
2µDr20
and s =
− 8µr
E , (4.47)
satisfy the constraint condition s+2n = 2λ− 1. We note that the parameter λ is potential dependent and s is related
to energy E. In Eq. (4.44), Lsn(y) is the associated Laguerre polynomial and N is the normalization constant [15]:
β(2λ− 2n− 1)Γ(n+ 1)
Γ(2λ− n)r0
. (4.48)
We are looking for the expectation value of linear momentum for a vibrating diatomic molecule, and its expression
〈p̂x〉 = −ih̄
u∗n(ξ)
un(ξ)dx . (4.49)
In terms of the changed variable ξ = 2λe−βx the integration limit changes to ∞ to 0 and the expectation value
becomes:
〈p̂x〉 = −ih̄
u∗n(ξ)
un(ξ)dξ
= ih̄N2
e−ξξs(Lsn(ξ))
2dξ +
e−ξξs−1(Lsn(ξ))
e−ξξsLsn(ξ)
Lsn(ξ)dξ
= ih̄N2
I2 + I3
. (4.50)
Integral I1 is the orthogonality relation of the associated Laguerre polynomials, which is:
e−ξξsLsn(ξ)L
m(ξ)dξ =
Γ(s+ n+ 1)
Γ(n+ 1)
δm,n . (4.51)
To evaluate the second integral one uses the normalization integral of Morse eigenstates. The normalization relation
u∗(ξ)u(ξ)dr =
|N |2r0
e−ξξs−1(Lsn(ξ))
2dξ = 1 . (4.52)
The above integral involving ξ, is explicitly I2. N , being the normalization constant as given in Eq. 4.48. Thus it is
very straight forward to evaluate I2 from the above relation as,
Γ(n+ s+ 1)
s Γ(n+ 1)
. (4.53)
The last integrand I3 includes a differentiation which can be written as [16]:
Lsn(ξ) = −Ls+1n−1(ξ) . (4.54)
Writing the right hand side of the above equation as a summation [17]:
Ls+1n =
Lsm , (4.55)
and substituting the derivative term in integral I3 we obtain:
I3 = −
e−ξξsLsn(ξ)L
m(ξ)dξ . (4.56)
In the above integral m 6= n because m can go only upto (n− 1). Thus the integral vanishes. Now let us see what is
the expectation value of momentum observable, after evaluating the three integrals above. Substituting the non-zero
values I1 and I2 in Eq. 4.50, it is clear that the expectation value of momentum is zero as has been expected.
C. Position expectation values for various potentials
After a thorough discussion about the momentum expectation values for various solvable one-dimensional potentials,
it is worth spending some time discussing about the average position of the particle inside the bound states. Among
all the above examples, in each case we had V (x) = V (−x) except the Morse potential as Morse potential is not an
example of a symmetric potential: V (x) 6= V (−x).
In deriving the expectation values of momentum for above symmetric cases, we often considered that the integrals of
odd functions over the symmetric limits vanishes. This result does not hold true for the asymmetric Morse potential.
Already we have shown that the momentum expectation value: < p >= 0 for all the above potentials. When it
comes to the expectation values of position, one can easily see that < x >= 0 for symmetric potentials whose centers
are at the origin. On the other hand if this is not the case, suppose the infinite square well is defined in the range
0 ≤ x ≤ L also then the expectation value of position does not vanish. It becomes L/2. Thus, more accurately the
average position of the particle is dependent on the symmetry of the potential where as the average momentum is
solely guided by the reality of it’s eigenvalues and consequently it is zero always.
Below we will briefly discuss how the asymmetry of the potential affects the expectation value of x in the case of
the Morse potential. The expectation value of the position operator is:
〈x̂〉 =
uλ∗n (ξ)xu
n(ξ)dx. (4.57)
The eigen function and the variables are respectively substituted from Eq. (4.44) and Eq. (4.45). We obtain
〈x̂〉 = N
ln(2λ)
e−ξξs−1(Lsn(ξ))
2dξ +
e−ξξs−1(Lsn(ξ))
2 ln(ξ)dξ
. (4.58)
The first integral is already been obtained in Eq. (4.53). This result is independent of the quantum number n. The
second integral (say I) is not that straight forward, because it contains associated Laguerre polynomial, logarithm,
exponential and monomial functions. Here at best we can evaluate the integral atleast for some specific n as, n = 0
or n = 1, when the Laguerre polynomial is respectively replaced by 1 and (−ξ + s + 1). For the ground state wave
function (n = 0), I would be
In=0 =
e−ξξs−1 ln(ξ)dξ, (4.59)
which can be written in terms of Ψ(s) and Γ function [18]:
In=0 = Γ(s)Ψ(s), (4.60)
where, Ψ(s) is the logarithmic factorial function, defined as
d(ln(s)!
= Ψ(s). For n = 0, first integral reduces
to Γ(s) from Eq. (4.53). Above two evaluations gives the ground state expectation value:
〈x̂〉n=0 =
[ln(s+ 1)−Ψ(s)] . (4.61)
For n = 1, one can proceed in the same way
〈x̂〉n=1 =
e−ξξs+1 ln(ξ)dξ + (s+ 1)2
e−ξξs−1 ln(ξ)dξ2(s+ 1)
e−ξξs ln(ξ)dξ
(4.62)
Γ(s+ 2)Ψ(s+ 2) + (s+ 1)2Γ(s)Ψ(s)− 2(s+ 1)Γ(s+ 1)Ψ(s+ 1)
which simplifies to give the expectation value corresponding to the second eigen state:
〈x̂〉n=1 =
ln(s+ 3)−Ψ(s+ 2) + 3
(s+ 2)
. (4.63)
Other expectation values for n > 1 can also be obtained in a similar fashion.
The important point which is to be noted here is, though the average momentum vanishes, the average position is
non-zero for Morse potential and remain so, irrespective of the choice of coordinate origin. This result is also true for
all eigen states of the same Hamiltonian.
D. Momentum expectation value for a three-dimensional slowly varying spherically symmetric potential
In three dimensions, for a spherically symmetric potential the solution of the Schrödinger equation is given in
Eq. (3.4). Here we have assumed that the variables can be separated. The expectation values of L̂θ and L̂φ have been
evaluated in section III. In this section we take the case of the Hydrogen atom and calculate the expectation value of
the radial component of the linear momentum.
1. The Hydrogen atom
In this case,
V (r̂) = −e
. (4.64)
where e is the electronic charge and r =
x2 + y2 + z2. Now we have to write Eq. (3.2) in spherical polar coordinates
and the solution of the time-independent Schrödinger equation is:
unLM (r, θ, φ) = Nr RnL(r)YLM (θ, φ) ,
= Nr e
−r/na0
L2L+1n−L−1
YLM (θ, φ) , (4.65)
where a0 =
is the Bohr radius and m is the reduced mass of the system comprising of the proton and the electron.
n is the principal quantum number which is a positive integer, L2L+1n−L−1(x) are the associated Laguerre polynomials,
YLM (θ, φ) are the spherical-harmonics, and Nr is the normalization arising from the radial part of the eigenfunction.
The values which L and M can take is discussed in section III. The radial normalization constant is given by:
(n− L− 1)!
(n+ L)!2n
. (4.66)
The spherical-harmonics are given by,
YLM (θ, φ) =
2L+ 1
(L −M)!
(L +M)!
PLM (cos θ)e
iMφ , (4.67)
where PML (cos θ) are the associated Legendre functions. It is noted that although the Coulomb potential is a real
potential but the solution in spherical polar coordinates is not real, eiMφ, is complex. The spherical-harmonics are
ortho-normalized according to the relation,
dθ dφ sin θ YLM (θ, φ)YL̃ M̃ (θ, φ) = δLL̃ δMM̃ . (4.68)
Let us write the eigenfunctions in terms of dimensionless quantity: ρ = 2r/na0 ≡ αr. Also we define k ≡ (2L + 1)
and nr ≡ (n − L − 1) for the sake of convenience. With this amount of notational machinery the eigenfunctions can
be written as:
unLM (r, θ, φ) = Nr RnL(ρ)YLM (θ, φ) . (4.69)
The radial momentum expectation value in this case is not given by −ih̄〈 ∂
〉, its form is (already discussed in section
III):
〈p̂ρ〉 = −ih̄Ñ2
dρ ρ2R∗nL(ρ)
RnL(ρ)
dΩ [YLM (θ, φ)]
2 . (4.70)
Where Ñ2 = N2r /α
2. The integral for the spherical harmonics yields identity. The radial expectation value then
becomes,
〈p̂ρ〉 = −ih̄Ñ2
e−ρρk+1[Lknr (ρ)]
2 + (L+ 1) e−ρρk[Lknr (ρ)]
2 + e−ρρk+1Lknr (ρ)
[Lknr (ρ)]
. (4.71)
Using the recurrence relation [16]:
Lknr (ρ) = ρ
−1 [nr Lknr (ρ)− (nr + k)L
nr−1(ρ)
, (4.72)
the expectation value integral acquires the form:
〈p̂ρ〉 = −ih̄Ñ2
e−ρρk+1[Lknr (ρ)]
2 + (nr + L+ 1) e
−ρρk[Lknr(ρ)]
2 + e−ρρkLknr (ρ)L
nr−1(ρ)
. (4.73)
The third contribution of the becomes zero from the orthogonality property of the associated Laguerre polynomials
as given in Eq. (4.51). The contribution from the second term can also be found similarly. To find the share of the
first term we make use of [19]:
dρ e−ρ ρk+1 [Lknr (ρ)]
(nr + k)!
(2nr + k + 1). (4.74)
Collecting all the contributions we get the radial expectation value to be zero as expected.
V. A DISCUSSION ON HEISENBERG’S EQUATION OF MOTION AND EHRENFEST THEOREM
The time evolution of any operator Ô in the Heisenberg picture is given by:
[Ô, Ĥ ] , (5.1)
where Ĥ is the Hamiltonian of the system. The Hamiltonian of a quantum system comprising of a particle of mass
m is given by:
+ V (x̂) . (5.2)
From the above two equations we can write the time evolution of the momentum operator in one dimension, in
Cartesian coordinates as:
[p̂x, Ĥ] = −
V (x̂) , (5.3)
which is the operator version of Newton’s second law in a time independent potential. Now if we take the expectation
values of both sides of Eq. (5.3) in any basis we get:
d〈p̂x〉
V (x̂)
, (5.4)
and historically the above equation is called the Ehrenfest theorem, which was deduced in a different way by P.
Ehrenfest. Using the Ehrenfest theorem we can deduce that the rate of change of the expectation value of the
momentum operator is zero in the case of the linear harmonic oscillator. In the case of the linear harmonic oscillator
we have:
V (x̂) = mω2x̂ , (5.5)
and it can be trivially shown that 〈x̂〉 = 0. This directly implies that,
d〈p̂x〉
= 0 , (5.6)
for the linear harmonic oscillator. The above equation shows that the expectation value of the momentum along x
direction is constant, and this constant is zero is known from other sources.
Next we focus on the Hydrogen atom. The Hamiltonian of the Hydrogen atom is:
Ĥ = − h̄
2 − e
, (5.7)
where,
2 = −h̄2
sin θ
sin θ
sin2 θ
, (5.8)
whose eigenvalues are of the form h̄
L(L+ 1) in the basis YLM (θ, φ). In the expression of the Hamiltonian m is the
reduced mass of the system comprising of the proton and electron. Next we try to apply Heisenberg’s equation to the
radial momentum operator. Noting that the first term of the Hamiltonian is nothing but p̂2r the Heisenberg equation
= − L̂
. (5.9)
The above equation is the operator form of Newton’s second law in spherical polar coordinates. Next we evaluate the
expectation value of both the sides of the above equation using the wave-functions given in Eq. (4.65). We know,
n3 a20(L +
, (5.10)
a30 n
3L(L+ 1
)(L + 1)
. (5.11)
Using the above expectation values in Eq. (5.9) and noting that 〈L̂2〉 = h̄2L(L+1) we see that the time derivative of
the expectation value of the radial momentum operator of the Hydrogen atom vanishes.
The above analysis shows that the form of the Ehrenfest theorem as given in Eq. (5.4) is only valid in Cartesian
coordinates. In the case of the Hydrogen atom if we used Eq. (5.4) we should have never got the correct result.
VI. CONCLUSION
In the present work we have emphasized on the reality of the momentum expectation value and using the reality
of the expectation value as a bench mark we did find out the form of the momentum operator in spherical polar
coordinate system. We found that most of the concepts which define the momentum operator in Cartesian coordinates
do not hold good in spherical polar coordinates and in general in any other coordinate system. The reason being
that whenever we do an integration in curvilinear coordinates the Jacobian of the coordinate transformation matrix
comes inside the picture and the Cartesian results start to falter if we do not change the rules appropriately. The
forms of the momentum along the radial direction and the form of the angular momentum operators are derived in
section III. The status of the angular variables was briefly discussed in the same section. We explicitly calculated the
expectation values of the momentum operator in various important cases and showed that the expectation value of
the momentum operator do really come out to be zero as expected. Although the expectation value of the momentum
operator vanishes in most of the bound states, with a real potential, the expectation value of the position is not
required to vanish. The expectation value of the position operator is directly related with the parity property of the
potential which was briefly discussed in subsection IVC. At the end we calculated the Heisenberg equation of motion
for the radial momentum operator for the Hydrogen atom and showed its formal semblance with Newton’s second
law. It was also shown that if we properly write the Heisenberg equation of motion in spherical polar coordinates
then Ehrenfest’s theorem follows naturally.
In short we conclude by saying:
1. the forms of the various momentum operators, in most of the coordinate systems, in quantum mechanics can be
obtained by imposing the condition of the reality of their eigenvalues. The form of the probability conservation
equation and Ehrenfest theorem must be modified in curvilinear coordinates to yield meaningful results.
2. There are obvious problems in elevating the status of angular variables to dynamical variables in quantum
mechanics.
3. For compact variables, if the variable is periodic the expectation value of the angular momentum conjugate to
it is non-zero. If the compact variable is not periodic then the angular momentum conjugate to it must vanish.
4. The momentum expectation values in cases of bound state motions vanish, whereas the position expectation
values in those cases depends on the symmetry of the potential.
Acknowledgements
The authors thank Professors D. P. Dewangan, S. Rindani, J. Banerji, P. Panigrahi and Ms. Suratna Das for
stimulating discussions and constant encouragements.
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http://mathworld.wolfram.com/LaguerrePolynomial.html, Equation 24.
http://mathworld.wolfram.com/LaguerrePolynomial.html
Introduction
Definition of the momentum operator and the reality of its expectation value
The expectation value of the momentum operator in Cartesian and spherical polar coordinates
Momentum expectation values in various bound states
Particle in one-dimensional stiff potential wells
Infinite square well potential
Finite square well potential
Dirac-delta potential
Particle in one-dimensional slowly varying potentials
Linear harmonic oscillator potential
Pöschl-Teller potential
Morse potential
Position expectation values for various potentials
Momentum expectation value for a three-dimensional slowly varying spherically symmetric potential
The Hydrogen atom
A discussion on Heisenberg's equation of motion and Ehrenfest theorem
Conclusion
References
|
0704.0376 | Environmental noise reduction for holonomic quantum gates | Environmental noise reduction for holonomic quantum gates
Daniele Parodi,1,2 Maura Sassetti,1,3 Paolo Solinas,4 and Nino Zangh̀ı1,2
1 Dipartimento di Fisica,
Università di Genova, Genova, Italy
2 Istituto Nazionale di Fisica Nucleare (Sezione di Genova),
Genova, Italy 3 INFM-CNR Lamia
Via Dodecaneso 33, 16146 Genova, Italy
4 Laboratoire de Physique Théorique de la Matière Condensée,
Université Pierre et Marie Curie,
Place Jussieu, 75252 Paris Cedex 05, France
(Dated: October 27, 2018)
We study the performance of holonomic quantum gates, driven by lasers, under the effect of
a dissipative environment modeled as a thermal bath of oscillators. We show how to enhance the
performance of the gates by suitable choice of the loop in the manifold of the controllable parameters
of the laser. For a simplified, albeit realistic model, we find the surprising result that for a long time
evolution the performance of the gate (properly estimated in terms of average fidelity) increases. On
the basis of this result, we compare holonomic gates with the so-called Stimulated Raman adiabatic
passage (STIRAP) gates.
PACS numbers: 03.67.Lx
I. INTRODUCTION
The major challenge for quantum computation is posed
by the fact that generically quantum states are very del-
icate objects quite difficult to control with the required
accuracy—typically, by means of external driving fields,
e.g., a laser. The interaction with the many degrees of
freedom of the environment causes decoherence; more-
over, errors in processing the information may lead to a
wrong output state.
Among the approaches aiming at overcoming these dif-
ficulties are those for which the quantum gate depends
very weakly on the details of the dynamics, in particu-
lar, the holonomic quantum computation (HQC) [1] and
the so-called Stimulated Raman adiabatic passage (STI-
RAP) [2, 3, 4]. In the latter, the gate operator is obtained
acting on the phase difference of the driving lasers dur-
ing the evolution, while in the former the same goal is
achieved by exploiting the non-commutative analogue of
the Berry phase collected by a quantum state during a
cyclic evolution. Concrete proposals have been put for-
ward, for both Abelian [5, 6] and non-Abelian holonomies
[7, 8, 9, 10, 11, 12]. The main advantage of the HQC
is the robustness against noise deriving from a imperfect
control of the driving fields [13, 14, 15, 16, 17, 18, 19, 20].
In a recent paper [21] we have shown that the dis-
turbance of the environment on holonomic gates can be
suppressed and the performance of the gate optimized
for particular environments (purely superohmic thermal
bath). In the present paper we consider a different sort
of optimization, which is independent of the particular
nature of the environment.
By exploiting the full geometrical structure of HQC,
we show how the performance of a holonomic gate can
be enhanced by a suitable choice of the loop in the man-
ifold of the parameters of the external driving field: by
choosing the optimal loop which minimizes the “error”
(properly estimated in terms of average fidelity loss). Our
result is based on the observation that there are different
loops in the parameter manifold producing the same gate
and, since decoherence and dissipation crucially depend
on the dynamics, it is possible to drive the system over
trajectories which are less perturbed by the noise. For a
simplified, albeit realistic model, we find the surprising
result that the error decreases linearly as the gating time
increases. Thus the disturbance of the environment can
be drastically reduced. On the basis of this result, we
compare holonomic gates with the STIRAP gates.
In Sec. II the model is introduced and the explicit
expression of the error is derived. In Sec. III we find
the optimal loop, calculate the error, make a comparison
with other approaches, and briefly sketch how to treat a
different coupling with the environment.
II. MODEL
The physical model is given by three degenerate (or
quasidegenerate) states, |+〉, |−〉, and |0〉, optically con-
nected to another state |G〉. The system is driven by
lasers with different frequencies and polarizations, acting
selectively on the degenerate states. This model describes
various quantum systems interacting with a laser radia-
tion, ranging from semiconductor quantum dots, such as
excitons [12] and spin-degenerate electron states [3], to
trapped ions [8] or neutral atoms [7].
The (approximate) Hamiltonian modeling the effect of
the laser on the system is (for simplicity, ~ = 1) [8, 12]
H0(t) =
j=+,−,0
ǫ|j〉〈j|+(e−iǫtΩj(t)|j〉〈G|+H.c)
, (1)
where Ωj(t) are the timedependent Rabi frequencies de-
http://arxiv.org/abs/0704.0376v2
pending on controllable parameters, such as the phase
and intensity of the lasers, and ǫ is the energy of the
degenerate electron states. The Rabi frequencies are
modulated within the adiabatic time tad, (which coin-
cides with the gating time), to produce a loop in the pa-
rameter space and thereby realize the periodic condition
H0(tad) = H0(0).
The Hamiltonian (1) has four time dependent eigen-
states: two eigenstates |Ei(t)〉 , i = 1, 2, called bright
states, and two eigenstates |Ei(t)〉, i = 3, 4, called dark
states. The two dark states have degenerate eigen-
value ǫ and the two bright states have timedependent
energies λ±(t) = [ǫ ±
ǫ2 + 4Ω2(t)]/2 with Ω2(t) =
i=±,0 |Ωi(t)|
2 [22].
The evolution of the state is generated by
Ut = Te
dt′H0(t
′), (2)
where T is the time-ordered operator. In the adiabatic
approximation, the evolution of the state takes place in
the degenerate subspace generated by |+〉, |−〉, and |0〉.
This approximation allows to separate the dynamic con-
tribution and the geometric contribution from the evolu-
tion operator. Expanding Ut in the basis of instantaneous
eigenstates of H0(t) (the bright and dark states), in the
adiabatic approximation, we have
Ut ∼=
′)dt′ |Ej(t)〉〈Ej(t)| Ut, (3)
where
Ut = Te
dτV (τ), (4)
here V is the operator with matrix elements Vij(t) =
〈Ei(t)|∂t|Ej(t)〉. The unitary operator Ut plays the role
of timedependent holonomic operator and is the funda-
mental ingredient for realizing complex geometric trans-
formation whereas
′)dt′ |Ej(t)〉〈Ej(t)| is the
dynamic contribution.
Consider Ut for a closed loop, i.e., for t = tad,
U = Utad . (5)
If the initial state |ψ0〉 is a superposition of |+〉 and |−〉,
then U|ψ0〉 is still a superposition of the same vectors (in
general, with different coefficients)[12]. Thus the space
spanned by |+〉 and |−〉 can be regarded as the “logi-
cal space” on which the “logical operator” U acts as a
“quantum gate” operator. Note that for t < tad, Ut|ψ0〉
has, in general, also a component along |0〉. However,
as it is easy to show [22], at any instant t < tad, Ut|ψ0〉
can be expanded in the twodimensional space spanned
by the dark states |E3(t)〉 and |E4(t)〉. It is important
to observe that U depends only on global geometric fea-
tures of the path in the parameter manifold and not on
the details of the dynamical evolution [1, 12].
To construct a complete set of holonomic quantum
gates, it is sufficient to restrict the Rabi frequencies
Ωj(t) in such a way that the norm Ω of the vector
~Ω = [Ω0(t),Ω+(t),Ω−(t)] is time independent and the
vector lies on a real three dimensional sphere [8, 12].
We parametrize the evolution on this sphere as Ω+(t) =
sin θ(t) cos φ(t), Ω−(t) = sin θ(t) sin φ(t) and Ω0(t) =
cos θ(t) with fixed initial (and final) point in θ(0) = 0,
the north pole By straightforward calculation we obtain
the analytical expression for V (t) in eq. (4), V (t) =
iσy cos[θ(t)]φ̇(t), where σy is the usual Pauli matrix
written in the basis of dark states. Thus, the oper-
ator (4) becomes Ut = cos[a(t)] − iσy sin[a(t)], here
a(t) =
dτφ̇(τ) cos θ(τ). Accordingly, the logical op-
erator U (5) is
U = cos a− iσy sin a, (6)
where
a = a(tad) =
∫ tad
dτφ̇(τ) cos θ(τ) (7)
is the solid angle spanned on the sphere during the evolu-
tion. Note that the are many paths on the sphere which
generate the same logical operator U , and span the same
solid angle a.
In a previous work we have studied how interaction
with the environment disturbs the logical operator U [21].
The goal of the present paper is to analyze whether and
how such a disturbance can be minimized for a given
U . To this end, we model the environment as a thermal
bath of harmonic oscillators with linear coupling between
system and environment [23]. The total Hamiltonian is
H = H0(t) +
α + cαxαA), (8)
where A is the system interaction operator called, from
now on, noise operator.
We now consider the time evolution of the reduced
density matrix of the system, determined by the Hamil-
tonian (8). We rely on the standard methods of the “mas-
ter equation approach,” with the environment treated in
the Born approximation and assumed to be at each time
in its own thermal equilibrium state at temperature T .
This allows to include the effect of the environment in
the correlation function (kB = 1)
g(τ) =
cos(ωτ) − i sin(ωτ)
Here the spectral density is
J(ω) =
δ(ω − ωα), (10)
at the low frequencies regimes, is proportional to ωs, with
s ≥ 0, i.e., s = 1 describes a Ohmic environments, typ-
ical of baths of conduction electrons, s = 3 describes
a super-Ohmic environment, typical of baths of phonons
[21, 24]. The asymptotic decay of the real part of g(τ) de-
fines the characteristic memory time of the environment.
Denoting with ρ̃(t) the time evolution of the reduced den-
sity matrix of the system in the interaction picture, e.g.,
ρ̃(t) = U
t ρUt, one has [24]
ρ̃(tad) = ρ(0) +
dτ{g(τ )[ÃÃ
ρ̃(t− τ )− Ã
ρ̃(t− τ )Ã]
+ g(−τ )[ρ̃(t− τ )Ã
Ã− Ãρ̃(t− τ )Ã
]. (11)
Here à and Ã′ stand for Ã(t) and Ã(t−τ), with the tilde
denoting the time evolution in the interaction picture.
In quantum information the quality of a gate is usually
evaluated by the fidelity F , which measures the closeness
between the unperturbed state and the final state,
F = 〈ψ0(0)|U
†ρ(tad)U|ψ0(0)〉, (12)
where |ψ0(0)〉 is the initial state, and ρ(tad) = U ρ̃(tad)U
is the reduced density matrix in the Schrödinger picture
starting from the initial condition ρ(0) = |ψ0(0)〉〈ψ0(0)|.
The average error is defined as the average fidelity loss,
i.e.,
δ =< 1−F >= 1− < 〈ψ0(0)|ρ̃(tad)|ψ0(0)〉 >, (13)
where < · · · > denotes averaging with respect to the
uniform distribution over the initial state |ψ0(0)〉.
The right-handside of Eq. (13) can be computed by
the following steps:
(1) solving Eq. (11) in strictly second order approxima-
tion; this approximation corresponds to replace ρ̃(t − τ)
with ρ(0);
(2) using the adiabatic approximation U(t − τ, t) ≈
exp(iτH0(t));
(3) expanding the scalar product in Eq. (13) with respect
to a complete orthonormal basis {|ϕn(t)〉}, n = 1, 2, 3,
orthogonal to |ψ0(t)〉. In this way, one obtains
∫ tad
dt G(t)|〈ψ0(t)|A|ϕn(t)〉|
, (14)
where
G(t) =
Re[g(τ)] cos(ω0nt) + Im[g(τ)] sin(ω0nt))
Here, ω0n = ω0−ωn are the energy differences associated
to the transition ψ0 ↔ φn, with ω0 = ǫ, ω1 = λ+, ω2 =
λ−, and ω3 = ǫ.
The interaction between system and environment is ex-
pressed by the noise operator A in Eq. (8). We shall now
make the assumption that A = diag{0, 0, 0, 1} in the |G〉,
|±〉, and |0〉 basis. In this case the transition between de-
generate states are forbidden, however the noise breaks
their degeneracy, shifting one of them. In spite of its sim-
ple form, this A is nevertheless a realistic noise operator
for physical semiconductor systems [4].
III. MINIMIZING THE ERROR
The problem can be stated in the following way: given
the noise operator A and the logical operator U , find a
path on the parameter space (the surface of the sphere,
described above) which minimizes the error δ.
The total error δ, given by Eq. (14), can be decom-
posed as
δ = δtr + δpd, (16)
where the transition error, δtr, is the contribution to the
sum of the nondegenerate states (ω0n 6= 0) and the pure
dephasing error δpd is the contribution of the degenerate
states (ω0n = 0). Thus
δpd =
∫ tad
sin(ωt)
sin2 2a(t)
sin4 θ(t) (17)
δtr =
n=+,−
1 + [(λn − ǫ)/Ω]2
∫ tad
sin2 2θ(t)dt,
where
Γ0n = J(|ω0n|)
|ω0n|
− sgn(ω0n)
correspond to the transition rates calculated by standard
Fermi golden rules, supposing, as usual, G(t) ≈ G(∞) for
g(τ) strongly peaked around τ = 0. In the following we
define for simplicity
n=+,−
1 + [(λn − ǫ)/Ω]2
Since we are interested at long time evolution, we start
discussing the transition error which dominates in this
regime [4, 25].
A. Transition rate
As explained in Sec. II, the holonomic paths are closed
curves on the surface of the sphere which start from the
north pole. It turns out that the curve minimizing δtr
can be found among the loops which are composed by a
simple sequence of three paths (see the Appendix): evo-
lution along a meridian (φ = const), evolution along a
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FIG. 1: The error δtr versus θM for two different a values:
a = π/2 (dashed line) and a = π/4 (full line) correspond to
NOT and Hadamard gate, respectively.
parallel (θ = const) and a final evolution along a merid-
ian to come back to the north pole.
The error δtr in (18), depends on a given by Eq. (7),
θM (the maximum angle spanned during the evolution
along the meridian), ∆φ (the angle spanned along the
parallel), and angular velocity v. We allow ∆φ ≥ 2π
which corresponds to cover more than one loop along
the parallel. The velocity along the parallel is v(t) =
φ̇(t) sin θ and that along the meridian is v(t) = θ̇(t). In
the following we assume that v is constant, and it cannot
exceed the maximal value of vmax, fixed by adiabatic
condition vmax ≪ Ω.
The parameters a, θM , and ∆φ are connected by the
relation a = ∆φ(1 − cos θM ). The error δtr is then
δtr = δ
tr + δ
tr, (20)
where
δMtr =
sin 4θM
is the contribution along the meridian and
δPtr = K
sin θM sin
2 2θM
1− cos θM
is the contribution along the parallel.
In Fig. 1 δtr is plotted for a = π/2 and a = π/4 (corre-
sponding to NOT and Hadamard gate, respectively) as a
function of θM . One can see that δtr has a local minimum
for θM = π/2 and a global minimum for θM = 0 where
the error vanishes. This suggests that the best choice is
to take θM as small as possible.
It is interesting to consider the dependence of δtr also
on the evolution time tad. For simplicity, we set the ve-
locity v = vmax. In this case, changing θM (and then
∆φ) corresponds to a change in the evolution time. We
obtain
θM = arccos
, (23)
5 10 15 20 25 30
vmaxtad
FIG. 2: The error δtr versus vmaxtad for two different a val-
ues: a = π/2 (dashed line) and a = π/4 (full line) correspond
to NOT and Hadamard gate, respectively. The dotted-dashed
line shows the value of the error at θ = π/2. The circles show
the critical value of vmaxtad above which the best loop is the
one with the minimal θM .
where
(vmaxtad)
2 + a2
. (24)
Using these relations, δMtr and δ
tr, given by (21) and
(22) become functions of tad, vmax, and a. Note that
m measures the space covered along the parallel, in fact
∆φ = 2πm.
In Fig. 2 we see the behavior of δtr as a function
of vmaxtad. The first minimum for both curves corre-
sponds to θM = π/2, then the curves for long tad de-
crease asymptotically to zero corresponding to the region
in which θM → 0. In this regime we have δtr ∝ 1/tad
which is drastically different from the results obtained
with other methods where δtr ∝ tad, (see Refs [4, 25]
and below Sec. III C). It should be observed that this
surprising results is a merit of holonomic approach which
allows to choose the loop in the parameter space, with-
out changing the logical operation as long as it subtends
the same solid angle. Observe that small θM and long
tad mean large value of m, i.e., multiple loops around the
north pole.
Figure 2 shows that, for a given gate, there is a criti-
cal value kc of vmaxtad which discriminate between the
choice of θM (e.g., k = 6 for the Hadamard gate and
k = 25 for the NOT gate). For vmaxtad < kc the best
choice for the loop is θM = π/2; For vmaxtad > kc the
best choice is the value of θM determined by eq. (23) and
(24).
Note that the region vmaxtad > kc is accessible with
physical realistic parameters [12]. For example, if we
choose the laser intensity Ω = 20 meV and vmax = Ω/50
(for which values the nonadiabatic transitions are forbid-
den), the critical parameter corresponds to the critical
time of 15 ps for the Hadamard gate and 42 ps for the
NOT gate.
B. Pure Dephasing
Until now we have ignored the pure dephasing effect
because we have assumed that it is negligible in com-
parison with the transition error for long evolution time.
Now, we check that the pure dephasing error contribu-
tion can indeed be neglected. We can write the pure
dephasing error using Eq. (17) and splitting to parallel
and meridian part as
δPpd =
∫ tad
Q[a(t)] sin ωt sin4 θM (25)
δMpd =
Q[a(t)]
sin ωt
sin4(vmaxt) + sin
vmaxt
,(26)
where Q[a(t)] = 1 + 1/2 sin2[2a(t)].
To estimate δpd we assume that tad is longer with re-
spect to the characteristic time of the bath. Remember-
ing that J(ω) ∝ ωs, the pure dephasing error behavior
along the parallel part at the temperature T is
δPpd ∝
, T ≪ 1/tad
, T ≫ 1/tad
while the along meridian is
δMpd ∝
. (28)
Then, we can conclude that the pure dephasing can al-
ways be neglected for long time evolution because it de-
creases faster than the transition error.
C. Comparison between HQC and STIRAP
We make a comparison between holonomic quantum
computation (HQC) and the STIRAP procedure which is
an analogous approach to process quantum information.
The STIRAP procedure ([2, 4]) is, in its basic points,
very similar to the holonomic information manipulation.
The level spectrum, the information encoding, the evolu-
tion produced by adiabatic evolving laser are exactly the
same. The fundamental difference is that in STIRAP
the dynamical evolution is fixed (we must pass through
a precise sequence of states) and then the correspond-
ing loop in the parameter space is fixed. In particular,
we go from the north pole to the south pole and back
to the north pole along meridians. Since the loop, as
in our model, is a sequence of meridian-parallel-meridian
path, we can calculate the error and make a direct com-
parison. In this case, the transition error results propor-
tional to δtr ∝ tad and grows linearly in time while for
HQC δtr ∝ 1/tad. Therefore, the HQC is fundamentally
the favorite for long application times with respect to the
STIRAP ones.
Moreover, we can show that the freedom in the choice
of the loop allows us to construct HQC which perform
better than the best STIRAP gates. In Ref. [4] the mini-
mum error (not depending on the evolution time) for STI-
RAP was obtained reaching a compromise between the
necessity to minimize the transition, pure dephasing error
and the constraint of adiabatic evolution. With realistic
physical parameters [21] (J(ω) = kω3e(−ω/ωc)
, Ω = 10
meV, ǫ = 1eV, vmax = Ω/50, k = 10
−2(meV)−2, ωc = 0.5
meV and for low temperature), the total minimum error
in Ref. [4] is δstirap = 10
−3. With the same parame-
ters, we still have the possibility to increase the evolution
time in order to reduce the environmental error. How-
ever, for evolution time tad = 50 ps we obtain a total
error δ = 1.5× 10−4 for the NOT gate and δ = 4× 10−5
for the Hadamard gate, respectively. As can be seen, the
logical gate performance is greatly increased.
D. More general noise
Until now we have discussed the possibility to minimize
the environmental error by choosing a particular loop in
the parameter sphere but the structure of the error func-
tional clearly depends on the system-environment inter-
action. Then one might wonder if the same approach can
be used for a different noise environment.
For this reason, we now briefly analyze the case of noise
matrix in the form A = diag{0, 1, 0,−1}. Again, for long
evolution we can neglect the contribution of the pure de-
phasing and focus on the transition error. In this case the
interesting part of the error functional takes the form
δtr = K[(
sin 2θ cos 2θ)2 + (sin θ sin 2φ)2]. (29)
Even if the analysis in this case is much more com-
plicated, it can be seen that δtr has an absolute mini-
mum for θM = 0. The long time behavior is the same
(δtr ∝ 1/tad) such that the results are qualitatively anal-
ogous to the above ones: for small θM loops (or long
evolution at fixed velocity) the holonomic quantum gate
presents a decreasing error. Then even in this case it is
possible to minimize the environmental error.
IV. CONCLUSIONS
In summary, we have analyzed the performance of
holonomic quantum gates in the presence of environmen-
tal noise by focusing on the possibility to have small
errors choosing different loops in the parameter mani-
fold. Due to the geometric dependence, we can imple-
ment the same logical gate with different loops. Since
different loops correspond to different dynamical evolu-
tions, we have used this freedom to construct an evolu-
tion through “protected” or “weakly influenced” states
leading to good holonomic quantum gates performances.
This allows to select (once that the physical parameter
are fixed) the best loop which minimizes the environ-
mental effect. (Note that this optimization procedure
is rather independent of the details of the simple model
we have considered and arguably, it could be extended
to more complicated systems without any substantial
modification.) We have shown that for long time evo-
lutions the noise decreases as 1/tad while in the other
cases it increases linearly with adiabatic time. We also
have shown that the same features can be found with
different kinds of noise suggesting the possibility to find
a way to minimize the environmental effect in the pres-
ence of any noise. These results open a new possibility
for implementation of holonomic quantum gates to build
quantum computation because they seem robust against
both control error and environmental noise.
Acknowledgment
The autors thank E. De Vito for useful discussions.
One of the authors (P. S.) acknowledges support from
INFN. Financial support by the italian MIUR via
PRIN05 and INFN is acknowledged.
APPENDIX A: MINIMIZING THEOREM
Let us consider the family Cn composed of the closed
curves generated by a sequence of n paths along a parallel
(θ = const) alternated with paths along a meridian (φ =
const). We call Cn a generic curve in this family. For
example, the family C1 contains all the closed curves com-
posed by the sequence of path meridian-parallel-meridian
while the family C2 contains the curves meridian-parallel-
meridian-parallel-meridian.
We argue that the closed curve minimizing the error
in Eq. (18) can be found in the C1 family. First, we
show that any closed curve in C2 spanning a solid angle
a on the sphere can be replaced by a closed curve in C1
spanning the same angle and producing a smaller error.
In analogous way any closed curve in C3 can be replaced
by a closed curve in C2 with smaller error and so on. By
induction we obtain that any closed curve in Cn can be
replaced by a curve in C1 spanning the same solid angle
but producing smaller error. Since the curve belonging
to Cn can approximate any closed curve on the sphere,
the best curve can be found in C1.
The crucial point is to show that any curve in C2 can be
replaced by a curve in C1. Let us consider a generic curve
C2 in C2 spanning a solid angle a: composed by a seg-
ment of a meridian (with θ going from 0 to θ1), a parallel
(spanning a ∆φ1 angle), meridian (with θ : θ1 → θ2), a
parallel (spanning a ∆φ2 angle), and finally a segment
to the north pole along a meridian. Let us consider two
closed curves C11 and C
1 in C1 subtending the same solid
angle a with, respectively, θ1 and θ2 as maximum angle
spanned during the evolution along the meridian. First
we analyze (20) along the meridian. Without losing gen-
erality, we can take θ1 < θ2; it is clear from Eq. (21) that
the value of δtr along the meridian for C
1 is smaller that
for C21 : δ
. We note from the Eq. (21), suitable
extended to C2, that the two paths along the meridians
depends only on θ2 and then produce the same error of
C21 ,
< δMC2
= δMC2 . (A1)
The difference between the contribution along the par-
allel is
δPC2−δ
= ∆φ1
sin θ1 sin
2 2θ1−
1− cos θ1
1− cos θ2
sin θ2 sin
2 2θ2
δPC2−δ
= ∆φ2
sin θ2 sin
2 2θ2−
1− cos θ2
1− cos θ1
sin θ1 sin
2 2θ1
Analysis of the positivity of the quantities given by
Eqs. (A2) and (A3) shows that δPC2 cannot be at the
same time smaller than δP
and δP
. In fact, there are
two possibilities: If δPC2 > δ
, from Eq. (A1) and (A3),
δC2 = δ
+ δPC2 > δ
+ δPC1
= δC1
, (A4)
and the best closed curve is C11 . If δ
, from Eqs.
(A1) and (A2),
δC2 = δ
+ δPC2 > δ
+ δPC2
= δC2
, (A5)
and the best closed curve is C21 .
In the same way it can be shown that any closed curve
in C3 can be replaced by a closed curve in C2 with smaller
error.
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is |E1〉 =
(Ω|e〉 +
Ωi|i〉) and |E2〉 =
(−Ω|e〉 +
Ωi|i〉); for the dark states is
|E3〉 = 1/(Ω
|Ω+|2 + |Ω−|2)[Ω0(Ω+|+〉 + Ω−|−〉) −
(Ω2 − |Ω0|
2)|0〉]) and |E4〉 = 1/
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|
0704.0377 | The lifetime of unstable particles in electromagnetic fields | arXiv:0704.0377v3 [hep-ph] 22 Dec 2008
The lifetime of unstable particles in electromagnetic fields
Daniele Binosi1 and Vladimir Pascalutsa1, 2
1ECT* Trento, Villa Tambosi, Villazzano, I-38050 TN, Italy
2Institut für Kernphysik, Johannes Gutenberg Universität, Mainz D-55099, Germany
(Dated: October 30, 2018)
Abstract
We show that the electromagnetic moments of unstable particles (resonances) have an absorptive
contribution which quantifies the change of the particle’s lifetime in an external electromagnetic
field. To give an example we compute here the imaginary part of the magnetic moment for the
cases of the muon and the neutron at leading order in the electroweak coupling. We also consider
an analogous effect for the strongly-decaying ∆(1232) resonance. The result for the muon is
Imµ = eG2Fm
3/768π3, with e the charge and m the mass of the muon, GF the Fermi constant,
which in an external magnetic field of B Tesla give rise to the relative change in the muon lifetime
of 3 × 10−15 B. For neutron the effect is of a similar magnitude. We speculate on the observable
implications of this effect.
PACS numbers: 13.40.Em, 13.35.-r, 12.15.Lk, 23.40.-s
http://arxiv.org/abs/0704.0377v3
I. INTRODUCTION
The electromagnetic (e.m.) moments of a particle are among the few fundamental quan-
tities which describe the particle properties and as such have thoroughly been studied. The
most renowned examples are the magnetic moments of the electron and the muon which
have been measured to unprecendented accuracy and yielded a number of physical insights,
see[1] for recent reviews. What is far lesser known is that the e.m. moments of unstable
particles are complex numbers in general [2, 3]. Their imaginary part reflects, of course, the
unstable nature of the particle, however, the precise interpretation has been missing. In this
paper we work out the relation, suggested first by Holstein [4], which should exist between
the imaginary part of the magnetic moment and the effect of an external magnetic field on
particle’s lifetime.
The argument for such a relation is very simple. The (self-)energy of the particle with a
lifetime τ has an absorptive part, which has an interpretation of the width Γ = 1/τ . The
particle’s magnetic moment ~µ in the presence of magnetic field ~B induces the change in
the energy: −~µ · ~B. The latter contribution can then also change the width, provided the
magnetic moment has an absorptive part (Imµ 6= 0).
The decay properties of unstable particles, such as muon or neutron are extremely
well studied and are widely used for the precise determination of the Standard Model
parameters[5, 6]. There are also a plethora of studies of how these particles behave in
e.m. fields. A well-known example is the search for the neutron’s electric dipole moment[7].
In view of these studies it is compelling to investigate how the decay properties of unstable
particles may be affected by e.m. fields.
The lifetime of unstable quantum-mechanical systems is known to be affected by an
e.m. field. Positronium provides a textbook example[8], where the effect arises due to the
admixture of para- (S = 0) and ortho- (S = 1) positronium states with orbital momentum
l = 0 by the magnetic field interacting with the magnetic moments of the constituents. As
the result, already in the field of B = 0.2 Tesla, the lifetime of ortho-positronium decreases
by almost a factor of 2.
It is far from obvious how the same kind of an effect can arise for an elementary unstable
particle, e.g., the muon. The above-mentioned relation between the imaginary part of the
magnetic moment and the lifetime change may, therefore, provide us with both an interpre-
µ νµ µ
FIG. 1: The muon self-energy contributing to its decay width.
tation for the imaginary part of the magnetic moment and the means to compute the effect
of the lifetime change.
In the following we examine in detail the case of the muon, compute the leading contri-
bution to Imµ and the corresponding effect on the lifetime. Then we will briefly discuss the
cases of the neutron and of the ∆-resonance.
II. MUON DECAY (µ → e νeνµ)
The leading contribution to the muon decay width arises at two-loop level, see Fig. 1. For
our purposes, the W propagators in this graph can safely be assumed to be static — Fermi
theory. We also neglect the mass of the electron in the loops, since it leads to an under-
percent correction of O(me/m); here and in what follows, m is the muon mass. The graphs
with other Standard Model fermions (e.g., quarks) in the loops need not to be considered
here, because they cannot give any contribution to the muon width.
Using dimensional regularization, we compute this graph in d = 4 − 2ǫ dimensions (in
the limit ǫ → 0+),[14]
Σ (p/) =
(2π)d
2γµ(1− γ5) (p/− k/) γν
(p− k)2 + iε
Πµν(k). (1)
where MW is the W -boson mass, g = |e|/ sin θW is the electroweak coupling related to the
Fermi constant by GF/
2 = g2/8M2W , e is the charge, θW is the Weinberg angle, and
Πµν(k) =
d (d− 2)
(4π)d/2(d− 1)
Γ (ǫ)Γ (1− ǫ)2
Γ (2− 2ǫ)
× (−k2)−ǫ
k2gµν − kµkν
is the one-loop correction to the polarization tensor of the W boson. The decay width can
then be found as Γ = −2 ImΣ (p/ = m). A brief calculation shows that the self-energy has
the following form:
Σ (p/) = v(s) p/ (1− γ5) , (3)
with s = p2 and the scalar function v given by:
v(s) = −
G2F s
3(4π)4
− 2γE − 2 ln
+O(ǫ)
, (4)
where γE = −Γ ′(1) is the Euler’s constant. The absorptive part of this function stems from
the logarithm [ln(−s− iε) = ln s− iπ, for s > 0]:
Im v(s) = −
G2F s
384π3
. (5)
Terefore, the width is Γ = −2m Im v(m2), and the muon lifetime:
τ = 192π3/(G2Fm
5) ≃ 2.187× 10−6 sec, (6)
This result is of course long-known due to the seminal work of Feynman and Gell-Mann on
Fermi theory[9]. It is in a percent agreement with the experimental value[5]:
τ (exp) = (2.19703± 0.00004) 10−6 sec, (7)
The discrepancy is due the neglect of the electron mass and some radiative corrections,
c.f.[10]. We now investigate the influence of the e.m. field on the leading contribution given
by Eq. (6).
Let us denote by Σ (x, y;Aµ) the self-energy in the presence of an external e.m. field Aµ.
It is obtained by minimal substitution (∂µ → ∂µ − ieAµ) of the derivatives of all charged
fields into the self-energy of Fig. 1. Expanding in the e.m. coupling, we obtain:
Σ [x, y;Aµ] = Σ (i∂/
x) δ4(x− y)
dz Λµ(x, y; z)Aµ(z) +O(e
2A2), (8)
where Σ (i∂/) is the already computed self-energy in the vacuum, while Λ is the e.m. vertex
correction of Fig. 2, with static W ’s.
Denoting p (p′) the 4-momentum of the initial (final) muon and assuming the on-shell
situation (p2 = p′
= p · p′ = m2), the vertex correction has in the momentum space the
following general form:
Λµ(p′, p) = e
F γµ +G
(p+ p′)µ
+ FA γ
, (9)
where F , G and FA are complex numbers. Note that eF/2m is the correction to the magnetic
moment, and eF + eG is the correction to the electric charge. The Ward-Takahashi (WT)
identity:
(p′ − p) · Λ(p′, p) = e [Σ (p/)− Σ (p/′)] (10)
µ µνµ
FIG. 2: Electromagnetic correction to the muon decay.
with the self-energy in Eq. (3) leads to the following conditions:
F +G = −v(m2)− 2m2v′(m2), FA = v(m2) . (11)
Therefore, the term FA is in fact necessary by the e.m. gauge invariance. The γ5 terms, in
both self-energy and the vertex, are shown to vanish when summing over all the fermions
in Standard Model[11]. However, this does not happen for the imaginary part because the
heavier fermions do not contribute.
The expression for the graph in Fig. 2 is (in Fermi theory) given by
Λµ(p′, p) = −
64M4W
(2π)d
(2π)d
γβ(1− γ5)
2γα(1− γ5) (p/′ − k/1) γµ (p/− k/1) γβ (k/1 − k/2)
(k1 − k2)2 (p− k1)2 (p′ − k1)2
After a lengthy calculation we obtain the following result:
ImF =
384π3
, ImG =
, ImFA = −
384π3
, (13)
hence satisfying the gauge-invariance conditions Eq. (11), for Im v given by Eq. (5).
We would like to emphasize here that, of course, not only the magnetic moment, but
also the charge operator receives an imaginary contribution, equal to e Im(F +G). However,
through the WT identity, this contribution is completely fixed by the momentum dependence
of the self-energy, and therefore is not independent. The same holds for FA. We thus
discuss only the effect of the absorptive part of the magnetic moment, here given by Imµ =
e ImF/2m = eG2Fm
3/768π3.
The energy of the magnetic moment interacting with the magnetic field is equal to −µBz,
with Bz being the projection of the field along the muon spin. Then the total energy, in the
muon rest-frame, is given by: m − (i/2)Γ− µBz. We thus deduce that the absorptive part
yields the following change in the muon width:
∆Γ = 2 ImµBz =
192π3
Bz , (14)
while the change in the lifetime is ∆τ = −(∆Γ/Γ) τ , for ∆Γ/Γ ≪ 1.
Given this result, we conclude that the positively-charged muons live shorter (longer) in
a uniform magnetic field if their spin is aligned along (against) the field. For the relative
change in the width we find:
|eBz|
<∼ 3× 10−15B T−1, (15)
where B is the strength of the field in Tesla. Therefore, in moderate magnetic fields the
change in the muon lifetime is tiny, well beyond the present experimental accuracy (which
is at the ppm level). We will dwell on this more in the concluding part of the paper, but for
now we turn to a more technical issue.
It is interesting to observe that the result of Eq. (13), can simply be obtained by the
minimal substitution into Eq. (3), rather than into the electron propagator in Eq. (1). To
show this we go to coordinate space and hence write the self-energy as Σ (x, y) = Σ (i∂/
) δ(x − y). The minimal substitution to the first order in e leads to the following vertex
correction:
Λ̃µ(x, y; z) = − δ/δAµ(z)Σ (i∂/ + eA/ ) δ(x− y) |A=0 . (16)
Note that in general this is different from the vertex function in Eq. (8), since in the latter
the minimal substitution is performed also in the internal lines. The general form of Eq. (9),
of course, applies here as well, but now the scalar functions are completely specified by the
self-energy:
F̃ = −v(m2), G̃ = −2m2 v′(m2), F̃A = v(m2) . (17)
Substituting the explicit form of Im v, we see that this method unambiguously leads to
exactly the same result [Eq. (13)] as the full calculation. We emphasize though, that this
method cannot always work (see, e.g., Ref.[12]), as will also be clear from the following
examples. Nevertheless, it is worthwhile to investigate this method further, since knowing
whether it is applicable a priori can enormously facilitate the calculations.
III. NEUTRON DECAY AND THE ∆-RESONANCE
We consider now the neutron β-decay. Assuming exact V − A interaction (gA = 1) and
neglecting the electron mass (but not the proton mass, mp), the corresponding two-loop
self-energy can still be written in the form of Eq. (3). We introduce δ = (s −m2p)/2s and
treat it as a small parameter, since in the physical case (where s = m2n), δ ≃ 1.293 × 10−3.
A simple calculation then yields:
Im v(s) = −G
F |Vud|2
s2 δ5, (18)
where Vud is the relevant quark-mixing (CKM) matrix element. We note in passing that
this result leads to the lifetime of τn ≈ 622 sec, to be compared with the experimental
value of 886 sec. This 30% disagreement is largely due to the fact that in reality the axial
coupling gA deviates from 1. However, for our order-of-magnitude estimate this discrepancy
is unimportant.
What is important is that the derivative of the self-energy is enhanced by one power of
Im v′(s) = −(GF |Vud|)
s δ4 . (19)
and this opens the possibility for the enhancement of the effect in the lifetime. Namely, the
relative change in the neutron width then goes as
|∆Γn|
∼ µN |Bz|
mn −mp
<∼ 3× 10−14B T−1, (20)
where µN ≃ 3.15× 10−14 MeV T−1 is the nuclear magneton. A more precise analysis of this
effect for the neutron is beyond the scope of this paper. We focus instead on the example
of the ∆-resonance, where such an enhancement will be shown to be even more dramatic,
at least qualitatively.
The ∆ resonance decays strongly into the pion and nucleon, ∆ → πN , and the cor-
responding self-energy, to leading order in chiral effective-field theory, yields the following
result for the absorptive part[3]:
ImΣ∆(p/) = −23πλ
3C2 (α p/+mN ) , (21)
where the isospin symmetry is assumed, e.g., mp = mn = mN . The constant C =
hAm∆/8πfπ ≃ 1.5, where hA represents the πN∆ coupling and is fitted to the empira-
cal width of the ∆, fπ ≃ 93 MeV is the pion-decay constant, and m∆ = 1232 MeV is the ∆
π(b)(a)
FIG. 3: The leading chiral-loop correction to the magnetic moment of the ∆.
mass. For simplicity we neglect the pion mass (i.e., take the chiral limit). Then, in Eq. (21),
λ = (s−m2N)/2s, α = 1− λ. For s = m∆, λ ≈ (m∆ −mN)/mN ∼ 1/3 is a small parameter
in the chiral effective-field theory with ∆’s (see Ref.[13] for a recent review), and will so be
treated here too.
The absorptive part of the magnetic dipole moment of the ∆ arises at this order from
graphs in Fig. 3. These graphs, computed in Ref.[3], in the chiral limit yield the following
result (upto λ4 terms):
ImF (a) = 4πC2(λ− 3λ2 + 43
λ3) ,
ImG(a) = 4πC2(−λ+ 4λ2 − 71
λ3) ,
ImF (b) = 4πC2(λ2 + 1
λ3) , (22)
ImG(b) = −32
πC2λ3 ,
where F and G correspond with the decomposition in Eq. (9), with the superscript referring
to the corresponding graphs in Fig. 3; FA is absent in this case, of course.
First of all, we observe that this result satisfies the WT conditions, Eq. (11), for each of
the four charge states of the ∆,
∆++ : Im [F (a) +G(a) + F (b) +G(b)] = −2 ImΣ ′∆ ,
∆+ : Im [1
(F (a) +G(a)) + 2
(F (b) +G(b))] = − ImΣ ′∆ ,
∆0 : Im [−1
(F (a) +G(a)) + 1
(F (b) +G(b))] = 0 , (23)
∆− : − Im [F (a) +G(a)] = ImΣ ′∆ ,
where Σ ′∆ = ∂/∂p/Σ∆(p/)|p/=m∆, and hence ImΣ ′∆ = 4πC2(−λ2 + 73λ
At the same time, the ‘naive’ minimal-substitution procedure [Eq. (16)], that happens
to work for the muon, fails here miserably. It would predict that the magnetic moment
contribution would go with the same power as the self-energy [Eq. (17)], which for the
absorptive part means Imµ ∼ ImΣ (m∆) ∼ λ3. In reality it goes as λ. E.g., for the ∆+:
Imµ∆+ = (e/2m∆) Im[
F (a) + 2
F (b)]
π µNC
2 λ+O(λ2). (24)
The fact that the self-energy goes as λ3, while Imµ as λ has as a consequence the enhance-
ment of the lifetime change in the magnetic field by two powers of λ.
Quantitatively such enhancements of the lifetime change over the lifetime by the phase-
space volume do not make much difference in the above examples. However, it shows that
it might be useful to look for manifestations of the lifetime change in the medium where the
phase-space volume can be varied.
IV. CONCLUSIONS AND OUTLOOK
We have examined her a concept of the ‘absorptive magnetic moment’ — an intrinsic
property of an unstable particle, together with the width or the lifetime. It manifests itself
in the change of the particle’s lefetime in an external magnetic field, see Eq. (25) below.
We have computed this quantity for the examples of muon, neutron and ∆-resonance to
leading order in couplings. In all the three cosidered cases the effect on the lifetime is tiny
for normal magnetic fields: in a uniform field of 1 Tesla the change in the lifetime is of order
of 10−13 percent, at most.
In the case of the muon we have computed this effect to the leading order in the elec-
troweak coupling; the change in the lifetime is
∆τ = −2 ImµBz τ 2 = −96π3eBz/(G2Fm7) , (25)
or, numerically, |∆τ | <∼ 6× 10−21 (B/T) sec. A direct measurement of this effect is therefore
beyond the present experimental precision. Nevertheless, it is worthwhile to investigate
the effect of the magnetic field on the differential decay rates, with the hope that some
asymmetries could show a significantly bigger sensitivity.
A notable feature of this effect is that the relative change of the lifetime is inversely
proportional to the phase space. It goes as (mn −mp)−1 in the neutron case, and as (m∆ −
−2 in the ∆-resonance case. (The difference in power is apparently because the neutron
decays solely into fermions while the ∆ has a boson in the decay product.) One can expect
that in the conditions where the phase-space is significantly reduced, e.g. for the neutron in
nuclear medium, the effect of the lifetime change may become measurable.
Especially interesting would be to evaluate the manifestations of this effect in neutron star
formations. Not only the phase-space of the neutron decay is shrinking, the protons decay
too, and all that occurs in magnetic fields as large as 1010 Tesla. Even larger fields can be
achieved in atomic or nuclear systems. Finally, it is worthwhile to point out that in lattice
QCD studies strong magnetic fields are standardly used to compute the electromagnetic
properties of hadrons. Combined with the lattice techniques of extracting the width, the
relation between the absorprive part and the lifetime change may allow to compute the
former on the lattice for unstable hadrons.
Acknowledgments
We thank Barry Holstein and Marc Vanderhaeghen for a number of insightful discus-
sions. The work of V.P. is partially supported by the European Community-Research In-
frastructure Activity under the FP6 ”Structuring the European Research Area” programme
(HadronPhysics, contract RII3-CT-2004-506078).
[1] M. Passera, J. Phys. G 31, R75 (2005); J. P. Miller, E. de Rafael and B. L. Roberts, Rept.
Prog. Phys. 70, 795 (2007).
[2] L. V. Avdeev and M. Y. Kalmykov, Phys. Lett. B 436, 132 (1998).
[3] V. Pascalutsa and M. Vanderhaeghen, Phys. Rev. Lett. 94, 102003 (2005); Phys. Rev. D 77,
014027 (2008).
[4] B. R. Holstein, unpublished.
[5] W. M. Yao et al. [Particle Data Group],“Review of particle physics,”J. Phys. G 33, 1 (2006).
[6] D. Tomono [RIKEN RAL R77 Collaboration], AIP Conf. Proc. 842, 906 (2006); K. R. Lynch,
AIP Conf. Proc. 870, 333 (2006); J. S. Nico, AIP Conf. Proc. 870, 132 (2006); A. P. Serebrov
et al., arXiv:nucl-ex/0702009.
[7] P. G. Harris et al., Phys. Rev. Lett. 82, 904 (1999); C. A. Baker et al., Phys. Rev. Lett. 97,
131801 (2006).
[8] J. L. Basdevant and J. Dalibard, “Quantum Mechanics Solver,” (Springer, Berlin, 2005).
[9] R. P. Feynman and M. Gell-Mann, Phys. Rev. 109, 193 (1958).
[10] T. van Ritbergen and R. G. Stuart, Phys. Rev. Lett. 82, 488 (1999).
[11] A. Czarnecki and B. Krause, Nucl. Phys. Proc. Suppl. 51C, 148 (1996).
[12] J. H. Koch, V. Pascalutsa and S. Scherer, Phys. Rev. C 65, 045202 (2002).
[13] V. Pascalutsa, M. Vanderhaeghen and S. N. Yang, Phys. Rept. 437, 125 (2007).
[14] Our conventions are: metric (+,−,−,−), ε0123 = +1, γ5 = iγ0γ1γ2γ3, γ’s stand for Dirac
matrices and their totally-antisymmetric products: γµν = 1
[γµ, γν ], γµνα = 1
{γµν , γα},
γµναβ = 1
[γµνα, γβ].
|
0704.0378 | An equilibrium problem for the limiting eigenvalue distribution of
banded Toeplitz matrices | AN EQUILIBRIUM PROBLEM FOR THE LIMITING
EIGENVALUE DISTRIBUTION OF BANDED TOEPLITZ
MATRICES
MAURICE DUITS AND ARNO B.J. KUIJLAARS
Abstract. We study the limiting eigenvalue distribution of n×n banded
Toeplitz matrices as n → ∞. From classical results of Schmidt-Spitzer
and Hirschman it is known that the eigenvalues accumulate on a spe-
cial curve in the complex plane and the normalized eigenvalue counting
measure converges weakly to a measure on this curve as n → ∞. In this
paper, we characterize the limiting measure in terms of an equilibrium
problem. The limiting measure is one component of the unique vector
of measures that minimes an energy functional defined on admissible
vectors of measures. In addition, we show that each of the other com-
ponents is the limiting measure of the normalized counting measure on
certain generalized eigenvalues.
1. Introduction
For an integrable function a : {z ∈ C | |z| = 1} → C defined on the unit
circle in the complex plane, the n× n Toeplitz matrix Tn(a) with symbol a
is defined by
Tn(a)
= aj−k, j, k = 1, . . . , n, (1.1)
where ak is the kth Fourier coefficient of a,
a(eiθ)e−ikθ dθ. (1.2)
In this paper we study banded Toeplitz matrices for which the symbol has
only a finite number of non-zero Fourier coefficients. We assume that there
exist p, q ≥ 1 such that
a(z) =
k, ap 6= 0, a−q 6= 0. (1.3)
Department of Mathematics, Katholieke Universiteit Leuven, Celestijnenlaan 200B,
3001 Leuven, Belgium. ([email protected], [email protected]).
The first author is a research assistant of the Fund for Scientific Research – Flanders.
The authors were supported by the European Science Foundation Program MISGAM.
The second author is supported by FWO-Flanders project G.0455.04, by K.U. Leuven
research grant OT/04/21, by Belgian Interuniversity Attraction Pole NOSY P06/02, and
by a grant from the Ministry of Education and Science of Spain, project code MTM2005-
08648-C02-01.
http://arxiv.org/abs/0704.0378v1
2 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Thus Tn(a) has at most p + q + 1 non-zero diagonals. As in [1, p. 263], we
also assume without loss of generality that
g.c.d. {k ∈ Z | ak 6= 0} = 1. (1.4)
We are interested in the limiting behavior of the spectrum of Tn(a) as
n → ∞. We use spTn(a) to denote the spectrum of Tn(a):
spTn(a) = {λ ∈ C | det(Tn(a)− λI) = 0}
Spectral properties of banded Toeplitz matrices are the topic of the recent
book [1] by Böttcher and Grudsky. We will refer to this book frequently,
in particular to Chapter 11 where the limiting behavior of the spectrum is
discussed.
The limiting behavior of spTn(a) was characterized by Schmidt and Spitzer
[10]. They considered the set
lim inf
spTn(a), (1.5)
consisting of all λ ∈ C such that there exists a sequence {λn}n∈N, with
λn ∈ spTn(a), converging to λ, and the set
lim sup
spTn(a), (1.6)
consisting of all λ such that there exists a sequence {λn}n∈N, with λn ∈
spTn(a), that has a subsequence converging to λ. Schmidt and Spitzer
showed that these two sets are equal and can be characterized in terms of
the algebraic equation
a(z) − λ =
k − λ = 0. (1.7)
For every λ ∈ C there are p+q solutions for (1.7), which we denote by zj(λ),
for j = 1, . . . , p+ q. We order these solutions by absolute value, so that
0 < |z1(λ)| ≤ |z2(λ)| ≤ · · · ≤ |zp+q(λ)|. (1.8)
When all inequalities in (1.8) are strict then the values zk(λ) are unambigu-
ously defined. If equalities occur then we choose an arbitrary numbering
so that (1.8) holds. The result by Schmidt and Spitzer [10], [1, Theorem
11.17], is that
lim inf
spTn(a) = lim sup
spTn(a) = Γ0 (1.9)
where
Γ0 := {λ ∈ C | |zq(λ)| = |zq+1(λ)|}. (1.10)
This result gives a description of the asymptotic location of the eigenvalues.
The eigenvalues accumulate on the set Γ0, which is known to be a disjoint
union of a finite number of (open) analytic arcs and a finite number of ex-
ceptional points [1, Theorem 11.9]. It is also known that Γ0 is connected
EIGENVALUES OF BANDED TOEPLITZ MATRICES 3
[13], [1, Theorem 11.19], and that C \ Γ0 need not be connected [1, Theo-
rem 11.20], [2, Proposition 5.2]. See [1] for many beautiful illustrations of
eigenvalues of banded Toeplitz matrices.
The limiting eigenvalue distribution was determined by Hirschman [5], [1,
Theorem 11.16]. He showed that there exists a Borel probability measure
µ0 on Γ0 such that the normalized eigenvalue counting measure of Tn(a)
converges weakly to µ0, as n → ∞. That is,
λ∈spTn(a)
δλ → µ0, (1.11)
where in the sum each eigenvalue is counted according to its multiplicity.
The measure µ0 is absolutely continuous with respect to the arclength mea-
sure on Γ0 and has an analytic density on each open analytic arc in Γ0,
which can be explicitly represented in terms of the solutions of the algebraic
equation (1.7) as follows. Equip every open analytic arc in Γ0 with an orien-
tation. The orientation induces ±-sides on each arc, where the +-side is on
the left when traversing the arc according to its orientation, and the −-side
is on the left. The limiting measure µ0 is then given by
dµ0(λ) =
zj+(λ)
zj−(λ)
dλ. (1.12)
where dλ is the complex line element on Γ0 (taken according to the orien-
tation), and where zj±(λ), λ ∈ Γ0, is the limiting value of zj(λ
′) as λ′ → λ
from the ± side of the arc. These limiting values exist for every λ ∈ Γ0,
with the possible exception of the finite number of exceptional points.
Note that the right-hand side of (1.12) is a priori a complex measure and
it is not immediately clear that it is in fact a probability measure. In the
original paper [5] and in the book [1, Theorem 11.16], the authors give a
different expression for the limiting density, from which it is clear that the
measure is non-negative. We prefer to work with the complex expression
(1.12), since it allows for a direct generalization which we will need in this
paper.
Note also that if we reverse the orientation on an arc in Γ0, then the ±-
sides are reversed. Since the complex line element dλ changes sign as well,
the expression (1.12) does not depend on the choice of orientation.
The following is a very simple example, which however serves as a moti-
vation for the results in the paper.
Example 1.1. Consider the symbol a(z) = z+1/z. In this case we find that
Γ0 = [−2, 2] and µ0 is absolutely continuous with respect to the Lebesgue
measure and has density
dµ0(λ)
4− λ2
, λ ∈ (−2, 2). (1.13)
4 MAURICE DUITS AND ARNO B.J. KUIJLAARS
This measure is well-known in potential theory and is called the arcsine
measure or the equilibrium measure of Γ0, see e.g. [9]. It has the property
that it minimizes the energy functional I defined by
I(µ) =
|x− y| dµ(x) dµ(y), (1.14)
among all Borel probability measures µ on [−2, 2]. The measure µ0 is also
characterized by the equilibrium condition
log |x− λ| dµ0(λ) = 0, x ∈ [−2, 2], (1.15)
which is the Euler-Lagrange variational condition for the minimization prob-
The fact that µ0 is the equilibrium measure of Γ0 is special for symbols
a with p = q = 1. In that case one may think of the eigenvalues of Tn(a) as
charged particles on Γ0, each eigenvalue having a total charge 1/n, that repel
each other with logarithmic interaction. The particles seek to minimize the
energy functional (1.14). As n → ∞, they distribute themselves according
to µ0 and µ0 is the minimizer of (1.14) among all probability measures
supported on Γ0.
The aim of this paper is to characterize µ0 for general symbols a of the
form (1.3) also in terms of an equilibrium problem from potential theory.
The corresponding equilibrium problem is more complicated since it involves
not only the measure µ0, but a sequence of p+ q − 1 measures
µ−q+1, µ−q+2, . . . , µ−1, µ0, µ1, . . . , µp−2, µp−1
that jointly minimize an energy functional.
2. Statement of results
2.1. The energy functional. To state our results we need to introduce
some notions from potential theory. Main references for potential theory in
the complex plane are [8] and [9].
We will mainly work with finite positive measures on C, but we will also
use ν1 − ν2 where ν1 and ν2 are positive measures. The measures need not
have bounded support. If ν has unbounded support then we assume that
log(1 + |x|) dν(x) < ∞. (2.1)
In that case the logarithmic energy of ν is defined as
I(ν) =
|x− y| dν(x)dν(y) (2.2)
and I(ν) ∈ (−∞,+∞].
EIGENVALUES OF BANDED TOEPLITZ MATRICES 5
Definition 2.1. We define Me as the collection of positive measures ν on
C satisfying (2.1) and having finite energy, i.e., I(ν) < +∞. For c > 0 we
define
Me(c) = {ν ∈ Me | ν(C) = c}. (2.3)
The mutual energy I(ν1, ν2) of two measures ν1 and ν2 is
I(ν1, ν2) =
|x− y| dν1(x)dν2(y). (2.4)
It is well-defined and finite if ν1, ν2 ∈ Me and in that case we have
I(ν1 − ν2) = I(ν1) + I(ν2)− 2I(ν1, ν2). (2.5)
If ν1, ν2 ∈ Me(c) for some c > 0, then
I(ν1 − ν2) ≥ 0, (2.6)
with equality if and only if ν1 = ν2. This is a well-known result if ν1 and ν2
have compact support [9]. For measures in Me(c) with unbounded support,
this is a recent result of Simeonov [11], who obtained this from a very elegant
integral representation for I(ν1 − ν2). It is a consequence of (2.6) that I is
strictly convex on Me(c), since
ν1 + ν2
(I(ν1) + I(ν2))− I
ν1 − ν2
(I(ν1) + I(ν2)) , for ν1, ν2 ∈ Me(c),
with equality if and only if ν1 = ν2.
Before we can state the equilibrium problem we also need to introduce
the sets
Γk := {λ ∈ C | |zq+k(λ)| = |zq+k+1(λ)|}, k = −q + 1, . . . , p− 1, (2.7)
which for k = 0 reduces to the definition (1.10) of Γ0. We will show that
each Γk is the disjoint union of a finite number of open analytic arcs and
a finite number of exceptional points. All Γk are unbounded, except for Γ0
which is compact.
The equilibrium problem will be defined for a vector of measures denoted
by ~ν = (ν−q+1, . . . , νp−1). The component νk is a measure on Γk satisfying
some additional properties that are given in the following definition.
Definition 2.2. We call a vector of measures ~ν = (ν−q+1, . . . , νp−1) admis-
sible if νk ∈ Me, νk is supported on Γk, and
νk(Γk) =
if k ≤ 0,
if k ≥ 0,
(2.8)
for every k = −q + 1, . . . , p − 1.
Now we are ready to state our first result. The proof is given in section
6 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Theorem 2.3. Let the symbol a satisfy (1.3) and (1.4), and let the curves
Γk be defined as in (2.7). For each k ∈ {−q + 1, . . . , p − 1}, define the
measure µk on Γk by
dµk(λ) =
zj+(λ)
zj−(λ)
dλ, (2.9)
where dλ is the complex line element on each analytic arc of Γk according
to a chosen orientation of Γk (cf. discussion after (1.12)). Then
(a) ~µ = (µ−q+1, . . . , µp−1) is admissible.
(b) There exist constants lk such that
log |λ− x| dµk(x) =
log |λ− x| dµk+1(x) +
log |λ− x| dµk−1(x) + lk,
(2.10)
for k = −q + 1, . . . , p − 1, and λ ∈ Γk. Here we let µ−q and µp be
the zero measures.
(c) ~µ = (µ−q+1, . . . , µp−1) is the unique minimizer of the energy func-
tional J defined by
J(~ν) =
k=−q+1
I(νk)−
k=−q+1
I(νk, νk+1) (2.11)
for admissible vectors of measures ~ν = (ν−q+1, . . . , νp−1).
The relations (2.10) are the Euler-Lagrange variational conditions for the
minimization problem for J among admissible vectors of measures.
It may not be obvious that the energy functional (2.11) is bounded from
below. This can be seen from the alternative representation
J(~ν) =
I(ν0) +
k(k + 1) I
ν−q+k
− ν−q+k+1
k + 1
k(k + 1) I
− νp−k−1
k + 1
. (2.12)
We leave the calculation leading to this identity to the reader. Under the
normalizations (2.8) it follows by (2.6) that each term in the two finite sums
on the right-hand side of (2.12) is non-negative, so that
J(~ν) ≥
I(ν0).
Since ν0 is a Borel probability measure on Γ0 and Γ0 is compact, we indeed
have that the energy functional is bounded from below on admissible vectors
of measures ~ν.
The alternative representation (2.12) will play a role in the proof of The-
orem 2.3.
EIGENVALUES OF BANDED TOEPLITZ MATRICES 7
Yet another representation for J is
J(~ν) =
j,k=−q+1
Ajk I(νj , νk) (2.13)
where the interaction matrix A has entries
Ajk =
1, if j = k,
, if |j − k| = 1,
0, if |j − k| ≥ 2.
(2.14)
The energy functional in the form (2.13) and (2.14) also appears in the
theory of simultaneous rational approximation, where it is the interaction
matrix for a Nikishin system [7, Chapter 5].
It allows for the following physical interpretation: on each of the curves Γk
one puts charged particles with total charge (q+k)/q or (p−k)/p, depending
on whether k ≤ 0 or k ≥ 0. Particles that lie on the same curve repel each
other. The particles on two consecutive curves interact in the sense that
they attract each other but in a way that is half as strong as the repulsion
on a single curve. Particles on different curves that are not consecutive do
not interact with each other in a direct way.
2.2. The measures µk as limiting measures of generalized eigenval-
ues. By (1.12) and Theorem 2.3 we know that the measure µ0 that appears
in the minimizer of the energy functional J is the limiting measure for the
eigenvalues of Tn(a). It is natural to ask about the other measures µk that
appear in the minimizer. In our second result we show that the measures
µk can be obtained as limiting counting measures for certain generalized
eigenvalues.
Let k ∈ {−q+1, . . . , p− 1}. We use Tn(z−k(a−λ) to denote the Toeplitz
matrix with the symbol z 7→ z−k(a(z) − λ). For example, for k = 1, q = 1
and p = 2, we have
−k(a−λ)) =
a1 a0 − λ a−1
a2 a1 a0 − λ a−1
a2 a1 a0 − λ a−1
. . .
. . .
. . .
. . .
a2 a1 a0 − λ a−1
a2 a1 a0 − λ
a2 a1
Definition 2.4. For k ∈ {−q + 1, . . . , p − 1} and n ≥ 1, we define the
polynomial Pk,n by
Pk,n(λ) = detTn(z
−k(a− λ)) (2.15)
and we define the kth generalized spectrum of Tn(a) by
spk Tn(a) = {λ ∈ C | Pk,n(λ) = 0}. (2.16)
8 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Finally, we define µk,n as the normalized zero counting measure of spk Tn(a)
µk,n =
λ∈spk Tn(a)
δλ (2.17)
where in the sum each λ is counted according to its multiplicity as a zero of
Pk,n.
Note that λ ∈ spk Tn(a) is a generalized eigenvalue (in the usual sense) for
the matrix pencil (Tn(z
−ka), Tn(z
−k)), that is, det(A − λB) = 0 with A =
−ka) and B = Tn(z
−k). If k = 0, then B = I and sp0 Tn(a) = spTn(a).
If k 6= 0, then B is not invertible and the generalized eigenvalue problem is
singular, causing that there are less than n generalized eigenvalues. In fact,
since Tn(z
−k(a−λ)) has exactly n−|k| entries a0−λ, we easily get that the
degree of Pk,n is at most n− |k| and so there are at most n− |k| generalized
eigenvalues. Due to the band structure of Tn(z
−k(a−λ)) the actual number
of generalized eigenvalues is substantially smaller.
Proposition 2.5. Let k ∈ {−q+1, . . . , p−1}. Let Pk,n(λ) = γk,nλdk,n + · · ·
have degree dk,n and leading coefficient γk,n 6= 0. Then
dk,n ≤
n, if k < 0,
n, if k > 0.
(2.18)
Equality holds in (2.18) if either k > 0 and n is a multiple of p, or k < 0
and n is a multiple of q, and in those cases we have
γk,n =
(−1)(k+1)na|k|n/q−q , if k < 0 and n ≡ 0 mod q,
(−1)(k+1)nakn/pp , if k > 0 and n ≡ 0 mod p.
(2.19)
We now come to our second main result. It is the analogue of the results
of Schmidt-Spitzer and Hirschman for the generalized eigenvalues.
Theorem 2.6. Let k ∈ {−q + 1, . . . , p− 1}. Then
lim inf
spk Tn(a) = lim sup
spk Tn(a) = Γk, (2.20)
φ(z) dµk,n(z) =
φ(z) dµk(z) (2.21)
holds for every bounded continuous function φ on C.
The key element in the proof of Theorem 2.6 is a beautiful formula of
Widom [14], see [1, Theorem 2.8], for the determinant of a banded Toeplitz
matrix. In the present situation Widom’s formula yields the following. Let
λ ∈ C be such that the solutions zj(λ) of the algebraic equation (1.7) are
mutually distinct. Then
Pk,n(λ) = detTn(z
−k(a− λ)) =
CM (λ) (wM (λ))
, (2.22)
EIGENVALUES OF BANDED TOEPLITZ MATRICES 9
where the sum is over all subsets M ⊂ {1, 2, . . . , p+ q} of cardinality |M | =
p− k and for each such M , we have
wM (λ) := (−1)p−kap
zj(λ), (2.23)
and (with M := {1, 2, . . . , p+ q} \M),
CM (λ) :=
zj(λ)
(zj(λ)− zl(λ))−1. (2.24)
The formula (2.22) shows that for large n, the main contribution comes from
those M for which |wM (λ)| is the largest possible. For λ ∈ C \ Γk there is a
unique such M , namely
M = Mk := {q + k + 1, q + k + 2, . . . , p + q} (2.25)
because of the ordering (1.8).
2.3. Overview of the rest of the paper. In section 3 we will state some
preliminary results about analyticity properties of the solutions zj of the al-
gebraic equation (1.7). These results will be needed in the proof of Theorem
2.3 which is given in section 4. In section 5 we will prove Proposition 2.5
and Theorem 2.6. Finally, we conclude the paper by giving some examples
in section 6.
3. Preliminaries
In this section we collect a number of properties of the curves Γk and the
solutions z1(λ), . . . , zp+q(λ) of the algebraic equation (1.7). For convenience
we define throughout the rest of the paper
Γ−q = Γp = ∅, and µ−q = µp = 0. (the zero-measure).
Occasionally we also use
z0(λ) = 0, zp+q+1(λ) = +∞.
3.1. The structure of the curves Γk. We start with a definition, cf. [1,
§11.2].
Definition 3.1. A point λ0 ∈ C is called a branch point if a(z) − λ0 = 0
has a multiple root. A point λ0 ∈ Γk is an exceptional point of Γk if λ0 is a
branch point, or if there is no open neighborhood U of λ such that Γk ∩U is
an analytic arc starting and terminating on ∂U .
If λ0 is a branch point, then there is a z0 such that a(z0) = λ0 and
a′(z0) = 0. Then we may assume that z0 = zq+k(λ0) = zq+k+1(λ0) for some
k and λ0 ∈ Γk. For a symbol a of the form (1.3), the derivative a′ has exactly
p+q zeros (counted with multiplicity), so that there are exactly p+q branch
points counted with multiplicity.
10 MAURICE DUITS AND ARNO B.J. KUIJLAARS
The solutions zk(λ) also have branching at infinity (unless p = 1 or q = 1).
There are p solutions of (1.7) that tend to infinity as λ → ∞, and q solutions
that tend to 0. Indeed, we have
zk(λ) =
−1/q(1 +O(λ−1/q)), for k = 1, . . . , q,
1/p(1 +O(λ−1/p)), for k = q + 1, . . . , p + q,
(3.1)
as λ → ∞. Here c1, . . . , cq are the q distinct solutions of cq = a−q (taken in
some order depending on λ), and cq+1, . . . , cp+q are the p distinct solutions
of cp = a−1p (again taken in some order depending on λ).
The following proposition gives the structure of Γk at infinity.
Proposition 3.2. Let k ∈ {−q+1, . . . , p−1}\{0}. Then there is an R > 0
such that Γk ∩ {λ ∈ C | |λ| > R} is a finite disjoint union of analytic arcs,
each extending from |λ| = R to infinity.
Proof. The proof is similar to the proof of [1, Proposition 11.8] where a
similar structure theorem was proved for finite branch points. We omit the
details. �
It follows from Proposition 3.2 that the exceptional points for Γk are in a
bounded set. Since the set of exceptional point is discrete we conclude that
there are only finitely many exceptional points. Then we have the following
result about the structure of Γk.
Proposition 3.3. For every k ∈ {−q + 1, . . . , p − 1}, the set Γk is the
disjoint union of a finite number of open analytic arcs and a finite number
of exceptional points. The set Γk has no isolated points.
Proof. This was proved for k = 0 in [10] and [1, Theorem 11.9]. For general
k, there are only finitely many exceptional points and the proof follows in a
similar way. �
3.2. The Riemann surface. From Proposition 3.3 it follows that the curves
Γk can be taken as cuts for the p + q-sheeted Riemann surface of the alge-
braic equation (1.7). We number the sheets from 1 to p+ q, where the kth
sheet of the Riemann surface is
Rk = {λ ∈ C | |zk−1(λ)| < |zk(λ)| < |zk+1(λ)|} = C \ (Γ−q+k−1 ∪ Γ−q+k).
(3.2)
Thus zk is well-defined and analytic on Rk.
The easiest case to visualize is the case where consecutive cuts are disjoint,
that is, Γ−q+k−1 ∩ Γ−q+k = ∅ for every k = 2, . . . , p+ q− 2. In that case we
have that Rk is connected to Rk+1 via Γ−q+k in the usual crosswise manner,
and zk+1 is the analytic continuation of zk across Γ−q+k.
The general case is described in the following proposition.
Proposition 3.4. Suppose A is an open analytic arc such that A ⊂ Γ−q+k,
for k = k1, . . . , k2, and A ∩ (Γ−q+k1−1 ∪ Γ−q+k2+1) = ∅. Then for k =
EIGENVALUES OF BANDED TOEPLITZ MATRICES 11
k1, . . . , k2+1, we have that the analytic continuation of zk across A is equal
to zk1+k2−k+1. Thus across A, we have that Rk is connected to Rk1+k2−k+1.
Proof. We have that
|zk1(λ)| = |zk1+1(λ)| = · · · = |zk2(λ)| = |zk2+1(λ)|
for λ ∈ A, with strict inequalities (<) for λ on either side of A. Choose an
orientation for A. Then there is a permutation π of {k1, . . . , k2 + 1} such
that zπ(k) is the analytic continuation of zk from the +-side of A to the
−-side of A.
Assume that there are k, k′ ∈ {k1, . . . , k2 + 1} such that k < k′ and
π(k) < π(k′). Take a regular λ0 ∈ A and a small neighborhood U of λ0 such
that A∩U = Γ−q+k ∩U = Γ−q+k′ ∩U and A∩U is an analytic arc starting
and terminating on ∂U . Then we have a disjoint union U = U+∪U−∪(A∩U)
where U+ (U−) is the part of U on the +-side (−-side) of A. The function
φ defined by
φ(λ) =
zk(λ)
zk′ (λ)
, for λ ∈ U+,
zπ(k)(λ)
zπ(k′)(λ)
, for λ ∈ U−,
has an analytic continuation to U , and satisfies |φ(λ)| < 1 for λ ∈ U+ ∪ U−
and |φ(λ)| = 1 for λ ∈ A ∩ U . This contradicts the maximum principle for
analytic functions. Therefore π(k) > π(k′) for every k, k′ ∈ {k1, . . . , k2 + 1}
with k < k′, and this implies that π(k) = k1 + k2 − k + 1 for every k =
k1, . . . , k2 + 1, and the proposition follows. �
3.3. The functions wk(λ). A major role is played by the functions wk,
which for k ∈ {−q + 1, . . . , p − 1}, are defined by
wk(λ) =
zj(λ), for λ ∈ C \ Γk. (3.3)
Note that wk = (−1)p−ka−1p w{1,...,k} in the notation of (2.23).
Proposition 3.5. The function wk is analytic in C \ Γk.
Proof. Since zj is analytic on Rj = C \ (Γ−q+j−1 ∪ Γ−q+j), see (3.2), we
obtain from its definition that wk is analytic in C \
j=1 Γ−q+j. Let A be
an analytic arc in Γ−q+j \ Γk for some j < k + q. Choose an orientation on
A. Since the arc is disjoint from Γk, we have that zj+(λ) = zπ(j)−(λ), for
λ ∈ A and j = 1, . . . , q+k, where π is a permutation of {1, . . . , q+k}. Since
wk is symmetric in the zj ’s for j = 1, . . . , q + k, it then follows that
wk+(λ) = wk−(λ), for λ ∈ A,
which shows the analyticity in C \Γk with the possible exception of isolated
singularities at the exceptional points of Γ−q+1, Γ−q+2, . . . , Γk−1. However,
each zj , and therefore also wk, is bounded near such an exceptional point,
so that any isolated singularity is removable. �
12 MAURICE DUITS AND ARNO B.J. KUIJLAARS
In the rest of the paper we make frequently use of the logarithmic de-
rivative w′k/wk of wk. By the fact that wk does not vanish on C \ Γk and
Proposition 3.5, it follows that w′k/wk is analytic in C \ Γk. By Proposition
3.4 it moreover has an analytic continuation across every open analytic arc
A ⊂ Γk. Near the exceptional points that are no branch points w′k/wk re-
mains bounded. At the branch points it can however have singularities of a
certain order.
Proposition 3.6. Let λ0 ∈ Γk be a branch point of Γk. Then there exists
an m ∈ N such that
w′k(λ)
wk(λ)
(λ− λ0)−m/(m+1)
, (3.4)
as λ → λ0 with λ ∈ C \ Γk.
Proof. Let 1 ≤ j ≤ q+k. We investigate the behavior of zj(λ) when λ → λ0
such that λ remains in a connected component of C \ (Γj−1 ∪ Γj). Then
zj(λ) → z0 for some z0 ∈ C with a(z0) = λ0. Let m0 + 1 be the multiplicity
of z0 as a solution of a(z) = λ0. Then
a(z) = λ0 + c0(z − z0)m0+1(1 +O(z − z0)), z → z0, (3.5)
for some nonzero constant c0. Therefore,
zj(λ) = z0 +O((λ− λ0)1/(m0+1)), (3.6)
z′j(λ) = O((λ− λ0)−m0/(m0+1)), (3.7)
for λ → λ0 such that λ remains in the same connected component of C \
(Γj−1 ∪ Γj). Let m be the maximum of all the multiplicities of the roots of
a(z) = λ0. Then it follows from (3.6) and (3.7) that
z′j(λ)
zj(λ)
= O((λ− λ0)−m/(m+1))
as λ → λ0 with λ ∈ C \ Γk. Then we obtain (3.4) in view of (3.3). �
We end this section by giving the asymptotics of w′k/wk for λ → ∞.
Proposition 3.7. As λ → ∞ with λ ∈ C \ Γk, we have
w′k(λ)
wk(λ)
− q+k
λ−1 +O
λ−1−1/q
, for k = −q + 1, . . . ,−1,
−λ−1 +O(λ−2), for k = 0,
λ−1 +O
λ−1−1/p
, for k = 1, . . . , p − 1.
(3.8)
Proof. This follows directly from (3.1) and (3.3). �
EIGENVALUES OF BANDED TOEPLITZ MATRICES 13
4. Proof of Theorem 2.3
We use the function wk introduced in (3.3). We define µk by the formula
(2.9) and we note that
dµk(λ) =
w′k+(λ)
wk+(λ)
w′k−(λ)
wk−(λ)
dλ. (4.1)
Proposition 4.1. For each k = −q + 1, . . . , p − 1, we have that µk is a
measure on Γk with total mass µk(Γk) = (q + k)/q if k ≥ 0, and µk(Γk) =
(p− k)/p if k ≥ 0.
Proof. We first show that µk is a measure, i.e., that it is non-negative on
each analytic arc of Γk. Let A be an analytic arc in Γk consisting only of
regular points. Let t 7→ λ(t) be a parametrization of A in the direction of
the orientation of Γk. Then
dµk(λ) =
w′k+(λ(t))
wk+(λ(t))
w′k−(λ(t))
wk−(λ(t))
λ′(t)dt
wk+(λ(t))
wk−(λ(t))
To conclude that µk is non-negative on A, it is thus enough to show that
Re log
wk+(λ)
wk−(λ)
= 0, for λ ∈ A, (4.2)
Im log
wk+(λ)
wk−(λ)
increases along A. (4.3)
Since |wk+(λ)| = |wk−(λ)| for λ ∈ A, we have (4.2) so that it only remains
to prove (4.3).
There is a neighborhood U of A such that U \ Γk has two components,
denoted U+ and U−, where U+ is on the +-side of Γk and U− on the −-side.
It follows from Proposition 3.4 that wk has an analytic continuation from
U− to U , which we denote by ŵk, and that |wk(λ)| < |ŵk(λ)| for λ ∈ U+,
and equality |wk+(λ)| = |ŵk(λ)| holds for λ ∈ A. Thus it follows that
Re log
wk(λ)
ŵk(λ)
≤ 0, for λ ∈ A,
where ∂
denotes the normal derivative to A in the direction of U+. Then by
the Cauchy-Riemann equations we have that Im log
wk+(λ)
ŵk+(λ)
is increasing
along A. Since ŵk+(λ) = wk−(λ) for λ ∈ A, we obtain (4.3). Thus µk is a
measure.
Next we show that µk is a finite measure, which means that we have to
show that
w′k+(λ)
wk+(λ)
w′k−(λ)
wk−(λ)
(4.4)
14 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Figure 1. Illustration for the proofs of Propositions 4.1 and
4.2. The solid line is a sketch of a possible contour Γk. The
dashed line is the contour Γ̃k,R and the dotted line is the
boundary of a disk of radius R around 0.
is integrable near infinity on Γk and near every branch point on Γk. This
follows from Propositions 3.7 and 3.6. Indeed, from Proposition 3.7 it follows
w′k+(λ)
wk+(λ)
w′k−(λ)
wk−(λ)
λ−1−δ
as λ → ∞, λ ∈ Γk. (4.5)
where δ = 1/q if k < 0 and δ = 1/p if k > 0. Since δ > 0 we see that
(4.4) is integrable near infinity. For a branch point λ0 of Γk, we have from
Proposition 3.6 that there exist an m ≥ 1 such that
w′k+(λ)
wk+(λ)
w′k−(λ)
wk−(λ)
(λ− λ0)−m/(m+1)
as λ → λ0, λ ∈ Γk. (4.6)
This shows that (4.4) is integrable near every branch point. Thus µk is a
finite measure.
Finally we compute the total mass of µk. LetD(0, R) = {z ∈ C | |z| < R}.
Then for R large enough, so that D(0, R) contains all exceptional points of
Γk and all connected components of C \ Γk (if any),
µk(Γk ∩D(0, R)) =
Γk∩D(0,R)
w′k+(λ)
wk+(λ)
Γk∩D(0,R)
w′k−(λ)
wk−(λ)
(4.7)
where we have used the behavior (4.6) near the branch points in order to
be able to split the integrals. Again using (4.6) we can then turn the two
integrals into a contour integral over a contour Γ̃k,R as in Figure 1. The
contour Γ̃k,R passes along the ±-sides of Γk ∩D(0, R) and if we choose the
orientation that is also shown in Figure 1 (and which is independent of the
EIGENVALUES OF BANDED TOEPLITZ MATRICES 15
choice of orientation for Γk), then
µk(Γk ∩D(0, R)) =
Γ̃k,R
w′k(λ)
wk(λ)
dλ. (4.8)
The parts of Γ̃k,R that belong to bounded components of C \Γk form closed
contours along the boundary of each bounded component. By Cauchy’s
theorem their contribution to the integral (4.8) vanishes. The parts of Γ̃k,R
that belong to the unbounded components of C \Γk can be deformed to the
circle ∂D(0, R) with the clockwise orientation. Thus if we use the positive
orientation on ∂D(0, R) as in Figure 1, then we obtain from (4.8)
µk(Γk ∩D(0, R)) = −
∂D(0,R)
w′k(λ)
wk(λ)
Letting R → ∞ and using Proposition 3.7, we then find that µk is a measure
on Γk with total mass µk(Γk) = (q + k)/q if k ≤ 0, and µk(Γk) = (p− k)/p
if k ≥ 0. �
The following proposition is the next step in showing that the measures
µk from (2.9) satisfy the equations (2.10).
Proposition 4.2. For k = −q + 1, . . . , p− 1, we have that
dµk(x)
x− λ =
w′k(λ)
wk(λ)
, for λ ∈ C \ Γk, (4.9)
log |λ− x| dµk(x) = − log |wk(λ)|+ αk, for λ ∈ C, (4.10)
where αk is the constant
log |a−q|+ kq log |a−q|, if k ≤ 0,
log |a−q| − kp log |ap|, if k ≥ 0.
(4.11)
Proof. To prove (4.9), we follow the same arguments as in the calculation
of µk(Γk) in the end of the proof of Proposition 4.1. Let λ ∈ C \ Γk, and
choose R > 0 as in the proof of Proposition 4.1. We may assume R > |λ|.
Then similar to (4.7) and (4.8) we can write
Γk∩D(0,R)
dµk(x)
x− λ =
Γ̃k,R
w′k(x)
wk(x)(x− λ)
where Γ̃k,R has the same meaning as in the proof of Proposition 4.1, see also
Figure 1. As in the proof of Proposition 4.1 we deform to an integral over
∂D(0, R), but now we have to take into account that the integrand has a
pole at x = λ with residue w′k(λ)/wk(λ). Therefore, by Cauchy’s theorem
Γk∩D(0,R)
dµk(x)
x− λ =
w′k(λ)
wk(λ)
∂D(0,R)
w′k(x)
wk(x)(x− λ)
dx. (4.12)
Letting R → ∞ and using Proposition 3.7 gives (4.9).
16 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Next we integrate (4.9) over a Jordan curve J in C \ Γk from λ1 to λ2.
x− λ dµk(x) dλ = −
∫ ∫ λ2
x− λ dλ dµk(x)
(log |λ1 − x| − log |λ2 − x|+ i∆J [arg(λ− x)]) dµk(x), (4.13)
where ∆J [arg(λ − x)] denotes the change in argument of λ − x as when λ
varies over J from λ1 to λ2. By (4.9) the integral (4.13) is equal to
w′k(λ)
wk(λ)
dλ = log |wk(λ2)| − log |wk(λ1)|+ i∆J [argwk(λ)]. (4.14)
Equating the real parts of (4.13) and (4.14) we get
(log |λ1 − x| − log |λ2 − x|) dµk(x) = − log |wk(λ1)|+ log |wk(λ2)|.
(4.15)
Since λ1 and λ2 can be taken arbitrarily in a connected component of C\Γk,
we find that there exists a constant αk ∈ R (which a priori could depend on
the connected component) such that
log |λ− x| dµk(x) = − log |wk(λ)|+ αk, (4.16)
for all λ in a connected component of C \ Γk. By continuity the equation
(4.16) extends to the closure of the connected component, which shows that
the same constant αk is valid for all connected components. Thus (4.16)
holds for all λ ∈ C.
The exact value of αk can then be determined by expanding (4.16) for
large λ. Suppose for example that k < 0. Then by (3.1) and (3.3)
|wk(λ)| =
|zj(λ)| = |a−q|(q+k)/q|λ|−(q+k)/q
1 +O(λ−1/q)
as λ → ∞. Thus
− log |wk(λ)| =
q + k
log |λ| − q + k
log |a−q|+O(λ−1/q). (4.17)
Since
log |λ− x| dµk(x) = log |λ|µk(Γk) + o(1) =
q + k
log |λ|+ o(1), (4.18)
as λ → ∞, the value (4.11) for αk follows from (4.16), (4.17), and (4.18). The
argument for k > 0 is similar. This completes the proof of the proposition.
To prove part (c) of Theorem 2.3 we also need the following lemma.
EIGENVALUES OF BANDED TOEPLITZ MATRICES 17
Lemma 4.3. Let ~ν1 = (ν1,−q+1 . . . , ν1,p−1) and ~ν2 = (ν2,−q+1 . . . , ν2,p−1) be
two admissible vectors of measures. Then J(~ν1 − ~ν2) is well defined and
J(~ν1 − ~ν2) ≥ 0, (4.19)
with equality if and only if ~ν1 = ~ν2.
Proof. Since both ~ν1 and ~ν2 have finite energy, we find that J(~ν1 − ~ν2) is
well defined. According to the alternative representation (2.12), we have
J(~ν1 − ~ν2) =
I(ν1,0 − ν2,0)
k(k + 1)I
ν1,−q+k
− ν2,−q+k
− ν1,−q+k+1
k + 1
ν2,−q+k+1
k + 1
k(k + 1)I
ν1,p−k
ν2,p−k
ν1,p−k−1
k + 1
ν2,p−k−1
k + 1
(4.20)
Using (2.6) and (2.8), we see that all terms in (4.20) are non-negative and
therefore (4.19) holds.
Suppose now that J(~ν1 − ~ν2) = 0. Then all terms in the right-hand side
of (4.20) are zero, so that
ν1,0 = ν2,0, (4.21)
ν1,−q+k
ν2,−q+k+1
k + 1
ν1,−q+k+1
k + 1
ν2,−q+k
, for k = 1, . . . , q − 1,
(4.22)
ν1,p−k
ν2,p−k−1
k + 1
ν1,p−k−1
k + 1
ν2,p−k
, for k = 1, . . . , p− 1.
(4.23)
Using (4.21) in (4.22) with k = q − 1, we find ν1,−1 = ν2,−1. Proceeding
inductively we then obtain from (4.22) that ν1,k = ν2,k for all k = −q +
1, . . . , 0. Similarly, from (4.21) and (4.23) it follows that ν1,k = ν2,k for
k = 0, . . . , p− 1, so that ~ν1 = ~ν2 as claimed. �
Now we are ready for the proof of Theorem 2.3.
Proof of Theorem 2.3. (a) In view of Proposition 4.1 it only remains to
show that µk ∈ Me for every k = −q + 1, . . . , p − 1. The decay estimate
(4.5) implies that
log(1 + |λ|) dµk(λ) < ∞.
The fact that I(µk) < +∞ follows from (4.10). Indeed,
I(µk) = −
log |λ− x|dµk(x)dµk(λ) =
(log |wk(λ)| − αk)dµk(λ)
18 MAURICE DUITS AND ARNO B.J. KUIJLAARS
and this is finite since µk is a finite measure on Γk with a density that decays
as in (4.5) and log |wk(λ)| is continuous on Γk and grows only as a constant
times log |λ| as λ → ∞. Thus ~µ is admissible and part (a) is proved.
(b) According to (4.10) we have
log |λ− x| dµk(x)−
log |λ− x| dµk+1(λ)−
log |λ− x| dµk−1(λ)
= −2 log |wk(λ)|+ 2αk + log |wk+1(λ)| − αk+1 + log |wk−1(λ)| − αk−1
= log
wk+1(λ)wk−1(λ)
wk(λ)
+ 2αk − αk+1 − αk−1
= log
zq+k+1(λ)
zq+k(λ)
+ 2αk − αk+1 − αk−1. (4.24)
Since |zq+k(λ)| = |zq+k+1(λ)| for λ ∈ Γk, we see from (4.24) that (2.10)
holds with constant
lk = 2αk − αk−1 + αk+1. (4.25)
Note that for k = −q + 1 and k = p − 1, we are using the convention that
µ−q = µp = 0, and we also have put α−q = αp = 0. This proves part (b).
(c) Let ~ν = (ν−q+1, . . . , νp−1) be any admissible vector of measures. From
the representation (2.13) we get
J(~ν) = J(~µ+ ~ν − ~µ)
= J(~µ) + J(~ν − ~µ) + 2
j,k=−q+1
AjkI(µj, νk − µk). (4.26)
Using (2.14), we find from (4.26)
J(~ν) = J(~µ) + J(~ν − ~µ) +
k=−q+1
I(2µk − µk−1 − µk+1, νk − µk) (4.27)
For each k = −q + 1, . . . , p − 1, we have
I(2µk − µk−1 − µk+1, νk − µk)
log |λ− x| d(2µk − µk−1 − µk+1)(x)
d(νk − µk)(λ) (4.28)
By (2.10) the inner integral in the right-hand side of (4.28) is constant for
λ ∈ Γk. Since νk and µk are finite measures on Γk with νk(Γk) = µk(Γk),
we find from (4.28) that
I(2µk − µk−1 − µk+1, νk − µk) = 0, for k = −q + 1, . . . , p− 1.
Then (4.27) shows that J(~ν) = J(~µ)+J(~ν−~µ), which by Lemma 4.3 implies
that J(~ν) ≥ J(~µ) and equality holds if and only if ~ν = ~µ. This completes
the proof of Theorem 2.3. �
EIGENVALUES OF BANDED TOEPLITZ MATRICES 19
5. Proofs of Proposition 2.5 and Theorem 2.6
5.1. Proof of Proposition 2.5. We will now prove Proposition 2.5, which
follows by a combinatorial argument.
Proof of Proposition 2.5. We prove (2.18) and (2.19) for k > 0. The case
k < 0 is similar. Let us first expand the determinant in the definition of
Pk,n(λ) = detTn(z
−k(a− λ)) =
(a− λ)j−π(j)+k. (5.1)
Here Sn denotes the set of all permutation on {1, . . . , n}. By the band struc-
ture of Tn(z
−k(a − λ)) it follows that we only have non-zero contributions
from permutations π that satisfy
k − p ≤ π(j) − j ≤ q + k, for all j = 1, . . . , n. (5.2)
Define for π ∈ Sn,
Nπ = {j | π(j) = j + k}. (5.3)
and denote the number of elements of Nπ by |Nπ|. For each π ∈ Sn we have
j=1(a−λ)j−π(j)+k is a polynomial in λ of degree at most |Nπ|. So by
(5.1)
dk,n = degPk,n ≤ max
|Nπ| (5.4)
where we maximize over permutations π ∈ Sn satisfying (5.2).
Let π ∈ Sn satisfying (5.2). We prove (2.18) by giving an upper bound
for |Nπ|. Since
j=1(π(j) − j) = 0 we obtain
(π(j) − j)+ =
(j − π(j))+, (5.5)
where (·)+ is defined as (a)+ = max(0, a) for a ∈ R. Each j ∈ Nπ gives a
contribution k to the left-hand side of (5.5). Therefore the left-hand side is
at least k|Nπ|. By (5.2) we have that each term in the right hand side is
at most p− k. Moreover, there are at most n− |Nπ| non-zero terms in this
sum. Combining this with (5.5) leads to
k|Nπ| ≤
(π(j) − j)+ =
(j − π(j))+ ≤ (n− |Nπ|)(p − k). (5.6)
Hence, if π is a permutation satisfying (5.2)
|Nπ| ≤
n(p− k)
. (5.7)
Now (2.18) follows by combining (5.7) and (5.4).
To prove (2.19), we assume that n ≡ 0 mod p. We claim that there exists
a unique π such that equality holds in (5.7). Then equality holds in both
20 MAURICE DUITS AND ARNO B.J. KUIJLAARS
inequalities of (5.6) and the above arguments show that this can only happen
π(j) = j + k, or π(j) = j − p+ k, (5.8)
for every j = 1, . . . , n. We claim that there exists a unique such permutation,
namely
π(j) =
j + k, if j ≡ 1, . . . , (p − k) mod p,
j − p+ k, if j ≡ (p− k + 1), . . . , p mod p.
(5.9)
To see this let π be a permutation satisfying (5.8). The numbers 1, . . . , p−
k can not satisfy π(j) = j−p+k and thus satisfy π(j) = j+k. On the other
hand, the numbers 1, . . . , k can not be the image of numbers j satisfying
π(j) = j + k, and thus π(j) = j − p + k for j = p − k + 1, . . . , p. So (5.9)
holds for j = 1, . . . , p. This means in particular that the restriction of π to
{p + 1, . . . , n} is again a permutation, but now on {p + 1, . . . , n}. By the
same arguments we then find that (5.9) holds for j = p + 1, . . . , 2p, and so
on. The result is that (5.9) is indeed the only permutation that satisfies
(5.8).
Finally, a straightforward calculation shows that the coefficient of λ(p−k)n/p
j=1(a − λ)j−π(j)+k with π as in (5.9) is nonzero and given by (2.19).
This proves the proposition. �
5.2. Proof of Theorem 2.6. Before we start with the proof of Theorem
2.6 we first prove the following proposition concerning the asymptotics for
Pk,n for n → ∞.
Proposition 5.1. Let Mk = {q + k + 1, . . . , p+ q}. We have that
Pk,n(λ) = (wMk(λ))
nCMk(λ) (1 +O(exp(−cKn)) , n → ∞, (5.10)
uniformly on compact subsets K of C \ Γk. Here cK is a positive constant
depending on K.
Proof. First rewrite (2.22) as
Pk,n(λ) = (wMk(λ))
nCMk(λ) (1 +Rk,n(λ)) . (5.11)
with Rk,n defined by
Rk,n(λ) =
M 6=Mk
(wM (λ))
nCM (λ)
(wMk(λ))
nCMk(λ)
. (5.12)
Let K be a compact subset of C \ Γk. If K does not contain branch points
then there exists A,B > 0 such that
A < |CM (λ)| < B (5.13)
for all λ ∈ K and M . Moreover, we have
wM (λ)
wMk(λ)
zq+k(λ)
zq+k+1(λ)
≤ sup
zq+k(λ)
zq+k+1(λ)
< 1, (5.14)
EIGENVALUES OF BANDED TOEPLITZ MATRICES 21
for all λ ∈ K and M 6= Mk. Therefore one readily verifies from (5.11)
that there exist cK such that |Rk,n(λ)| ≤ exp(−cKn) for all λ ∈ K and n
large enough. This proves the statement in case K does not contain branch
points.
Suppose that K does contain branch points. Without loss of generality
we can assume that all branch points lie in the interior of K (otherwise we
replace K by a bigger compact set). The boundary ∂K of K is a com-
pact set with no branch points and therefore (5.10) holds for ∂K by the
above arguments. Since wMk and CMk are analytic in K, we find by (5.11)
that Rk,n is analytic in K. The maximum modulus principle for analytic
functions states that supz∈K |Rk,n(z)| = supz∈∂K |Rk,n(z)| and thereby we
obtain that (5.10) also holds for K with the same constant cK = c∂K . �
We now state two particular consequences of (5.10).
Corollary 5.2. Let k ∈ {−q + 1, . . . , p − 1}. For every compact set K ⊂
C \ Γk we have that µk,n(K) = 0 for n large enough.
Proof. Let K be a compact subset of C \ Γk. By (5.10) it follows that Pk,n
has no zeros in K for large n. Since nµk,n(K) equals the number of zeros of
Pk,n in K the corollary follows. �
Corollary 5.3. Let k ∈ {−q + 1, . . . , p− 1}. We have that
dµk,n(x)
x− λ =
dµk(x)
x− λ , (5.15)
uniformly on compact subsets of C \ Γk.
Proof. Let K be a compact subset of C \ Γk. Note that
dµk,n(x)
x− λ =
λi∈spk Tn(a)
λi − λ
P ′k,n(λ)
nPk,n(λ)
, (5.16)
for all λ ∈ K. With Mk and cK as in Proposition 5.1 we obtain from (5.10)
P ′k,n(λ)
nPk,n(λ)
w′Mk(λ)
wMk(λ)
+O(1/n), n → ∞, (5.17)
uniformly on K. Let us rewrite the right-hand side of (5.17). By expanding
both sides of zq(a(z)− λ) = ap
j=1(z − zj(λ)) and collecting the constant
terms we obtain
(−zj(λ)) =
. (5.18)
Since λ /∈ Γk, we can split this product in two parts, take the logarithmic
derivative and use (3.3) and (2.23) to obtain
z′j(λ)
zj(λ)
j=q+k+1
z′j(λ)
zj(λ)
w′k(λ)
wk(λ)
w′Mk(λ)
wMk(λ)
. (5.19)
22 MAURICE DUITS AND ARNO B.J. KUIJLAARS
Combining (5.16), (5.17) and (5.19), we obtain
dµk,n(x)
x− λ =
w′k(λ)
wk(λ)
(5.20)
uniformly on K. Then (5.15) follows from (5.20) and (4.9). �
Now we are ready for the proof of Theorem 2.6.
Proof of Theorem 2.6.
First we prove (2.21). By Proposition 2.5 and the fact that ~µ is admissible,
we get (see (2.8))
µk,n(C) =
degPk,n ≤ µk(C), (5.21)
for every n ∈ N.
Let C0(C) be the Banach space of continuous functions on C that vanish
at infinity. The dual space C0(C)
∗ of C0(C) is the space of regular complex
Borel measures on C. By (5.21) the sequence (µk,n)n∈N belongs to the ball in
C0(C)
∗ centered at the origin with radius µk(C), which is weak
∗ compact by
the Banach-Alaoglu theorem. Let µk,∞ be the limit of a weak
∗ convergent
subsequence of (µk,n)n∈N.
By weak∗ convergence and Corollary 5.2 we obtain that µk,∞ is supported
on Γk. Combining this with (5.15) and the weak
∗ convergence leads to
dµk(x)
x− λ =
dµk,∞(x)
x− λ , (5.22)
for every λ ∈ C \ Γk. The integrals in (5.22) are known in the literature as
the Cauchy transforms of the measures µk and µk,∞. The Cauchy transform
on Γk is an injective map that maps measures on Γk to functions that are
analytic in C \ Γk (one can find explicit inversion formulae, see for example
the arguments in [9, Theorem II.1.4] or the Stieltjes-Perron inversion formula
in the special case Γk ⊂ R). Thus it follows from (5.22) that µk,∞ = µk.
Therefore
µk,n = µk (5.23)
in the sense of weak∗ convergence in C0(C)
∗. Thus (2.21) holds if φ is a
continuous function that vanishes at infinity.
From (5.21) and (5.23) it also follows that
µk,n(C) = µk(C), (5.24)
Then the sequence (µk,n)n∈N is tight. That is, for every ε > 0 there exists a
compact K such that µk,n(C \K) < ε for every n ∈ N. By a standard ap-
proximation argument one can now show that (2.21) holds for every bounded
continuous function φ on C.
Having (2.21) and Proposition 5.1, we can prove (2.20) as in [1, Theo-
rem 11.17]. Indeed, the sets lim infn→∞ spk Tn(a) and lim supn→∞ spk Tn(a)
equal the support of µk, which is Γk. �
EIGENVALUES OF BANDED TOEPLITZ MATRICES 23
–2 –1.5 –1 –0.5 0.5 1 1.5
lambda
–2 –1.5 –1 –0.5 0.5 1 1.5
lambda
Figure 2. Illustration for Example 1: The densities of the
measures µ0 (left) and µ1 (right) for a =
4(z+1)3
6. Examples
6.1. Example 1. As a first example consider the symbol a defined by
a(z) =
4(z + 1)3
. (6.1)
In this case we have p = 2 and q = 1. So we obtain two contours Γ0 and
Γ1 with two associated measures µ0 and µ1. This example appeared in [3],
in which the authors gave explicit expressions for Γ0 and µ0. The following
proposition also contains expressions for Γ1 and µ1. In what follows we take
the principal branches for all fractional powers.
Proposition 6.1. With a as in (6.1), we have that Γ0 = [0, 1] and
dµ0(λ) =
dλ. (6.2)
Moreover, Γ1 = (−∞, 0] and
dµ1(λ) =
)1/3 −
1− λ− 1
(−λ)2/3
dλ. (6.3)
Proof. A straightforward calculation shows that λ = 0 and λ = 1 are the
branch points.
Let λ ∈ Γ0 ∪ Γ1 and assume that λ is not a branch point. There exist
y1, y2 ∈ C such that y1 6= y2, |y1| = |y2| and a(y1) = a(y2) = λ. Then it
follows from (6.1) that |y1+1| = |y2+1|. Therefore y1 and y2 are intersection
points of a circle centered at −1 and a circle centered at the origin. Since
y1 6= y2, this means that y1 = y2 and therefore λ = a(y1) = a(y2) = a(y1) =
λ, so that λ ∈ R. A further investigation shows that a(z)−λ has 3 different
real zeros if λ > 1. If λ < 1 and λ 6= 0 then a(z) − λ has precisely 1 real
zero and 2 conjugate complex zeros. Therefore, Γ0 ∪ Γ1 = (−∞, 1].
24 MAURICE DUITS AND ARNO B.J. KUIJLAARS
−2 −1.5 −1 −0.5 0 0.5 1 1.5 2
k = 0
−2 −1.5 −1 −0.5 0 0.5 1 1.5 2
k = 1
Figure 3. Illustration for Example 1: The spectrum
spT50(a) (top) and the generalized spectrum sp1T50(a) (bot-
tom), for the symbol a =
4(z+1)3
Now we will show that Γ0 = [0, 1] and Γ1 = (−∞, 0]. By Cardano’s
formula the solutions of the algebraic equation a(z) = λ are given by
zj(λ) = −1−
3λ1/3
1 + (1− λ)1/2
+ ω−j
1− (1− λ)1/2
(6.4)
for λ ∈ [0, 1] and
zj(λ) = −1+
3(−λ)1/3
1 + (1− λ)1/2
− ω−j−2
(1− λ)1/2 − 1
(6.5)
for λ ∈ (−∞, 0]. Here ω = e2πi/3. One can check that |z1(λ)| = |z2(λ)| <
|z3(λ)| for λ ∈ (0, 1] and |z1(λ)| < |z2(λ)| = |z3(λ)| for λ ∈ (−∞, 0). More-
over, for λ = 0 we have z1(0) = z2(0) = z3(0) = −1. Therefore Γ0 = [0, 1]
and Γ1 = (−∞, 0].
The density (6.2) was already given in [3] and (6.3) follows in a similar
way. �
In Figure 2 we plotted the densities of µ0 and µ1. Note that, due to the
interaction between µ0 and µ1 in the energy functional, there is more mass
of µ0 near 0 than near 1. We also see that the singularities of the densities
for µ0 and µ1 are of order O(|λ|−2/3) for λ → 0, whereas the typical nature
of a singularity in each of the measures is a square root singularity. The
stronger singularity is due to the fact that a(z) − λ has a triple root for
λ = 0.
EIGENVALUES OF BANDED TOEPLITZ MATRICES 25
–4 –2 2 4
lambda
–4 –2 2 4
lambda
Figure 4. Illustration for Example 2: The densities of the
measures µ0 (left) and µ1 = µ−1 (right) for a(z) = z
2 + z +
z−1 + z−2.
In Figure 3 we plotted the eigenvalues and generalized eigenvalues for
n = 50. It is known that the eigenvalues are simple and positive [3, §2.3],
which we also see in Figure 3.
6.2. Example 2. For the symbol a defined by
a(z) = z2 + z + z−1 + z−2. (6.6)
we have p = q = 2. From the symmetry a(1/z) = a(z) it follows that
Γ−1 = Γ1 and µ−1 = µ1.
The interesting feature of this example is that the contours Γ0 and Γ±1
overlap. To be precise, the interval (−9/4, 0) is contained in all three con-
tours Γ−1,Γ0 and Γ1. This can be most easily seen by investigating the
image of the unit circle under a. Consider
a(eit) = 2 cos 2t+ 2cos t, for t ∈ [0, 2π). (6.7)
A straightforward analysis shows that for every λ ∈ (−9/4, 0), the equation
a(eit) = λ has four different solutions for t in [0, 2π). This means that the
four solutions of the equation a(z) = λ are on the unit circle, and so in
particular have the same absolute value.
The equation a(z) − λ = 0 can be explicitly solved by introducing the
variable y = z + 1/z. In exactly the same way as in the previous example
one can obtain the limiting measures. We will not give the explicit formulas,
but only plot the densities in Figure 4. The branch points are λ = −9/4,
λ = 0 and λ = 4. The contours are given by
Γ0 = [−9/4, 4], Γ−1 = Γ1 = (−∞, 0]. (6.8)
The densities have singularities at the branch points in the interior of their
supports. The singularities are only felt at one side of the branch points.
Consider first µ0, whose density has a singularity at 0. However the limiting
value when 0 is approached from the positive real axis is finite. The change
in behavior of µ0 has to do with the fact that z1 is analytic on (0, 4) but not
26 MAURICE DUITS AND ARNO B.J. KUIJLAARS
on (−9/4, 0). Therefore we find by (1.12) that
dµ0(λ) =
z1+(λ)
z2+(λ)
z1−(λ)
z2−(λ)
dλ (6.9)
on (−9/4, 0), and
dµ0(λ) =
z2+(λ)
z2−(λ)
dλ (6.10)
on (0, 4).
For µ−1 = µ1 a similar phenomenon happens at λ = −9/4. This is a
consequence of the fact that z1 has an analytic continuation into z2 when
we cross (−∞,−9/4), but it has an analytic continuation into z4 when we
cross (−9/4, 0).
6.3. Example 3. As a final example, consider the symbol
a(z) = zp + z−q, (6.11)
with p, q ≥ 1 and gcd(p, q) = 1. This example appeared in [10], where the
authors mentioned that Γ0 is given by the star
Γ0 = {rωj | j = 1, . . . , p+ q, 0 ≤ r ≤ R} (6.12)
with ω = e2πi/(p+q) and R = (p + q)p−p/(p+q)q−q/(p+q). The other contours
also have a star shape, namely
Γk = {(−1)krωj | j = 1, . . . , p+ q, 0 ≤ r < ∞} (6.13)
for k 6= 0. Note that the star Γk for k 6= 0 is unbounded.
In Figure 5 we plotted the eigenvalues and the generalized eigenvalues
for p = 2, q = 3 and n = 50. All the (generalized) eigenvalues appear
to lie exactly on the contours. In the special case p = 1 it is known that
the eigenvalues of Tn(a) lie indeed precisely on the star (6.12) and are all
simple (possibly except for 0) [4, Theorem 3.2], see also [6] for a connection
to Chebyshev-type quadrature.
6.4. Numerical stability. In Figure 3 and Figure 5 the eigenvalues and
the generalized eigenvalues of T50(a) were computed numerically. To control
the stability of the numerical computation of the eigenvalues one needs to
analyze the pseudo-spectrum. For banded Toeplitz matrices the pseudo-
spectrum is well understood [12, Th. 7.2]. To this date, a similar analysis
of the pseudo-spectrum for the matrix pencil (Tn(z
−ka), Tn(z
−k)) has not
been carried out. See [12, §X.45] for some remarks on the pseudo-spectrum
for the generalized eigenvalue problem.
EIGENVALUES OF BANDED TOEPLITZ MATRICES 27
−5 0 5
k = −2
−5 0 5
k = −1
−5 0 5
k = 0
−5 0 5
k = 1
Figure 5. Illustration for Example 3: The contours Γk and
the eigenvalues and generalized eigenvalues for T50(a) for the
symbol a = z2 + z−3.
References
1. A. Böttcher and S. M. Grudsky, Spectral Properties of Banded Toeplitz Matrices,
SIAM, Philadelphia, PA, 2005.
2. A. Böttcher and S. M. Grudsky, Can spectral values sets of Toeplitz band matrices
jump?, Linear Algebra Appl., 351-352 (2002), pp. 99-116.
3. E. Coussement, J. Coussement and W. Van Assche, Asymptotic zero distribution for
a class of multiple orthogonal polynomials, Trans. Amer. Math. Soc., (to appear)
4. M. Eiermann and R. Varga, Zeros and local extreme points of Faber polynomials
associated with hypocycloidal domains, Electron. Trans. Numer. Anal., 1 (1993), pp.
49-71.
5. I. I. Hirschman, Jr., The spectra of certain Toeplitz matrices, Illinois J. Math., 11
(1967), pp. 145-159.
6. A. Kuijlaars, Chebyshev quadrature for measures with a strong singularity, J. Comput.
Appl. Math., 65 (1995), pp. 207-214.
7. E. Nikishin and V. Sorokin, Rational Approximations and Orthogonality, Translations
of Mathematical Monographs 92, American Mathematical Society, Providence, RI,
(1991).
8. T. Ransford, Potential Theory in the Complex Plane, London Mathematical Society
Student Texts 28, Cambridge University Press, Cambridge, 1995.
9. E.B. Saff and V. Totik, Logartihmic Potentials with External Fields, Grundlehren der
Mathematischen Wissenschaften 316, Springer-Verlag, Berlin, 1997.
10. P. Schmidt and F. Spitzer, The Toeplitz matrices of an arbitrary Laurent polynomial,
Math. Scand., 8 (1960), pp. 15-38.
28 MAURICE DUITS AND ARNO B.J. KUIJLAARS
11. P. Simeonov, A weigthed energy problem for a class of admissible weights, Houston
J. Math., 31 (2005), pp. 1245-1260.
12. L.N. Trefethen and M. Embree, Spectra and Pseudospectra, Princeton University
Press, Princeton, NJ, 2005.
13. J.L. Ullman, A problem of Schmidt and Spitzer, Bull. Amer. Math. Soc., 73 (1967),
pp. 883-885.
14. H. Widom, On the eigenvalues of certain Hermitean operators, Trans. Amer. Math.
Soc., 88 (1958), pp. 491-522.
1. Introduction
2. Statement of results
2.1. The energy functional
2.2. The measures k as limiting measures of generalized eigenvalues
2.3. Overview of the rest of the paper
3. Preliminaries
3.1. The structure of the curves k
3.2. The Riemann surface
3.3. The functions wk()
4. Proof of Theorem 2.3
5. Proofs of Proposition 2.5 and Theorem 2.6
5.1. Proof of Proposition 2.5
5.2. Proof of Theorem 2.6
6. Examples
6.1. Example 1
6.2. Example 2
6.3. Example 3
6.4. Numerical stability
References
|
0704.0379 | Capturing knots in polymers | FIG. 1. Knotted bead-spring polymer: Starting configuration with N=16384 beads; after 6 reduction steps (N=265); final
configuration after 15 iterations (N=8) with the knotted (trefoil) region circled in red; and magnified.
(enhanced online)
Capturing knots in polymers
Peter Virnau, Mehran Kardar
Department of Physics, MIT, Cambridge, MA
02139-4307, USA
Yacov Kantor
School of Physics and Astronomy, Tel Aviv
University, 69978 Tel Aviv, Israel
(received, published)
[DOI: 10.1063/1.2130690]
Visualizing topological properties is a particularly
challenging task. Although algorithms can usually
determine if a loop contains a knot, finding its
exact location is difficult (and not necessarily
well-defined).
Here, we apply a reduction method by Koniaris
and Muthukumar
, which was originally proposed
to simplify polymers before calculating knot
invariants. We start with one end and consider
consecutive triangles formed by three adjacent
monomers. If the triangle is not crossed by any of
the remaining bonds, the particle in the middle is
removed. Going back and forth between both ends
we proceed until the configuration cannot be
reduced any further (see Fig.1).
Although the method is not perfect (sometimes
entangled, but unknotted regions remain), it
provides us with a valuable impression on the
typical number of knots, their respective location
and sizes
This work was supported by the DFG grant
Vi237/1.
P. Virnau, Y. Kantor, and M. Kardar, J. Am. Chem. Soc.,
in press (2005).
W. G. Taylor, Nature 406, 916 (2000).
K. Koniaris and M. Muthukumar, J.Chem.Phys. 95, 2873
(1991).
Pictures and movie were generated using the VMD
visualization package; see W. Humphrey, A. Dalke, and K.
Schulten,, J. Molec. Graphics 14, 33 (1996).
Copyright (2005) American Institute of Physics.
This article may be downloaded for personal use
only. Any other use requires prior permission of
the author and the American Institute of Physics.
The following article appeared in the Gallery of Images in
Chaos 15, 041103 (2005)
and may be found at
http://chaos.aip.org/chaos/gallery/toc_Dec05.jsp
This version also contains a movie of the algorithm.
|
0704.0380 | Exponential growth rates in a typed branching diffusion | Exponential growth rates in a typed branching diffusion
The Annals of Applied Probability
2007, Vol. 17, No. 2, 609–653
DOI: 10.1214/105051606000000853
c© Institute of Mathematical Statistics, 2007
EXPONENTIAL GROWTH RATES IN A TYPED
BRANCHING DIFFUSION
By Y. Git, J. W. Harris1 and S. C. Harris
Cambridge University, University of Bristol and University of Bath
We study the high temperature phase of a family of typed branch-
ing diffusions initially studied in [Astérisque 236 (1996) 133–154] and
[Lecture Notes in Math. 1729 (2000) 239–256 Springer, Berlin]. The
primary aim is to establish some almost-sure limit results for the long-
term behavior of this particle system, namely the speed at which the
population of particles colonizes both space and type dimensions, as
well as the rate at which the population grows within this asymp-
totic shape. Our approach will include identification of an explicit
two-phase mechanism by which particles can build up in sufficient
numbers with spatial positions near −γt and type positions near
t at large times t. The proofs involve the application of a variety
of martingale techniques—most importantly a “spine” construction
involving a change of measure with an additive martingale. In ad-
dition to the model’s intrinsic interest, the methodologies presented
contain ideas that will adapt to other branching settings. We also
briefly discuss applications to traveling wave solutions of an associ-
ated reaction–diffusion equation.
1. Introduction. In this article we will consider a certain family of typed
branching diffusions that have particles which move (independently of each
other) in space according to a Brownian motion with variance controlled by
the particle’s type process. The type of each particle evolves as an Ornstein–
Uhlenbeck process and this type also controls the rate at which births occur.
The particular form of this model permits many explicit calculations, but
throughout we will strive to develop techniques that rely on general prin-
ciples as much as possible, so they might readily adapt to other situations.
This model was previously considered in [12, 13]; these papers form essential
foundations for this work, although we will recall various results as necessary.
Received December 2004; revised November 2006.
1Supported in part by an EPSRC studentship.
AMS 2000 subject classification. 60J80.
Key words and phrases. Spatial branching process, branching diffusion, multi-type
branching process, additive martingales, spine decomposition.
This is an electronic reprint of the original article published by the
Institute of Mathematical Statistics in The Annals of Applied Probability,
2007, Vol. 17, No. 2, 609–653. This reprint differs from the original in pagination
and typographic detail.
http://arxiv.org/abs/0704.0380v1
http://www.imstat.org/aap/
http://dx.doi.org/10.1214/105051606000000853
http://www.imstat.org
http://www.imstat.org
http://www.imstat.org/aap/
http://dx.doi.org/10.1214/105051606000000853
2 Y. GIT, J. W. HARRIS AND S. C. HARRIS
We will make some significant applications of the spine theory for branch-
ing processes. Inspired by the series of papers Lyons, Pemantle and Peres
[19], Lyons [18] and Kurtz, Lyons, Pemantle and Peres [16], spine techniques
have been instrumental in recent years in providing intuitive and elegant
proofs of many important classical and new results in the theory of branch-
ing processes. In this article we use the recent reformulation of the spine
method presented in [8], which follows in similar spirit to the branching
Brownian motion study of Kyprianou [17]. For a selection of other applica-
tions of spine techniques, for example, see [1, 6, 7, 23] and references therein.
1.1. The branching model. We define Nt to be the set of particles alive
at time t ≥ 0. For a particle u ∈ Nt, Xu(t) ∈ R is its spatial position, and
Yu(t) ∈ R is the type of u. We will label offspring using the Ulam–Harris
convention where, for example, if u=∅21 then particle u is the first child
of the second child of the initial ancestor, and we will write v > u if particle
v is a descendant of particle u. The configuration of the branching diffusion
at time t is given by the point process Xt := {(Xu(t), Yu(t)) :u ∈Nt}.
A particle’s type evolves as an Ornstein–Uhlenbeck process with an invari-
ant measure given by the standard normal density φ(y) and an associated
differential operator (generator)
Qθ :=
− y ∂
where θ > 0 is considered to be the temperature of the system. The spatial
motion of a particle of type y is a driftless Brownian motion on R with
variance
A(y) := ay2, where a≥ 0.
A particle of type y particle is replaced by two offspring at a rate
R(y) := ry2 + ρ, where r, ρ≥ 0.
Each offspring inherits its parent’s current type and spatial position, and
then moves off independently of all others. We use P x,y and Ex,y with x, y ∈
R to represent probability and expectation when the Markov process starts
with a single particle at position (x, y).
We will find the almost sure rate of exponential growth, D(γ,κ), of par-
ticles which are found simultaneously with spatial positions near −γt and
type positions near κ
t at large times t. From this we can deduce the speed
of extremal particles and hence the asymptotic shape of the particle system.
The main effort is required in identifying D(γ,κ) as the almost sure limit of
t−1 log
1{Xu(t)≤−γt;Yu(t)≥κ
A TYPED BRANCHING DIFFUSION 3
In particular, the convergence properties of two different families of addi-
tive martingales associated with the branching diffusion will lead directly to
the spatial exponential growth rates and an upper bound on the space-type
growth. For the remaining lower bound, we describe an explicit two-phase
mechanism for amassing the required number of particles with prescribed
space-type positions. The first phase involves building up an “excess” num-
ber of particles, each covering a certain proportion of the required spatial
distance. During their second phase, enough of these particles must succeed
in making a difficult and rapid ascent into the required position. The latter
phase is proved using an intuitive change of measure technique that induces
a spine construction.
The family of models we are considering is specific but nevertheless have
some features of fundamental significance that motivate the choices for Qθ,
R and A. If the spatial motion is ignored, we have investigated a binary
branching Ornstein–Uhlenbeck process in a quadratic breeding potential. In
contrast, Enderle and Hering [5] considered a branching Ornstein–Uhlenbeck
with constant branching rate but random offspring distribution. A quadratic
breeding potential is a critical rate for explosions in the population of parti-
cles. In a branching Brownian motion on R with binary splitting occurring
at rate xp at position x, the population will explode almost surely in finite
time if p > 2, whereas for p = 2 the expected number of particles explodes
while the total population remains finite for all time with probability 1 (see
[15], Chapter 5.12). The Ornstein–Uhlenbeck process is not only a canoni-
cal ergodic diffusion, but this type-motion has exactly the right drift to help
counteract the quadratic breeding rate. For high temperatures, θ > 8r, there
is a sufficiently strong mean-reversion in the type processes to ensure that
the expected total population size does not blow up; but for temperatures
θ ≤ 8r, the quadratic breeding overpowers the pull toward the origin, the
expected population blows up in a finite time and particles behave very dif-
ferently. Throughout this paper we consider only high temperatures θ > 8r,
deferring the low and critical temperature regimes to future work. Given
other choices, the quadratic spatial diffusion coefficient now becomes very
natural, enabling us to find explicit families of (fundamental) additive mar-
tingales since the linearized traveling-wave equation can be linked to the
classical harmonic oscillator equations from physics. The binary branching
mechanism was taken for simplicity; in principle our approach could extend
to general offspring distributions, although new features would arise from
possible extinctions and necessary offspring moment conditions. All these
choices make the models rich in structure, possessing some very challenging
features whilst remaining sufficiently tractable.
4 Y. GIT, J. W. HARRIS AND S. C. HARRIS
1.2. Application to reaction–diffusion equations. Following in the foot-
steps of McKean [20], the solution of the reaction–diffusion equation
+R(y)u(u− 1) + θ
− y∂u
with initial condition f(x, y)∈ [0,1] for all x, y ∈R, can be represented by
u(t, x, y) =Ex,y
f(Xu(t), Yu(t))
Of great importance for reaction–diffusion equations are traveling-wave so-
lutions (e.g., see [21]). In the present context, a solution to equation (1) of
the form u(t, x, y) := w(x− ct, y) is said to be a traveling-wave of speed c,
where w(x, y) solves the traveling-wave equation
+R(y)w(w − 1) + θ
− y∂w
= 0.(3)
Fundamental to our study of the branching diffusion are two families of “ad-
ditive” martingales, Z±λ (t) [defined at (6)], which are linked to the lineariza-
tion of (1). When θ > 8r, Harris and Williams [13] determined when Z−λ is
uniformly integrable (see Theorem 17) and then wλ(x, y) :=E
x,y exp(−Z−λ (∞))
yields a traveling wave of speed c−λ . This gives the existence of traveling
waves for all speeds c greater than some threshold c̃(θ) := inf c−λ .
Furthermore, combining the McKean representation (2) with the almost-
sure convergence result established in [12] (look ahead to Theorem 18) can
give results on the attraction toward traveling waves from given initial data.
For example, if − lnf(x, y)∼ eλxg(y) uniformly in y as x→∞ for some suit-
able g ∈ L2(φ), the solution u(x, y) to (1) with initial conditions f satisfies
u(t, x− c−λ t, y)→wλ(x+ x̂, y) as t→∞, where x̂ is some constant that can
be determined from g.
In future work we hope to develop the approach used for standard BBM
and the FKPP equation in [11], and prove that traveling waves of a given
speed c > c̃(θ) are unique (up to translation) and that no traveling waves
exist for speeds c < c(θ). We anticipate that our new results on the growth
rates of particles will aid in establishing some difficult estimates on the tail
behavior of any traveling wave, and hence assist in proving the conjectured
uniqueness. In addition, we expect our growth rate results will be essential
in obtaining broader classes of initial conditions that are attracted toward
traveling waves. In each of these problems, difficulties arise from the un-
bounded type space where, for example, some control must be gained over
the possible contributions to
u∈Nt log f(Xu(t) − ct, Yu(t)) from particles
that have large type positions in addition to large spatial positions.
A TYPED BRANCHING DIFFUSION 5
2. Main results. In this section, we will present our main results that
identify the growth rates found within the branching diffusion. We will give
an overview for our proofs, identifying the key ideas and techniques used,
as well as introducing some intuition for the dominant behavior of particles
that underpins our approach.
2.1. Martingales. The principal tools used throughout this paper are
two fundamental families of “additive” martingales, which were introduced
in [13].
Before defining the martingales we give some key definitions. Let
λmin :=−
θ− 8r
Let λ ∈R, with the following convention which we always use for λ:
λmin <λ< 0.
Also, define
µλ :=
θ(θ− 8r− 4aλ2), ψ±λ :=
E±λ := ρ+ θψ
λ , c
λ :=−E
λ /λ.(5)
Will will occasionally write E±λ as E
±(λ) in order to emphasize that E±λ are
really functions of λ; the ± superscripts will always distinguish these from
expectation operators. Note that λmin is the point beyond which µλ is no
longer a real number.
The martingales are Z−λ and Z
λ , defined for λ ∈ (λmin,0] as
Z±λ (t) :=
v±λ (Yu(t))e
λXu(t)−E±λ t,(6)
where v±λ (y) := exp(ψ
2) are strictly-positive eigenfunctions of the opera-
Qθ + 12λ
2A+R,
with corresponding eigenvalues E−λ <E
λ and A,R are the functions defined
in Section 1.1. This operator is self-adjoint on L2(φ) with the inner product
〈·, ·〉φ where 〈f, g〉φ :=
fgφdy and φ is the standard normal density. Note
that v−λ ∈ L2(φ), whereas v
λ /∈L2(φ) so is not normalizable.
The calculations of Section 3 make it easy to see these are martingales,
and throughout the paper we will need a variety of martingale convergence
results which are gathered together in Section 8. In particular, we will need
to know precisely when Z−λ is uniformly integrable with a strictly positive
limit, some further strong convergence results for other closely related sums
over particles (also identifying which particles contribute nontrivially to their
limits), and the rate of convergence to zero of the Z+λ martingales.
6 Y. GIT, J. W. HARRIS AND S. C. HARRIS
2.2. The asymptotic growth-rate of particles along spatial rays. As an
essential initial step toward determining the growth rate of particles in the
two-dimensional space-type domain, we first look at the growth rate of par-
ticles in the spatial dimension only.
For γ ≥ 0 and C ⊂R, define
Nt(γ;C) :=
1{Xu(t)≤−γt;Yu(t)∈C}.(7)
The limit giving the expected rate of growth,
t−1 logE(Nt(γ;R))
can be shown to exist and its value can be calculated to be
∆(γ) := inf
λ∈(λmin,0)
{E−λ + λγ}
a−1(θ − 8r)(4γ2 + θa).
An outline for this expectation calculation is given in Section 3.
It is now tempting to guess that the asymptotic speed of the spatially
left-most particle, c̃(θ), is given by
c̃(θ) := sup{γ :∆(γ)> 0}
r+ ρ+
2(2r+ ρ)2
θ− 8r
Recall that c̃(θ) = infλ∈(λmin,0) c
λ is also the minimum threshold for traveling
waves. In this particular situation, the guess that “expectation” and “almost
sure” right-most particle speeds agree was first proved rigorously using a
martingale change of measure technique in [13]. In this paper, we extend
this connection and prove that the “expected” and “almost sure” rates of
growth of particles with given speeds (Theorem 1) and given space-type
locations (Theorem 3) agree.
Theorem 1. Let γ ≥ 0 and y0 < y1. Under each P x,y law, the limit
D(γ) := lim
t−1 logNt(γ; [y0, y1])
exists almost surely and is given by
D(γ) =
∆(γ), if 0≤ γ < c̃(θ),
−∞, if γ ≥ c̃(θ).
A TYPED BRANCHING DIFFUSION 7
Note that symmetry in the process means there is a corresponding re-
sult for particles with spatial velocities greater than +γ (corresponding to
positive λ values). We may occasionally make use of such process symme-
tries without further comment. Then, since Nt(γ;R) is integer valued, the
asymptotic speed of the right-most particle follows immediately:
Corollary 2. Almost surely,
t−1 sup{Xu(t) :u ∈Nt}= c̃(θ).
This spatial growth rate result is proved in Section 10 using the martin-
gale results from Section 8. In fact, it is very easy to obtain the upper bound
by first dominating the indicator function with exponentials to reveal that
Nt(γ;R) ≤ exp{(E−λ + λγ)t}Z
λ (t), recalling that Z
λ is a convergent mar-
tingale, and then optimizing over the choice of λ. For the lower bound, we
will use a strong convergence result obtained in [12], combined with the idea
that each uniformly integrable martingale Z−λ essentially “counts” only the
particles of corresponding velocity −γ.
2.3. The asymptotic shape and growth of the branching diffusion. The
main result of this paper is the almost-sure rate of growth of particles which
are in the vicinity of −γt in space and near κ
t in type position at large
times t. For γ, κ≥ 0, it can be shown that the limit
t−1 logE(Nt(γ; [κ
t,∞)))(10)
exists and takes the value
∆(γ,κ) := inf
λ∈(λmin,0)
{E−λ + λγ − κ
2ψ+λ }
(θ− κ2)
θ(θ− 8r)(4aθγ2 + a2(θ+ κ2)2).
An outline of this expectation calculation is given in Section 3. Once again,
we will find that the “almost sure” rate of growth of particles agrees with
this “expected” rate exactly where there is growth in particle numbers.
Theorem 3. Let γ,κ ≥ 0 with ∆(γ,κ) 6= 0. Under each P x,y law, the
limit
D(γ,κ) := lim
t−1 logNt(γ; [κ
t,∞))
exists almost surely and is given by
D(γ,κ) =
∆(γ,κ), if ∆(γ,κ)≥ 0,
−∞, if ∆(γ,κ)< 0.(12)
8 Y. GIT, J. W. HARRIS AND S. C. HARRIS
To prove the tricky lower bound of Theorem 3, which amounts to the
major work of this paper, we will exhibit an explicit two-phase mechanism by
which the branching diffusion can build up at least the required exponential
number of particles near to −γt in space and κ
t in type position by large
times t.
During the first phase, over a large time t the process builds up an ini-
tial excess of approximately exp(∆(α)t) particles with spatial position at
least −αt, as is already known from Theorem 1. In this “ergodic” phase,
“typical” particles found near −αt in space will have drifted with a steady
spatial speed of α whilst their type histories will have behaved roughly like
OU processes with inward drift of µλy for a certain optimal choice λ(α) of
parameter λ.
For the second phase, we will show that the probability any individual
particle has at least one descendant that makes a “rapid ascent” in both
space and type dimensions from initial position (0,0) to final position near
(−βt,κ
t) is approximately exp(−Θ(β,κ)t), where
Θ(β,κ) =
θ(θ− 8r)(a2κ4 +4aθβ2)
,(13)
and the time taken for this “rapid ascent” is an interval [0, τ ]. We show that
this time τ can be chosen such that 2µλτ ∼ log t, and hence the additional
time is asymptotically negligible in comparison with t. Intuitively, we will
see that given an offspring that has successfully made such a difficult “rapid
ascent,” it will most likely have roughly had its type process behaving like an
OU process with an outward drift of µλy and the Brownian motion driving
its spatial motion will have had a drift λ [corresponding to a real time spatial
drift λA(y) that increases in strength as the type position y increases], for
some optimal choice λ(β,κ) of parameter λ. The precise result required will
be formulated rigorously as a large-deviation lower bound in Theorem 7
of Section 5, and is proved using a “spine” change of measure technique
intimately related to the Z+λ martingales.
Combining these two phases and using independence of the particles, we
can see that the number of particles near (−αt,0) at time t that subse-
quently proceed to have at least one descendant near (−(α + β)t, κ
t) is
approximately Poisson with mean
exp({∆(α)−Θ(β,κ)}t).
Optimizing for a fixed overall spatial speed γ, some calculus reveals that
α+β=γ
α,β>0
{∆(α)−Θ(β,κ)}=∆(ᾱ)−Θ(β̄, κ) = ∆(γ,κ),(14)
A TYPED BRANCHING DIFFUSION 9
with optimal parameters
ᾱ= γ
θ+ κ2
and β̄ = γ
θ+ κ2
.(15)
Thus we will be able to demonstrate an explicit two-phase mechanism pro-
ducing the required number of particles, with this outline argument later
guiding our rigorous proof. In addition, it is interesting to note that the
optimal choices for λ over each phase then also coincide at a single value
λ̄= λ(ᾱ) = λ(β̄, κ).
An informative large deviation heuristic for the rapid ascent can also be
found in Section 4, with this section also containing some essential optimal
path calculations. We actually prove the two-phase mechanism for the lower
bound of Theorem 3 in Section 5, although we defer proving the large-
deviation lower bound until Section 7 after presenting the necessary “spine”
background in Section 6.
We prove the upper bound for the space-type growth rate in Section 9,
again making crucial use of martingale results from Section 8. Similarly to
the spatial growth case, we can find an upper bound using the Z+λ mar-
tingales, that is to say Nt(γ; [κ
t,∞)) ≤ exp{(E+λ + λγ − κ2ψ
λ )t}Z
λ (t).
However, as each Z+λ martingale converges to zero, we must show that its
exponential decay rate is (E+λ −E
λ ) before being able to optimize over the
choice of λ to obtain the required upper bound.
Given Theorem 3, and noting symmetries, it becomes straightforward to
retrieve the following:
Corollary 4. For any F ⊂R2, define
Nt(F ) :=
1{(Xu(t)/t,Yu(t)/
t)∈F}.
If B ⊂R2 is any open set and C ⊂R2 is any closed set, then almost surely
under any P x,y
lim inf
logNt(B)≥ sup
(γ,κ)∈B
D(γ,κ),
lim sup
logNt(C)≤ sup
(γ,κ)∈C
D(γ,κ),
with the growth rate D(γ,κ) given at equation (12).
We can also recover the almost sure asymptotic shape of the region occu-
pied by the particles in the branching diffusion.
10 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Corollary 5. Let B ⊂R2 be any open set. Almost surely, under each
P x,y law,
Nt(B)→
0, if S ∩B =∅,
+∞, if S ∩B 6=∅,
where S ⊂R2 is the set given by
S := {(γ,κ) ∈R2|∆(γ,κ)> 0}.
3. Some expectation calculations. This section discusses how the ex-
pected growth rates given in the previous section may be obtained. For this,
we use the “many-to-one” lemma (see, e.g., [8]) and one-particle changes of
measure. In the process we shall start to gain valuable intuition into how
particles within the branching diffusion behave, as well as seeing hints as to
which are the “correct” martingales to use to prove the almost-sure growth
rate results.
For simplicity, we assume throughout this section that the branching dif-
fusion starts with one particle at the origin in both space and type at time
zero, unless otherwise stated. We also introduce a family of single particle
probability measures Pµ,λ with associated expectations Eµ,λ where, under
Pµ,λ, η is an Ornstein–Uhlenbeck process with variance θ and drift µ, and
ξt =B(
0 A(ηs)ds) where B is a Brownian motion with drift λ.
Lemma 6 (Many-to-one). If f :R2 7→R is Borel measurable then
f(Xu(t), Yu(t)) = Eθ/2,0
R(ηs)ds
f(ξt, ηt)
.(16)
Using the many-to-one lemma, and changing measure to alter the drift of
Brownian motion, we see that
f(Xu(t), Yu(t))
= Eθ/2,0
R(ηs)ds
f(ξt, ηt)
= Eθ/2,0
e−λξt exp
R(ηs) +
A(ηs)
× f(ξt, ηt) · eλξt−λ
A(ηs)ds
= Eθ/2,λ
−λξt +
R(ηs) +
A(ηs)
f(ξt, ηt)
A TYPED BRANCHING DIFFUSION 11
To perform a further change of measure on the OU process to get rid of the
time integrals in the exponential of the expectation, we recall that
dPµλ,·
dPθ/2,·
µλ,θ/2
:= exp
ψ−λ η
t −E−λ t+
R(ηs) +
λ2A(ηs)
and then
f(Xu(t), Yu(t))
= Eθ/2,λ(exp(−λξt −ψ−λ η
λ t)f(ξt, ηt) ·M
µλ,θ/2
t )(17)
= Eµλ,λ(exp(−λξt − ψ
λ t)f(ξt, ηt)).
Note that the many-to-one lemma, combined with the branching prop-
erty, immediately suggests how to get “additive” martingales for the branch-
ing diffusion from single particle martingales—for example, taking f(x, y) =
exp{λx+ ψ−λ y2} in equation (17) quickly leads to the martingale Z
We may now proceed to calculate the expected growth rates. However,
for both clarity and brevity we will leave rigorous details to the interested
reader, noting that the intuition we will gain from our rough calculations
will later be invaluable in guiding our rigorous proof of the corresponding
almost-sure growth rates.
3.1. The expected rate of growth along spatial rays. We first give the
outline of some calculations to find the rate of growth in the expected number
of particles near −γt in space at time t.
Using the formula from (17), for λ ∈ (λmin,0) and any ε > 0 we have
1{t−1Xu(t)+γ∈(−ε,ε)}
= Eµλ,λ(e
−λξt−ψ−λ η
1{t−1ξt+γ∈(−ε,ε)})
≤ e(E
+λγ−λε)tEµλ,λ
+ γ ∈ (−ε, ε)
≥ e(E
+λγ+λε)tEµλ,λ
η2t ;
+ γ ∈ (−ε, ε)
where, with some abuse of notation that we shall continue to use throughout
this section, we will abbreviate this to
1{Xu(t)∼−γt} = Eµλ,λ(e
−λξt−ψ−λ η
1{ξt∼−γt})
∼ e(E
+λγ)t
Eµλ,λ(e
t ; ξt ∼−γt)
12 Y. GIT, J. W. HARRIS AND S. C. HARRIS
with the understanding that any subsequent arguments to identify expo-
nential growth rates can readily be made rigorous by using the appropriate
upper and lower bounds, and so on.
Now, considering E−(λ) :=E−λ as a function of λ, we have from (8) that
∆(γ) = infλ∈(λmin,0){E−(λ) + λγ}=E−(λγ) + λγγ, where λγ satisfies
(λγ) =−γ, hence λγ =−γ
(θ− 8r)
θa2 +4aγ2
.(19)
Of course, choosing this optimal λγ value in (18) means that we must have si-
multaneously maximized the expectation Eµλ,λ(exp(−ψ
t ); ξt ∼−γt), and
to confirm that this value is not exponentially decaying in t is now relatively
straightforward. Under Pµλ,λ, η is an Ornstein–Uhlenbeck process with an
invariant measure given by the probability density, φλ, of the normal distri-
bution N(0, θ/(2µλ)); and ξt =B(
0 A(ηs)ds), where B is a BM with drift
λ. Note also that by differentiating
〈(Qθ + (1/2)λ2A+R−E−λ )v
λ , v
λ 〉φ = 0
with respect to λ, using self-adjointness, and observing that φλ ∝ (v−λ )2φ,
we find that
〈λAv−λ , v
〈v−λ , v
A(y)φλ(y)dy.
Then almost surely under Pµλ,λ,
0 A(ηs)ds)
0 A(ηs)ds
0 A(ηs)ds
Aφλ dy =
,(20)
and so when we use the optimal λγ value we get exactly the desired drift,
since ∂E
(λγ) =−γ. Then
Eµλγ ,λγ (e
η2t ; ξt ∼−γt) → lim
Eµλγ ,λγ(e
η2t )
φλγ (y)dy.
In this way, we can obtain the exact rate of exponential growth for the
expectation,
t−1 logE(Nt(γ;R)) = ∆(γ).
The changes of measure used above are actually suggesting a great deal
about the dominant particles that are found in the vicinity of a given ray
in space. An alternative discussion of this expectation result, involving a
dual approach via large deviation theory for occupation densities, can also
be found in [13].
A TYPED BRANCHING DIFFUSION 13
3.2. The expected asymptotic shape. We give a rough outline of calcula-
tions that will yield the correct exponential growth in the expected number
of particles both near −γt in space and κ
t in type at large times t. Using
the formula from (17) and abusing notation throughout in the same way as
Section 3.1, we find that
1{Xu(t)∼−γt;Yu(t)≥κ
= Eµλ,λ(e
−λξt−ψ−λ η
1{ξt∼−γt;ηt≥κ
∼ e(E
+λγ−κ2ψ−
Pµλ,λ(ξt ∼−γt;ηt ≥ κ
Now, from standard bounds on the tail of the normal distribution,
Pµλ,λ(ξt ∼−γt;ηt ≥ κ
= Pµλ,λ(ηt ≥ κ
t)Pµλ,λ(ξt ∼−γt|ηt ≥ κ
t)(21)
∼ e−µλ/θκ2tPµλ,λ(ξt ∼−γt|ηt ≥ κ
and, since ψ−λ + (µλ/θ) = ψ
λ , this yields
1{Xu(t)∼−γt;Yu(t)≥κ
∼ e(E
+λγ−κ2ψ+
Pµλ,λ(ξt ∼−γt|ηt ≥ κ
Recalling that ∆(γ,κ) := infλ∈(λmin,0){E
λ +λγ−κ2ψ
λ }, simple calculus re-
veals this infimum is attained at a λ value of
λ̄(γ,κ) =−γ
θ(θ− 8r)
a2(κ2 + θ)2 + 4aγ2θ
∈ (λmin,0),(23)
and using this optimal value in equation (22) will lead to the upper bound
limsup
t−1 logE
1{Xu(t)∼−γt;Yu(t)≥κ
t} ≤∆(γ,κ).
It is also clear from equation (22) that when minimizing E−λ + λγ − κ2ψ
we simultaneously maximize the probability Pµλ,λ(ξt ∼ −γt |ηt ≥ κ
t). In
particular, to get a matching lower bound, we do not want this probability
to have any exponential decay in time when we choose the optimal parameter
for λ.
In fact, at least up to the exponential decay rate in time, it can be shown
using large-deviations arguments that
Pµλ̄,λ̄
(ξt ∼−γt;ηt ≥ κ
t)∼ exp
14 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Indeed, we immediately gain the required upper bound from (21). For the
lower bound, consider the following heuristics where we break paths into two
sections: normal ergodic behavior over large time period [0, t] followed by a
rapid ascent out to type position κ
t over a much shorter period [t, t+ τ ].
(i) Ergodic behavior. Over a large time t, the occupation density of η
will most likely have settled close to the invariant measure. Hence for large
t, almost surely under Pλ,µλ ,
η2s ds→
(ii) Rapid ascent. Over a large time τ , but where τ = o(t), the probability
that η starts close to the origin and ends near to κ
t, having followed close
to the path y over the entire time period τ , is roughly given by
{ẏ(s) + µλy(s)}2 ds
under the Pλ,µλ law. See, for example [24], Chapter 6, or [4], Chapter 5.6.
After some Euler–Lagrange optimization, the path
y(s) = κ
sinhµλs
sinhµλτ
gives
0 y(s)
2 ds≈ κ2t/(2µλ), with the probability of this path being roughly
exp(−(µλ/θ)κ2t).
Combining these two types of behavior, we can find paths with final positions
ηt+τ ∼ κ
t, ξt+τ ∼ λa
∫ t+τ
η2s ds∼ λa
and, moreover, when substituting the optimal λ value of λ̄(γ,κ) and sim-
plifying, this actually gives ξt+τ ∼ −γt. Further, one of these paths occurs
with a probability of roughly exp(−(µλ̄/θ)κ2t), and note that t+ τ ∼ t since
τ = o(t). Thus we see that to exponential order, the probability Pµλ̄,λ̄
(ξt ∼
−γt;ηt ≥ κ
t) must be at least exp(−(µλ̄/θ)κ2t), as required.
This heuristic argument can be made rigorous to prove, as claimed, that
t−1 logE
1{Xu(t)≤−γt;Yu(t)≥κ
=∆(γ,κ).
If we scale all spatial coordinates by t−1 and all type coordinates by
t)−1 at time t, the expected asymptotic shape can be considered to be the
region S := {(γ,κ) :∆(γ,κ) ≥ 0} where, on average, we have growth in the
numbers of (scaled) particles.
A TYPED BRANCHING DIFFUSION 15
4. Short climb large deviation heuristics. In this section, we give a heuris-
tic calculation that suggests why the probability a single particle manages
to have at least one descendant in the vicinity of (−βt,κ
t ) near time τ
is roughly exp(−Θ(β,κ)t) for very large t, where Θ(β,κ) is given at equa-
tion (13). For these heuristics, we will think of τ as large and fixed, but of
smaller order than t (later on, in our rigorous approach, we will choose τ
proportional to log t). We emphasize that the heuristics in this section are
neither meant to be precise nor made rigorous, yet they will provide invalu-
able intuition, guidance and motivation for our rigorous approach later on.
Of particular importance will be the optimization problem that the heuris-
tics suggest. Indeed, many of the exact calculations in Sections 4.2 and 4.3
will be essential later in the paper.
Suppose we start the branching diffusion with a single particle at (0,0).
First, we wish to know the probability that there is at least one particle at
time τ that has a spatial position near −βt having followed close to the path
x(s) for 0≤ s≤ τ and a type position near κ
t having closely followed the
path y(s) for 0≤ s≤ τ for t arbitrarily large.
We recall from large deviation theory of Ventcel–Freidlin (see [24], Chap-
ter 6, or [4], Chapter 5.6) that the probability a single particle manages to
follow closely both the type path y(s) and the spatial path x(s) for 0≤ s≤ τ
is roughly given by
ẏ(s) +
ds− 1
ẋ(s)2
ay(s)2
when x(0) = 0, x(τ) = −βt, y(0) = 0, y(τ) = κ
t and t is very large. This
probability will typically be very small, but if such paths are followed by
particles in the branching diffusion, we have to also take account of the
large breeding rates that are found far from the type origin.
If we let X(s) represent the numbers of particles in the branching diffusion
that are alive at time s and have traveled “close” to the path (x(u), y(u))
for 0 ≤ u ≤ s, then we can get a rough idea of how X might behave by
considering the following birth–death process.
4.1. A birth–death process. For given fixed paths x(·) and y(·), let M be
a time-dependent birth–death process where at time s particles either give
birth to single offspring with breeding rate λ(s) given by
λ(s) = ρ+ ry(s)2,
or particles die with death rate µ(s) given by
µ(s) =
ẏ(s) +
ẋ(s)2
ay(s)2
16 Y. GIT, J. W. HARRIS AND S. C. HARRIS
(Note that the probability the initial particle of this birth–death process sur-
vives the entire time period [0, τ ] is consistent with the rough large deviation
probability for the branching diffusion at equation (24).)
An important quantity is the effective total death rate up to time t which
is defined by ν(s) :=
0 {µ(w)− λ(w)}dw, so here
ν(s) = J(x, y, s)
ẏ(w) +
ẋ(w)2
ay(w)2
− ry(w)2 − ρ
The distribution for total number of offspring surviving, M(τ), for the
time-dependent birth–death process is well known, for example, see [14].
Then defining
Wτ := e
−ν(τ)
µ(s) eν(s) ds
Uτ := 1− e−ν(τ)W−1τ ,
Vτ := 1−W−1τ ,
we have
P(M(τ) = 0) = Uτ ,
P(M(τ) = n) = (1−Uτ )(1− Vτ )V n−1τ , n= 1,2, . . .
with EM(τ) = e−ν(τ) and E(M(τ)|M(τ)≥ 1) =Wτ .
In our particular case, we have
E(M(τ)) = exp(−J(x, y, τ)).
Define the largest effective total death rate prior to time τ by
L(x, y, τ) := sup
s∈[0,τ ]
J(x, y, s)≥ 0.
If we are in a case when L(x, y, τ) is very large, suggesting a high chance of
extinction, then
P(M(τ)≥ 1) = 1
0 µ(s)e
ν(s) ds
∼Kτ exp(−L(x, y, τ)),(25)
where K−1τ :=
0 µ(s) exp(−{L(x, y, τ) − J(x, y, s)})ds. If there is at least
one particle alive, we would then expect to have
E(M(τ)|M(τ)≥ 1)∼K−1τ exp(L(x, y, τ)− J(x, y, τ)).
Thus, we might guess that the probability any particles in the branching
diffusion manage to make the difficult, rapid ascent along path (x, y) to finish
up near (−βt,κ
t ) can, very roughly, be estimated by exp(−L(x, y, τ)). [To
A TYPED BRANCHING DIFFUSION 17
help see this, try writing x(s) = tf(s) and y(s) =
tg(s), thinking of f, g as
fixed paths and recall that t is very large and τ = o(t), then the role of Kτ
in (25) is insignificant next to exp(−L(x, y, τ)).]
We might then further guess that the chance any particles manage to
stay near position (−βt,κ
t) during a very small interval of time close to τ
should roughly look like
− inf
L(x, y, τ)
where we permit all possible paths x and y satisfying x(0) = 0, x(τ) =−βt
and y(0) = 0, y(τ) = κ
t for the fixed time τ . (We will state and prove a
precise lower bound that corresponds to this guess at Theorem 7.)
4.2. Finding the optimal path and probability. We proceed to calculate
L(x, y, τ)
over paths x and y satisfying x(0) = 0, x(τ) =−βt and y(0) = 0, y(τ) = κ
for the fixed time τ .
We first note that
L(x, y, τ) = inf
s∈[0,τ ]
J(x, y, s)≥ inf
J(x, y, τ)(26)
and we now proceed to calculate infx,y J(x, y, τ).
We can easily optimize over the choice of function x given y, finding that
ẋ(s)∝ ay(s)2 ⇒ x(s) = λa
y(u)2 du
where λ is the constant of proportionality and must satisfy
0 y(s)
,(27)
yielding
ẋ(s)2
ay(s)2
0 y(s)
This is exactly as anticipated since, when following the path y in type space,
the spatial position of a particle is following a Brownian motion with total
amount of variance over period τ given by a
0 y(s)
2 ds. Hence, the proba-
bility that a particle following the path y in type space will also be found
near to βt in space at time τ is roughly
( −β2t2
0 y(s)
18 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Introducing the notation
I(y) :=
ẏ(s) +
− ry(s)2
we are left to find
I(y) +
0 y(s)
= inf
I(y)− 1
y(s)2 ds− λβt
≥ sup
I(y)− 1
y(s)2 ds− λβt
where the first equality is trivially true by maximizing the quadratic in
λ, the introduction of which conveniently removes the awkward integral in
the denominator. Some further Euler–Lagrange optimization now gives the
optimal path as
yλ(s) = κ
sinhµλs
sinhµλτ
(0≤ s≤ τ),(29)
where
θ(θ− 8r− 4aλ2)
and then
I(y)− 1
y(s)2 ds− λβt
= sup
cothµλτ
− λβt
The optimal parameter choice λ̂ (which depends on τ as well as the model
parameters) then satisfies
= κ2t
cothµ
2 sinh2 µ
(s)2 ds.(30)
Then we have shown that
0 yλ̂(s)
≥ inf
I(y) +
0 y(s)
= inf
I(y)− 1
y(s)2 ds− λβt
≥ sup
I(y)− 1
y(s)2 ds− λβt
≥ I(y
(s)2 ds− λ̂βt,
A TYPED BRANCHING DIFFUSION 19
and, in fact, we see that the left- and right-hand sides of the above are equal
by recalling (30). It follows that the preceding supremum and infimum can
be freely interchanged, actually preserving equality at the inequality (28).
Then, with the optimal spatial path
xλ(s) := λa
yλ(u)
2 du=−βt sinh2µλs− 2µλs
sinh2µλτ − 2µλτ
,(31)
and defining x̂ := x
, ŷ := y
, we have
J(x, y, τ) = J(x̂, ŷ, τ)
= t sup
cothµλτ
cothµ
− λ̂β
− ρτ.
Finally, it is easy to check that J(x̂, ŷ, τ) = L(x̂, ŷ, τ), whence
J(x, y, τ)≥ inf
L(x, y, τ),
and, combining with equation (26), we have found that
L(x, y, τ) = inf
J(x, y, τ) = J(x̂, ŷ, τ).
4.3. An important note on the optimal paths. As τ →∞, we have
cothµλτ
↑ sup
{κ2ψ+λ − λβ}= κ
− λ̄ β,
where the optimizing parameters of the supremums also converge with
λ̂→ λ̄=−β
θ(θ− 8r)
a2κ4 + 4aθβ2
κ2 + θ
.(32)
Note the agreement with previous optimal values at equations (23) and (15).
Then letting
Θ(β,κ) := sup
{κ2ψ+λ − λβ}
θ(θ− 8r)(a2κ4 +4aθβ2)
and writing x̄ := xλ̄ and ȳ := yλ̄, we note that for all ε, δ > 0 there exist
τ̃ , µ > 0 such that for all t > 0 and τ > τ̃
− inf
J(x, y, τ)
≥ exp(−J(x̄, ȳ, τ))
= exp
cothµλ̄τ
− λ̄β
≥ exp(−t(Θ(β,κ) + ε)).
20 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Further (when κ > 0), for all s ∈ [τ − µ, τ ],
ȳ(s)≥ (κ− δ)
t, x̄(s)≤−(β − δ)t.
In particular, the paths stay close to the required positions for some fixed
length of time with corresponding probability at least as large as required.
5. Proof of Theorem 3. Lower bound. In this section we will state a pre-
cise short climb probability result and show how to combine it with almost
sure spatial (only) growth rates to prove the lower bound of the growth rate
in Theorem 3. This will make rigorous the two-phase mechanism described
in Section 2 and suggested by the expectation calculations in Section 3.
The first phase requires knowledge of the almost-sure rates of growth of
particles in the spatial dimension only. To this end, we will already make full
use of Theorem 1 throughout this section, deferring its proof until Section 10.
The second phase requires a lower bound for the probability that a single
particle makes a rapid ascent in type-space over the time interval [0, τ ]. This
is the lower bound found in the heuristics of Section 4, but we require some
further notation before the precise result can be stated. Note, throughout
this section, we will only be interested in the optimal parameter value λ= λ̄
as introduced in Section 4.3.
We wish to fix the relationship between sufficiently large t and τ as
θ/(2µλ̄) e
µλ̄τ = κ
t(34)
and so define τ = τ(t) by
τ(t) :=
(2µλ̄)
−1 log(2µλ̄t/θ), for 2µλ̄t > θ,
0, otherwise.
Recall the optimal paths (x̄, ȳ) over s ∈ [0, τ ], where
ȳ(s) = κ
sinhµλ̄s
sinhµλ̄τ
,(36)
x̄(s) = aλ̄
ȳ(w)2 dw=−βt sinh2µλ̄s− 2µλ̄s
sinh2µλ̄τ − 2µλ̄τ
,(37)
with fixed end points ȳ(τ) = κ
t and x̄(τ) =−βt.
For large times t and δ, ε > 0, let
t (u) :=
s∈[0,τ(t)]
|Yu(s)− ȳ(s)|< ε
t; sup
s∈[0,τ(t)]
|Xu(s)− x̄(s)|< δt
.(38)
We will use the notation
u∈Nτ(t)
t (u)(39)
A TYPED BRANCHING DIFFUSION 21
for the event that there exists a particle in the branching diffusion that
makes the short climb. Finally, recalling Θ(β,κ) given at (33), we can now
state the short climb theorem:
Theorem 7. Fix any y1 > y0 > 0, x ∈ R, and let ε0 > 0. Then for any
ε, δ > 0, there exists T > 0 such that for all y ∈ [y0, y1],
t−1 logP x,y(Aε,δt )≥−(Θ(β,κ) + ε0)
for all t > T .
We will prove Theorem 7 using a spine change of measure. This requires
us to introduce the notation for the spine set-up in detail before proceeding,
so this and further technical issues are postponed to Sections 6 and 7.
Remark 8. We note that Theorem 7 is actually a stronger result than
needed to prove Theorem 3 because we identify the specific paths followed
by particles that are near position (βt,κ
t ) at time t+ τ , rather than just
considering the particle’s positions close to time t+ τ .
In combining the two phases, we will have a huge number of independent
trials each with a small probability of success, intuitively giving rise to a
Poisson approximation for a large number of successful particles. In fact,
in our proof of the lower bound of Theorem 3 below, we will actually use
the following result about the behavior of sequences of sums of independent
Bernoulli random variables.
Lemma 9. For each n, define the random variable Bn :=
u∈Fn 1En(u)
where the events {En(u) :u ∈ Fn} are independent. Let pn(u) := P (En(u))
and Sn :=
u∈Fn pn(u) and suppose that, for some ν ∈ (1/2,1),
(Sn)2ν−1
<∞.(40)
Then the sequence of (possibly dependent) random variables {B1,B2, . . .}
has |Bn − Sn|> (Sn)ν for only finitely many n, almost surely.
In particular, for any ε > 0, there exists some (random) N ∈N such that,
with probability one,
> 1− ε for all n >N.(41)
Proof. For ν ∈ (1/2,1), Chebyshev’s inequality yields
P(|Bn − Sn|> (Sn)ν)≤
u∈Fn pn(u)(1− pn(u))
2ν−1 ,
22 Y. GIT, J. W. HARRIS AND S. C. HARRIS
and hence the Borel–Cantelli lemmas, combined with hypothesis (40), imply
|Bn − Sn|> (Sn)ν
for only finitely many n, almost surely. Equation (41) now follows on division
by Sn, and noticing the assumption (40) implies that limn→∞Sn =∞. �
Proof of Theorem 3. Lower bound. Define f−1(t) := t− τ(t), not-
ing that both f(t)/t→ 1 and f−1(t)/t→ 1 as t→∞. Also, for n ∈ N and
µ > 0, define Tn := (n+ 1)µ. We want to estimate the number of particles
that are near the large position (−(α + β)Tn, κ
Tn) during time interval
[Tn−1, Tn]. For this, we will consider particles that travel with a velocity
−α over time period [0, f−1(Tn)] before commencing their rapid ascent of
(relatively short) duration τ(Tn) to be in final position at time Tn. Then
s∈[Tn−1,Tn]
Ns((α+ β − δ)Tn; [(κ− δ)
Tn,∞))
u∈NTn
s∈[Tn−1,Tn]
{Xu(s)≤−(α+β−δ)Tn ;Yu(s)≥(κ−δ)
Tn}}(42)
u∈Fαn
1{N̄β,κn (u)>0}
where
Fαn := {u ∈Nf−1(Tn) :Xu(f
−1(Tn))≤−αTn, Yu(f−1(Tn)) ∈ [y0, y1]}
and, for u ∈ Fαn ,
N̄β,κn (u) :=
v∈NTn
s∈[Tn−1,Tn]
{Xv(s)−Xv(f−1(Tn))≤−(β−δ)Tn ;Yv(s)≥(κ−δ)
Tn}}.
We will now show that the sum at (42) grows as fast as anticipated:
Lemma 10. For any ε > 0, we may choose µ > 0 such that there exists
a random N ∈N where
u∈Fαn
1{N̄β,κn (u)>0}
≥∆(α)−Θ(β,κ)− ε
for all n >N with probability one.
Proof. We will be able to apply Lemma 9 given sufficient information
about the growth of |Fαn | and decay of the probabilities
pβ,κn (u) := P (N̄
n (u)> 0|Ff−1(Tn)),
A TYPED BRANCHING DIFFUSION 23
where u ∈ Fαn ⊂Nf−1(Tn).
It follows easily from Theorem 1, f−1(Tn)/Tn → 1 and the continuity of
∆(α) that
log |Fαn |
≥∆(α)− ε
for all sufficiently large n.
The definition of N̄β,κn (u) and spatial translation invariance implies that,
for each u ∈ Fαn , the rapid ascent probability pβ,κn (u) depends only on the
initial type position Yu(f
−1(Tn)).
For δ,µ > 0, define
t (u) :=
s∈[τ(t)−µ,τ(t)]
{Xu(s)−Xu(0)<−(β − δ)t;Yu(s)≥ (κ− δ)
u∈Nτ(t)
t (u).(43)
Recalling the comments of Section 4.3, there exist ε′, δ′ > 0 and we may
choose µ> 0 sufficiently small, such that
pβ,κn (u) = P
0,Yu(f
−1(Tn))(B
)≥ P 0,Yu(f−1(Tn))(Aε
) =: p̄n(u)
for all u ∈ Fαn whenever n is sufficiently large. Together with Theorem 7 and
since Yu(f
−1(Tn)) ∈ [y0, y1] for u ∈ Fαn , this reveals
log pβ,κn (u)
≥ log p̄n(u)
≥−Θ(β,κ)− ε
for all for u ∈ Fαn and all sufficiently large n, almost surely. Then we may
combine the observations above to obtain
u∈Fαn
pβ,κn (u)≥∆(α)−Θ(β,κ)−
Taking this last line together the assertion of Lemma 9 at equation (41)
gives the result. �
It is now straightforward to combine Lemma 10 with the inequality at
(42) to see that, given ε, δ > 0, there exists µ > 0 and a random time T such
t−1logNt((α+ β − δ)t; [(κ− δ)
t,∞))≥∆(α)−Θ(β,κ)− ε
for all t > T , almost surely. Since ε and δ can be taken arbitrarily small,
using the optimal ᾱ and β̄ according to equations (14)–(15), we find
lim inf
t−1 logNt(γ, [κ
t,∞))≥∆(γ,κ) almost surely,
24 Y. GIT, J. W. HARRIS AND S. C. HARRIS
as required. [It is also interesting to note that λ̄= λᾱ = λ̄(γ,κ) from equa-
tions (19), (23) and (32), so the optimal parameters are in agreement with
those of the expectation calculations in Section 3 and the path large devia-
tions in Section 4.] �
6. The “spine” setup and results. In this section, we describe how to
construct an enriched branching diffusion with an identified “spine” or “back-
bone” particle and discuss how to perform some extremely useful changes
of measure (closely related to the additive martingales) that will essentially
“force” the spine perform the short climb, whilst giving birth at an accel-
erated rate to offspring that behave as if under the original measure. These
spine techniques are at the very heart of our proof of Theorem 7 in Sec-
tion 5. Spine ideas were first seen for branching Brownian motion in [3] and
developed for Galton–Watson processes in [16, 18, 19]. Kyprianou [17] and
Englander and Kyprianou [6], developed the technique for some families of
branching diffusions; and more recently the spine approach has been signif-
icantly improved in [8]. This approach uses several different filtrations on
an enlarged probability space carrying the branching diffusion, and permits
some very useful techniques and results to be developed. For example, “addi-
tive” (many-particle) martingales can be represented as suitable conditional
expectations of “spine” (single-particle) martingales and consequently there
are clear interpretations for any changes of measure and all measures in-
volved in our “spine” setup are probability measures with intuitive construc-
tions. Following Hardy and Harris [8], we will first outline the notation and
then describe the changes of measure. The notation described in this section
is generalized to allow each particle u to have 1 +Au offspring, where each
Au is an independent copy of a random variable with values in {0,1,2, . . .}.
The spine techniques developed in this paper could readily be generalized
to such models.
All probability measures are to be defined on the space T̃ of marked
Galton–Watson trees with spines; before defining precisely what this space
is we need to set up some other notation. We recall the set of Ulam–Harris
labels, Ω, defined by Ω := {∅}∪
n∈N(N)
n, where N := {1,2,3, . . .}. For two
words u, v ∈Ω, uv denotes the concatenated word, where we take u∅=∅u=
u. So Ω contains elements such as “∅412,” which represents “the individual
being the 2nd child of the 1st child of the 4th child of the initial ancestor
∅.” For labels u, v ∈Ω the notation v < u means that v is an ancestor of u,
and |u| denotes the length of u.
We define a Galton–Watson tree to be a set τ ⊂Ω such that:
(i) ∅ ∈ τ , so there is the unique initial ancestor;
(ii) if u, v ∈Ω, then vu ∈ τ ⇒ v ∈ τ , so τ contains all of the ancestors of
its nodes;
A TYPED BRANCHING DIFFUSION 25
(iii) for all u ∈ τ , there exists Au ∈ {0,1,2, . . .} such that for j ∈N, uj ∈ τ
if and only if 1≤ j ≤ 1 +Au.
The set of all such trees is T, and we will use the symbol τ for a particular
tree. As our work concerns branching diffusions we shall often refer to the
labels of τ as particles. Note that for the binary branching mechanism in this
paper, P (Au = 1) ≡ 1; of course, here there is only one τ ∈ T—the binary
tree.
A Galton–Watson tree by itself only records the family structure of the
individuals, so to each individual u ∈ τ we give a mark (Xu, Yu, σu) which
contains the following information:
• σu ∈ [0,∞) is the lifetime of particle u, which also determines the fission
time of the particle as Su :=
v≤u σv . We may also refer to the Su as
death times;
• the function Xu(t) : [Su − σu, Su)→R describes the particle’s spatial mo-
tion in R during its lifetime;
• the function Yu(t) : [Su− σu, Su)→R describes the evolution of the parti-
cle’s type in R during its lifetime.
For clarity we must decide whether or not a particle is in existence at its
death time: our convention will be that a particle dies “infinitesimally be-
fore” its death time—this is why Xu and Yu are defined on [Su−σu, Su) and
not [Su−σu, Su]—so that at time Su the particle u has disappeared and has
been replaced by its two children.
We denote a particular marked tree by (τ,X,Y,σ), or the abbreviation
(τ,M), and the set of all marked Galton–Watson trees by T . For each
(τ,X,Y,σ) ∈ T , the set of particles alive at time t is defined as Nt := {u ∈
τ :Su − σu ≤ t < Su}. For any given marked tree (τ,M) ∈ T we can distin-
guish individual lines of descent from the initial ancestor: ∅, u1, u2, u3, . . . ∈
τ , where ui is a child of ui−1 for all i ∈ {2,3, . . .} and u1 is a child of the
initial individual ∅. We call such a line of descent a spine and denote it by
ξ. In a slight abuse of notation we refer to ξt as the unique node in ξ that is
alive at time t, and also for the position of the particle that makes up the
spine at time t; that is, ξt :=Xu(t), where u ∈ ξ∩Nt. However, although the
interpretation of ξt should always be clear from the context, we introduce
the following notation for use where some ambiguity may still arise:
• nodet((τ,M, ξ)) := u if u ∈ ξ is the node in the spine alive at time t.
It is natural to think of the spine as a single diffusing particle ξt, or, strictly
speaking, the pair (ξt, ηt), where ηt is the type of the spine at time t.
We define nt to be a counting function that tells us which generation of
the spine is currently alive, or equivalently the number of fission times there
have been on the spine:
nt = |nodet(ξ)|.
26 Y. GIT, J. W. HARRIS AND S. C. HARRIS
The collection of all marked trees with a distinguished spine is the space
T̃ on which our probability measures will eventually be defined, but first we
define four filtrations on this space that contain different levels of information
about the branching diffusion.
• Filtration (Ft)t≥0. We define a filtration of T̃ made up of the σ-algebras
Ft := σ((u,Xu, Yu, σu) :Su ≤ t;
(u,Xu(s), Yu(s) : s ∈ [Su − σu, t]) : t ∈ [Su − σu, Su)),
which means that Ft is generated by the information concerning all par-
ticles that have lived and died before time t, and also those that are
still alive at time t. Each of these σ-algebras is a subset of the limit
F∞ := σ(
t≥0Ft).
• Filtration (F̃t)t≥0. We define the filtration (F̃t)t≥0 by augmenting the
filtration Ft with the knowledge of which node is the spine at time t; that
is, (F̃t)t≥0 := σ(Ft,nodet(ξ)) and F̃∞ := σ(
t≥0 F̃t), so that this filtration
knows everything about the branching diffusion and everything about the
spine.
• Filtration (Gt)t≥0. (Gt)t≥0 is a filtration of T̃ defined by Gt := σ(ξs : 0 ≤
s≤ t), and G∞ := σ(
t≥0 Gt). These σ-algebras are generated only by the
spine’s motion and so do not contain the information about which nodes
of the tree τ make up the spine.
• Filtration (G̃t)t≥0. As we did in going from Ft to F̃t we create (G̃t)t≥0
from (Gt)t≥0 by including knowledge of which nodes make up the spine:
(G̃t)t≥0 := σ(Gt,nodet(ξ)) and G̃∞ := σ(
t≥0 G̃t). This means that G̃t also
knows when the fission times on the spine occurred, whereas Gt does not.
Now that we have defined the underlying space and filtrations, we can de-
fine the probability measures of interest. We let the typed branching diffusion
be as described in Section 1.1, with the probability measures {P x,y :x, y ∈R}
on (T̃ ,F∞) representing the law of this typed branching diffusion when ini-
tially started with a single particle at (x, y).
We recall from [18] that, if f is an F̃t-measurable function, we can write
fu1{ξt=u},(44)
where fu is Ft-measurable. Now we can extend P x,y to a measure P̃ x,y on
(T̃ , F̃∞) by choosing the particle that continues the spine uniformly each
time there is a birth on the spine; more precisely, for any f ∈ mF̃t with
representation like (44), we have:
f dP̃ x,y(τ,M, ξ) :=
dP x,y(τ,M).
A TYPED BRANCHING DIFFUSION 27
We construct the F̃t-measurable martingale ζ̃(t) as
ζ̃(t) := v+λ (ηt)e
{R(ηs)+1/2λ2A(ηs)}ds−E+λ t × 2nte−
R(ηs)ds
× eλξt−1/2λ
A(ηs)ds(45)
= v+λ (ηt)2
nteλξt−E
Observe that this is a product of single-particle martingales, details of which
can be found in [17] or [10]. One can think of these as h-transforms of the P̃ -
law of the spine: the first makes η an outward-drifting Ornstein–Uhlenbeck
process with drift parameter µλ; the second increases the breeding rate on
the spine to 2R(·); and the third adds a spatial drift to ξ.
Using the martingale ζ̃(t) we may define a measure Q̃
λ on (T̃ , F̃∞) by
dP̃ x,y
ζ̃(t)
ζ̃(0)
v+λ (y)
v+λ (ηt)2
nteλξt−E
t.(46)
And since ζ̃(t) is a product of h-transforms, under Q̃
λ the process may be
re-constructed path-wise according to the following description:
• starting from spatial position x and type y the spine (ξt, ηt) diffuses spa-
tially as a Brownian motion with infinitesimal variance A(ηt) and infinites-
imal drift λA(ηt);
• the type of the spine, ηt, begins at y and moves in type space as an
outward-drifting Ornstein–Uhlenbeck process with generator
+ µλy
• the spine branches at rate 2R(ηt), producing 2 particles;
• one of these particles is selected uniformly at random;
• the chosen offspring repeats stochastically the behavior of its parent;
• the other offspring particle initiates a P ·,·-BBM from its birth position
and type.
The change of measure (46) projects onto the sub-algebra Ft as a condi-
tional expectation:
dP̃ x,y
v+λ (y)
P̃ x,y(v+λ (ηt)2
nteλξt−E
t|Ft),
and it is a short calculation using the methods of, for example, Hardy and
Harris [10] to show that:
28 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Theorem 11. If we define Q
λ := Q̃
λ |F∞ , then Q
λ is a measure on
F∞ that satisfies
dP x,y
= Ẑ+λ (t) :=
Z+λ (t)
Z+λ (0)
Moreover under Q
λ , the path-wise construction of the branching diffusion
is the same as under Q̃λ.
Although the path-wise construction of the branching diffusion is the same
under Q
λ and Q̃
λ , only the measure Q̃
λ “knows” about the spine. It is
clear, however, that we have Q̃
λ (A) =Q
λ (A) for any A ∈ F∞.
Under the measure Q̃
λ only the behavior of the spine is altered, and
combining this observation with conditioning on the spine’s path and fission-
times gives us a very useful representation for Z+λ (t) under Q̃
λ that we shall
refer to as the spine decomposition:
λ (t)|G̃∞) =
v+λ (ηSu)e
λξSu−E
Su + v+λ (ηt)e
λξt−E+λ t.(47)
Throughout the rest of this article we will refer to the two pieces of this de-
composition as the “sum term” and the “spine term.” This decomposition
is discussed in detail for a wide variety of branching diffusions in [9], but
to derive it we simply note that the contributions to Z+λ (t) from the sub-
trees that branch off the spine have constant Q̃
λ -expectation because they
behave as if under the original measure P , and we know that Z+λ (t) is a P -
martingale. The spine decomposition reduces many calculations about the
behavior of Z+λ (t) under Q̃
λ to one-particle calculations about the spine,
and this observation is exploited in the spine proofs of Lp-bounds for some
families of additive martingales in [9].
7. Proof of Theorem 7. The short climb probability. With the spine
foundations firmly established in Section 6, we may proceed with the proof
of the short climb probability lower bound from Theorem 7.
First, recall definitions (38) and (39), where A
t is the event that there
exists a particle that makes the short climb along optimal path (x̄, ȳ), and
t (ξ) is the event that the spine makes the short climb. Note that ε controls
the proximity to x̄ and δ the proximity to ȳ. Importantly, we will only be
interested in taking λ= λ̄ throughout this section, although we will usually
just write λ for notational simplicity. Also recall throughout that t and τ
are related through (θ/(2µλ)) exp(2µλτ) = κ
A TYPED BRANCHING DIFFUSION 29
Proof of Theorem 7. The key step in the proof of this is the following
use of the spine change of measure: for any function g :R+ →R+ we have
P x,y(A
t ) =Q
Ẑ+λ (τ)
Ẑ+λ (τ)
≥ Q̃x,yλ
Ẑ+λ (τ)
≥ Q̃x,yλ
Ẑ+λ (τ)
; sup
s∈[0,τ ]
Ẑ+λ (s)≤ g(τ)
≥ g(τ)−1Q̃x,yλ
t (ξ); sup
s∈[0,τ ]
Ẑ+λ (s)≤ g(τ)
Essentially we just have to make the “correct” choice for both λ and g in
expression (48), although there will still remain a number of technicalities
to resolve.
The first idea is to ensure the (originally rare) event A
t actually occurs
under the new measure Q̃
λ by making the spine follow close to the required
path (x̄, ȳ); this is achieved by choosing the optimal value λ̄ for λ and choos-
ing τ to be on the natural time scale it would take the spine to reach position
t. In particular, this choice will mean that in the first line of the above set
of inequalities there is no significant loss of mass when replacing the event
t with A
t (ξ). Next, we wish to choose the smallest possible g that will
still leave some positive probability on the last line of the above argument.
Hence, we wish to identify the rate of growth of the martingale Z+λ under
λ , and this will essentially be governed by the contribution from the spine
itself.
With this is mind, and recalling the various properties of the optimal
paths and parameters from Section 4, for ε0 > 0 we define
gε0(τ) := exp
ψ+λ +
e2µλτ − (ψ+λ y
2 + λx)
and recall from (35) that the scaling between t and τ is fixed throughout,
where κ2t= (θ/(2µλ̄))e
2µλ̄τ for large t, hence t+ τ ∼ t. Note that since we
are only considering the optimal value λ= λ̄, we have
e2µλ̄τ = (κ2ψ+
− λ̄β)t=Θ(β,κ).
Then from (48) we have
P x,y(A
t )≥ gε0(τ)−1Q̃
t (ξ); sup
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
.(49)
Our strategy for the rest of this proof is to show that the Q̃
λ -probability in
(49) is at least some ε′ > 0 for all sufficiently large t, uniformly for y ∈ [y0, y1],
30 Y. GIT, J. W. HARRIS AND S. C. HARRIS
so that the decay rate part of (49) matches the desired rate in the statement
of the theorem.
Conditioning on the spine’s path and birth times, G̃∞, and then making
use of some standard properties of conditional expectation we have
t (ξ); sup
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
t (ξ); sup
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
since A
t (ξ) is G̃∞-measurable. We next observe that, conditional on G̃∞,
we can write Ẑ+λ (t) as
Ẑ+λ (t) = e
y2+λx)
λ (t− Su) + f(t)
,(50)
where the Z
λ are independent copies of Z
λ started from a single particle
at (ξSu , ηSu); and f(t) is the contribution to Z
λ (t) from the spine, which,
conditional on G̃∞, is a known function of t. Now if we could show, for
0< ε̃0 < ε0,
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
and sup
s∈[0,τ ]
(Ẑ+λ (s)− f̂(s))≤
gε0(τ)
where f̂(t) := e−(ψ
y2+λx)f(t), we would have sups∈[0,τ ] Ẑ
λ (s)≤ gε0(τ). Hence,
defining Ẑ+λ (s) := Ẑ
λ (s)− f̂(s), we have
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
≥ Q̃x,yλ
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
; sup
s∈[0,τ ]
Ẑ+λ (s)≤
gε0(τ)
1{sups∈[0,τ ] f̂(s)≤gε̃0 (τ)/2}
s∈[0,τ ]
Ẑ+λ (s)≤
gε0(τ)
since, conditional on G̃∞, the supremum of f̂ on [0, τ ] is known.
We see from (50) that, conditional on G̃∞, Ẑ+λ (t) is a submartingale. This
is because the Q̃
λ -conditional expectation of each of the Z
λ in the sum
y2+λx)
λ (t− Su)(51)
A TYPED BRANCHING DIFFUSION 31
is constant, so the expectation of the sum cannot decrease, and in fact this
expectation increases every time there is a birth on the spine. Then by
Doob’s submartingale inequality we have
s∈[0,τ ]
Ẑ+λ (s)≤
gε0(τ)
= 1− Q̃x,yλ
s∈[0,τ ]
Ẑ+λ (s)≥
gε0(τ)
≥ 1− 2
gε0(τ)
λ (Ẑ
λ (τ)|G̃∞).
We must note here that the expectation on the above line is not a priori
finite. However, the expectation of each term in the sum (51) is bounded by
sups∈[0,τ ] f̂(s), which we have control over via an indicator function and so
we do not have to worry about this expectation blowing up.
So we need to show that for all sufficiently large τ and all y ∈ [y0, y1],
(ξ)∩{sups∈[0,τ ] f̂(s)≤gε̃0 (τ)/2}
gε0(τ)
λ (Ẑ
λ (τ)|G̃∞)
> ε′,
and hence also
t (ξ); sup
s∈[0,τ ]
Ẑ+λ (s)≤ gε0(τ)
as required. This will follow by combining both parts of the following result.
Lemma 12. Fix y1 > y0 > 0 and ε0 > ε̃0 > 0.
(i) For all sufficiently small ε, δ > 0, there exists some ε′ > 0 and T̃ > 0
such that for all y ∈ [y0, y1] and all t > T̃ ,
t (ξ); sup
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
> ε′.
(ii) As τ →∞,
gε0(τ)
λ (Ẑ
λ (τ)|G̃∞); sup
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
uniformly over y ∈ [y0, y1].
Then we have shown that, for any ε0 > 0, y1 > y0 > 0, and sufficiently
small ε, δ > 0, there exists a T > 0 such that, for all y ∈ [y0, y1] and all t > T ,
t−1 logP x,y(Aε,δt )≥−(Θ(β,κ) + ε0).
32 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Finally, we observe that the probability P x,y(A
t ) is trivially monotone
increasing in both ε and δ, and so it follows that if the result is true for all
sufficiently small ε and δ, it is in fact true for all ε, δ > 0. This completes
the proof of Theorem 7. �
Proof of Lemma 12(i). We will prove Lemma 12(i) in a sequence of
other lemmas, using a convenient coupling for the spine’s type process.
First recall that, under Q̃
λ , ηs solves the SDE
dηs =
θ dBs + µληs ds,
where Bs is a Q̃λ-Brownian motion. Noting that d(e
−µλsηs) = e
θ dBs,
we can construct e−µλsηs as a time-change of a Brownian motion with
e−µλsηs − η0 =
e−µλw dBw =
B̃(1− e−2µλs),
where B̃ is also a Q̃
λ -Brownian motion started at the origin.
In this way, we will construct ηy under P from a Brownian motion By
started at y
2µλ/θ where, for s ∈ [0,∞),
ηy(s) =
eµλsBy(1− e−2µλs).
To construct simultaneously all type processes ηy under the same measure
P, we first construct the process By0 as an independent Brownian motion
started at y0
2µλ/θ. Second, we construct the process B
y1 by running an
independent Brownian motion started at y1
2µλ/θ until it first hits the
path of By0 , at which point we couple the two processes together. Next, for
any other y ∈ (y0, y1), we run an independent Brownian motion By until it
first meets with either the process By0 below or By1 above, at which point
we couple it to the process it first hits.
Finally, we construct all the corresponding spatial processes ξy under P
from a single Brownian motion W by defining
ξy(s) =W
ηy(w)2 dw
ηy(w)2 dw,(52)
where W is started at x and is independent of the By processes.
Constructed in this way, for each y ∈ [y0, y1], the P-law of (ξy, ηy) is the
same as the Q̃
λ -law of (ξ, η).
Fixing µ ∈ (0,1) and K >max{y1,1}, we define the events and stopping
times
Ayε :=
By(s) ∈
,∀s ∈ (1− µ,1]
T0 := inf{t :By0(t) = 0}, TK := inf{t :By1(t) =K},
Ãε,K :=A
ε ∩Ay1ε ∩ {T0 > 1} ∩ {TK > 1}.
A TYPED BRANCHING DIFFUSION 33
Then, clearly P(Ãε,K)> 0 and, on the event Ãε,K , the coupling gives
0< ηy0(s)≤ ηy(s)≤ ηy1(s)≤K
eµλs,
for all s≥ 0 and y ∈ [y0, y1]. Note that our construction also ensures that if
event Ay0ε ∩Ay1ε occurs then so must Ayε for any y ∈ [y0, y1], hence Ayε ⊃ Ãε,K .
Lemma 13. Let ε > 0. On event Ãε,K , there exists a deterministic time
s0 = s0(ε)> 0 such that for all τ > s0,
s∈[0,τ ]
|ηy(s)− ȳ(s)| ≤ ε
for all y ∈ [y0, y1].
Proof. Set s1 =− 12µλ logµ and then, on event Ãε,K , for all τ ≥ s > s1
we have
ηy(s)−
eµλs,
for all y ∈ [y0, y1]. Writing
ȳ(s) =
1− e−2µλs
1− e−2µλτ
eµλs,
we see that there exists s2 = s2(ε)> 0 such that, for τ ≥ s > s2,
ȳ(s)−
eµλs.
Taking s3(ε) = max{s1, s2(ε)} now yields
|ηy(s)− ȳ(s)|< ε
eµλs ≤ ε
t(53)
for all τ ≥ s > s3 and all y ∈ [y0, y1].
Now consider s ∈ [0, s3]. On Ãε,K we have
|ηy(s)− ȳ(s)| ≤
eµλs3(1 +K),
and hence for some s4(ε)> 0 we have |ηy(s)− ȳ(s)| ≤ ε
t for all τ > s4, all
s ∈ [0, s3], and all y ∈ [y0, y1]. Taking s0(ε) =max{s3, s4} yields the result.
34 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Lemma 14. Let δ > 0. Then for all sufficiently small ε, there exists a
deterministic τ0 = τ0(ε, δ)> 0 such that, on Ãε,K , we have
s∈[0,τ ]
ηy(w)2 dw−
ȳ(w)2 dw
< δt(54)
for all τ > τ0 and all y ∈ [y0, y1].
Proof. Given any δ > 0, we first fix an ε > 0 sufficiently small such
that ε(2+ ε
; this yields a corresponding s3 = s3(ε), which is chosen
as at equation (53). Given this s3, we find τ1 = τ1(ε, δ)> 0 such that, for all
τ > τ1,
(K2 + 1)
e2µλw dw <
We now set τ0 = τ0(ε, δ) = max{s3, τ1}. With this choice of ε and τ0, we
proceed to show that the inequality (54) is satisfied. Note that τ0 is deter-
ministic and independent of y.
From equation (53) we see that, on Ãε,K and for s > s3,
ηy(w)2 dw ≥
ηy(w)2 dw+
ȳ(w)− ε
ȳ(w)2 dw−
ȳ(w)2 dw− 2
e2µλw dw
ȳ(w)2 dw−
e2µλw dw− (2ε) κt
ȳ(w)2 dw− δ
for all τ > τ0 and all y ∈ [y0, y1]. Similarly
ηy(w)2 dw ≤
ηy(w)2 dw+
ȳ(w) +
ȳ(w)2 dw+
ȳ(w)2 dw+K2
e2µλw dw+ ε
ȳ(w)2 dw+
A TYPED BRANCHING DIFFUSION 35
for all τ > τ0 and all y ∈ [y0, y1]. Finally, for s ∈ [0, s3], on Ãε,K we have
ηy(w)2 dw−
ȳ(w)2 dw
ηy(w)2 dw+
ȳ(w)2 dw
≤ (K2 + 1)
e2µλw dw < δt
for all τ > τ0 and all y ∈ [y0, y1]. �
Lemma 15. Let δ > 0. Then for all sufficiently small ε > 0, there exists
P-almost everywhere on Ãε,K a random time S0 = S0(δ, ε)<∞ such that
s∈[0,τ ]
ξy(s)− λa
ηy(w)2 dw
< δt,
for all y ∈ [y0, y1] and all τ > S0.
Proof. Given δ > 0, choose any δ′, δ′′ > 0 such that δ′(|β/λ|+ δ′′)< δ.
Recalling the construction of ξy at (52), we see from standard properties of
Brownian motion that there almost surely exists some S1 = S1(δ
′)<∞ such
s∈[0,t]
|W (s)|< δ′ for all t > S1.
s∈[0,τ ]
ηy(w)2 dw
ηy(w)2 dw
for all τ such that a
y(w)2 dw > S1, and by the coupling construction,
on Ãε,K this is true for all y ∈ [y0, y1] if a
y0(w)2 dw > S1. Then there
exists (P-almost everywhere on Ãε,K) a random time S2 = S2(δ
′)<∞, which
depends on By0 and S1, such that a
y(w)2 dw > S1 for all y ∈ [y0, y1] when
τ > S2.
Now by Lemma 14, given δ′′ and a sufficiently small ε, there exists a
deterministic τ0 = τ0(ε, δ
′′)> 0 such that, on Ãε,K ,
ηy(w)2 dw ≤ a
ȳ(s)2 ds+ δ′′t=
+ δ′′
t(56)
for all τ > τ0 and all y ∈ [y0, y1]. Combining the inequalities at (55) and (56),
we now see that, for τ > S0 = S0(ε, δ
′, δ′′) = max{S2, τ0},
s∈[0,τ ]
ξy(s)− λa
ηy(w)2 dw
= sup
s∈[0,τ ]
ηy(w)2 dw
for all y ∈ [y0, y1]. �
36 Y. GIT, J. W. HARRIS AND S. C. HARRIS
On combining Lemmas 14 and 15 and recalling the definition of optimal
path x̄ at (37), we obtain the following:
Lemma 16. Let δ > 0. Then for all sufficiently small ε > 0, there exists
P-almost everywhere on Ãε,K a random time S̃0 = S̃0(δ, ε)<∞ such that
s∈[0,τ ]
|ξy(s)− x̄(s)|< δt,
for all y ∈ [y0, y1], and all τ > S̃0.
We may now draw everything together to finish the proof of Lemma 12(i).
First we observe that since λ < 0, on event A
t (ξ),
s∈[0,τ ]
y2+λx) exp(ψ+λ η
s + λξs −E+λ s)
≤ e−(ψ
y2+λx) exp(ψ+λ (κ+ ε)
2t+ λ(−β − δ)t),
and so, given ε̃0, we can choose first δ and then ε sufficiently small so that
t (ξ)⊂
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
and, from Lemmas 13 and 16, there exists a random time T̃ = T̃ (δ, ε) <∞
such that on Ãε,K we have
s∈[0,τ ]
|ηy(s)− ȳ(s)|< ε
s∈[0,τ ]
|ξy(s)− x̄(s)|< δt
for all τ > T̃ and all y ∈ [y0, y1]. That is, Ãε,K ∩ {T̃ < τ} ⊂Aε,δt (ξy) for each
y ∈ [y0, y1], with the slight abuse of notation that
s∈[0,τ(t)]
|ηy(s)− ȳ(s)|< ε
t; sup
s∈[0,τ(t)]
|ξy(s)− x̄(s)|< δt
Note also that P(Ãε,K)> ε
′ for some ε′ > 0.
Combining the above, for any y ∈ [y0, y1] we have
t (ξ); sup
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
t (ξ)) = P(A
≥ P(Ãε,K; T̃ < τ)→ P(Ãε,K)
as τ →∞, as required. �
A TYPED BRANCHING DIFFUSION 37
Proof of Lemma 12(ii). Consider the expectation of the “sum term.”
We have
λ (Ẑ
λ (τ)|G̃∞) = e
y2+λx)
λ (t− Su)
= e−(ψ
y2+λx)
λ (t− Su)|G̃∞)
≤ e−(ψ
y2+λx)nτ max{eψ
η(Su)
2+λξ(Su)−E+λ Su :u < ξτ}
≤ nτ sup
s∈[0,τ ]
f̂(s).
Hence
gε0(τ)
λ (Ẑ
λ (τ)|G̃∞); sup
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
≤ Q̃x,yλ
gε̃0(τ)
gε0(τ)
; sup
s∈[0,τ ]
f̂(s)≤ gε̃0(τ)
≤ e−(ε0−ε̃0)tQ̃x,yλ (nτ ),
and we can now calculate Q̃
λ (nτ ) = Q̃
λ (Q̃
λ (nτ |G∞)), where G∞ the
σ-algebra generated by the path of the spine (not including the birth times).
Conditional on G∞, nτ is a Poisson random variable with mean given by
0 2(rη
s + ρ)ds, and using Fubini’s theorem we have
2(rη2s + ρ)ds
s)ds+2ρτ
e2µλτ −
− rθτ
+ 2ρτ
2y2κµλ
t+ o(τ).
So the Q̃
λ -expectation of nτ only grows linearly in t. Then since ε0− ε̃0 > 0,
the expression at (57) tends to 0 as t→∞. Moreover, the expectation at (57)
is bounded by the Q̃
λ -expectation, and hence the convergence is uniform
over y ∈ [y0, y1], as claimed. �
8. Martingale results. In this section we recall some existing and prove
some new martingale results that are intermediate steps in the proofs of
38 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Theorem 1 and the upper bound of Theorem 3. We recall from [13] that E−λ
[also written E−(λ)] and ∆(γ) are Legendre conjugates with
∆(γ) = inf
{E−(λ) + λγ}, E−(λ) = sup
{∆(γ)− γλ}.(58)
If, for λmin <λ< 0, we write γλ for the γ value which achieves the supremum
on the right-hand side of equation (58), then the functions λ 7→ γλ from
(−λmin,0) to (0,∞), and γ 7→ λγ from (0,∞) to (−λmin,0) are inverses of
each other and, of course, λγ is the λ value which achieves the infimum on
the left-hand side of equation (58). In addition, we note that
γλ =−
E−(λ) =
θa2λ2
θ− 8r− 4aλ2 ,(59)
that E−(λ) and ∆(γ) are convex functions, and that
c̃(θ) = sup{γ :∆(γ)> 0}= inf{−E−(λ)/λ :λmin <λ< 0}
= inf{c−λ :λmin < λ< 0}= c
λ̃(θ)
where
c−λ :=−E
λ /λ and
λ̃(θ) :=−
2(θ− 8r)(θρ+ 2ρ2 + rθ)
a(θ+4ρ)2
∈ (λmin,0).
A formula for c̃(θ) is given in equation (9). The following fundamental con-
vergence result for the Z−λ martingale was first partly proved in [13], but
also see [9] for a more complete proof using “spine” techniques.
Theorem 17. Suppose λ ∈ (λmin,0].
(i) If λ ∈ (λ̃(θ),0], the martingale Z−λ is uniformly integrable and has
an almost sure strictly positive limit.
(ii) If λ≤ λ̃(θ), then Z−λ (∞) = 0 almost surely.
The following convergence result was proved in [12] using martingales
based on Hermite polynomials.
Theorem 18. Let λ ∈ (λ̃(θ),0] and α< 1/4. For each P x,y starting law
and every continuous bounded function f :R 7→R, we have
f(Yu(t))e
αYu(t)
2+λ(Xu(t)+c
λ (∞),
A TYPED BRANCHING DIFFUSION 39
where
f0 :=
)1/4 ∫
f(y)eαy
y2φ(y)dy(62)
and φ(y) is the standard normal density.
In this paper, we require a corollary to this theorem which specifies more
precisely which particles contribute to the final limit.
Corollary 19. Let λ ∈ (λ̃(θ),0] and α < 1/4. For each P x,y starting
law and every continuous bounded function f :R 7→ R, we have for every
ε > 0
f(Yu(t)) e
αYu(t)
2+λXu(t)−E−λ t
1{|Xu(t)/t+γλ|<ε}−→a.s. f0Z
λ (∞)(63)
where γλ =− ∂∂λE
−(λ) and f0 is given at equation (62).
This last result will enable us to show in Section 10 that the almost sure
growth rate is at least as large as the expected growth rate, D(γ)≥∆(γ). It
is easy to see from Corollary 19 that when Z−λ (∞)> 0, there must exist at
least one particle near to −γλt in space. Further, because of the decay rate
of each term in the sum over particles at equation (63), it will be relatively
straightforward to improve this to get the required exponential numbers of
particles, exp(∆(γ)t), near −γλt for large times [as long as Z−λ (∞)> 0].
The following result concerns the rate at which the martingales Z+λ and
Z−λ converge to zero.
Theorem 20. Let λ ∈ (λmin,0). For every starting law, P x,y,
logZ±λ (t)
→ λ(c±λ − c
λ) a.s.
where c±λ is given at (5), and
c∗λ :=
c̃(θ), if λmin <λ≤ λ̃(θ),
c−λ , if λ̃(θ)≤ λ < 0.
Corollary 21. If λ ∈ (λmin,0), then Z+λ (t)→ 0 P x,y-almost surely.
The rate of convergence of the Z+λ martingale in part (i) of Theorem 20
is crucial in Section 9 to obtain the upper bound on the almost sure growth
rate, D(γ,κ) ≤ ∆(γ,κ). We also comment that if Corollary 19 were true
for all α < ψ+λ , then we could have gained this upper bound at that point.
40 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Although Corollary 19 is only proven for α < 1/4 (where we can utilize
suitable Hermite expansions), we conjecture that it holds for all α <ψ+λ .
Proof of Corollary 19. Let ε > 0 be small, µ := λ−ε, λ,µ ∈ (λ̃(θ),0),
f be a positive, continuous bounded function, α< 1/4 and note that γµ > γλ.
Then we have
f(Yu(t))e
αYu(t)
2+λXu(t)−E−λ t
1{Xu(t)<−γµt}
≤ e(E
µ −E−λ −εγµ)t
f(Yu(t))e
αYu(t)
2+µXu(t)−E−µ t
1{Xu(t)<−γµt}
f(Yu(t))e
αYu(t)
2+µXu(t)−E−µ t
−E−µ +(λ−µ)γµ)t.
Recall that E−(λ) is convex with ∂
E−(λ)≥ 0 and ∂
E−(λ) = γλ, so, from
the Taylor expansion,
E−λ −E
µ + (µ− λ)
E−(λ)
(µ− λ)2
E−(λ) + o((µ− λ)2).
Then taking ε > 0 small enough so that E−λ −E−µ +(λ−µ)γµ > 0, and using
Theorem 18, we find that for any δ > 0
limsup
f(Yu(t))e
αYu(t)
2+λXu(t)−E−λ t
1{Xu(t)<−(γλ+δ)t} = 0.
Similarly, we can show
limsup
f(Yu(t))e
αYu(t)
2+λXu(t)−E−λ t
1{Xu(t)>−(γλ−δ)t} = 0,
and hence the only contribution to the limit comes from the particles near
−γλt in space. Combining this with Theorem 18 we have
λ (∞) = limt→∞
f(Yu(t))e
αYu(t)
2+λXu(t)−E−λ t
= lim
f(Yu(t))e
αYu(t)
2+λXu(t)−E−λ t
1{|Xu(t)/t+γλ |<δ}.
Proof of Theorem 20. We use a useful technique brought to our
attention in [22]. Let p ∈ (0,1) so that, by Jensen’s inequality, Z±λ (t)p is a
supermartingale; then for u, v > 0 we have
(u+ v)p ≤ up + vp,
A TYPED BRANCHING DIFFUSION 41
and hence
Z±λ (t)
Yu(t)
2+λ(Xu(t)+c
Yu(t)
2+pλ(Xu(t)+c
For any ε > 0, Doob’s supermartingale inequality says
s≤w≤s+t
Z±λ (w)
p > εp
λ (s)
≤ ε−p
Yu(s)
2+pλ(Xu(s)+c
and then
s≤w≤s+t
eδwZ±λ (w)> ε
s≤w≤s+t
Z±λ (w)
p > e−pδ(s+t)εp
≤ ε−pepδt
Yu(s)
2+pλ(Xu(s)+c
p(λ(c±
)+δ)s
Now, if we can choose p ∈ (0,1) such that λ(c±λ − c
pλ) + δ < 0 and pψ
ψ+pλ, we must have e
δuZ±λ (u)→ 0 almost surely by using a familiar Borel–
Cantelli argument. [The condition pψ±λ < ψ
pλ guarantees that the expec-
tation in the last line above tends to a finite limiting value, hence stays
bounded over all times s, as can be checked by using formula (17), for ex-
ample.]
For all 0 ≤ p < 1 we find pψ±λ < ψ
pλ. Considering the graph of c
quickly see that, for λ ∈ [λ̃(θ),0), taking p as close to 1 as we like gives the
best rate. For λ ∈ [λmin, λ̃(θ)) we can choose p so that pλ= λ̃(θ), which gives
the best rate.
Recall from Theorem 17 that Z−λ (∞)> 0 when λ ∈ (λ̃(θ),0). Then, so far,
we have proved the following:
Lemma 22. For every starting law, P x,y, and for all ε > 0, if λ ∈ (λmin,0)
e−εt e−λ(c
)tZ±λ (t)→ 0 a.s.
42 Y. GIT, J. W. HARRIS AND S. C. HARRIS
where
c∗λ :=
c̃(θ), if λmin <λ≤ λ̃(θ),
c−λ , if λ̃(θ)≤ λ < 0.
It is clear that this gives the required upper bound of
lim sup
logZ±λ (t)
≤ λ(c±λ − c
Now, for any ε > 0, if λ ∈ (λmin, λ̃(θ)] then
Yu(t)
2+λ(Xu(t)+c̃(θ)t) ≥ eλ(Lt+c̃(θ)t)+εt →∞ a.s.
since we know that Lt := inf{Xu(t) : u ∈ N(t)} satisfies Lt/t→ −c̃(θ) a.s.
Otherwise, with λ∈ (λ̃(θ),0),
Yu(t)
2+λ(Xu(t)+c
t) ≥ eεtZ−λ (t)→∞ a.s.
since here Z−λ (∞)> 0 a.s. Thus, in all cases,
lim inf
logZ±λ (t)
≥ λ(c±λ − c
which completes the proof of Theorem 20. �
9. Proof of Theorem 3. Upper bound. The idea for the upper bound
proof is to overestimate indicator function by exponentials, and then re-
arrange the expressions to form martingale terms.
Simply observe that for λ ∈ (λmin,0),
Nt(γ, [κ
t,∞)) =
1{Xu(t)≤−γt;Yu(t)≥κ
1{Xu(t)≤−γt;Yu(t)2≥κ2t}e
(Yu(t)
2−κ2t)+λ(Xu(t)+γt)
≤ e(E
−κ2ψ+
+λγ)t
Yu(t)
2+λXu(t)−E+λ t
≤ e−λ(c
)tZ+λ (t)e
+λγ−κ2ψ+
where E±λ =−λc
Recall from equations (11) and (32) that E−λ + λγ − κ2ψ
λ has a minimal
value of ∆(γ,κ) achieved when λ= λ̄(γ,κ). Since c̃(θ) is the minimal value
of cλ, Theorem 20 implies that
lim sup
t−1 logZ+λ (t)≤ λ(c
λ − c
λ )(65)
A TYPED BRANCHING DIFFUSION 43
almost surely for all λ ∈ (λmin,0).
In cases where ∆(γ,κ)< 0, we can use the optimal value for λ, Theorem 20
and trivially note that Nt(γ, [κ
t,∞)) is integer valued to deduce that
1{Yu(t)≥κ
t;Xu(t)≤−γt} = 0
eventually, almost surely. Hence, D(γ,κ) =−∞ almost surely if ∆(γ,κ)< 0.
Otherwise we have ∆(γ,κ)≥ 0, which in fact guarantees that γ ∈ (0, c̃(θ)]
and hence λ̄(γ,κ) ∈ [λ̃(θ),0). Then since
lim sup
t−1 logNt(γ, [κ
t,∞))
≤ lim sup
t−1 log(e−λ(c
)tZ+λ (t)) + (E
λ + λγ − κ
2ψ+λ )
we can again make use of Theorem 20 and the minimizing λ value, λ̄(γ,κ),
to get the bound
limsup
t−1 logNt(γ, [κ
t,∞))≤∆(γ,κ) almost surely,
as desired.
Notice that, when ∆(γ,κ) = 0, the right-hand side of the inequality at
(64) will tend to infinity (see Corollary 19). Then, on the boundary, we have
only shown that lim sup t−1 logNt(γ, [κ
t,∞))≤ 0.
10. Proof of Theorem 1. The spatial growth rate. We first bound the
spatial growth rate above. Suppose that C ⊂ R is Borel-measurable with
y2φ(y)dy > 0. Let λ ∈ (λmin,0), then
1{Xu(t)≤−γt;Yu(t)∈C} ≤
1{Yu(t)∈C}e
λ(Xu(t)+γt)
= e(E
+λγ)t
1{Yu(t)∈C} e
λXu(t)−E−λ t
≤ e(E
+λγ)tZ−λ (t).
Recalling equations (8) and (19), we therefore have
1{Xu(t)≤−γt;Yu(t)∈C} ≤ e
∆(γ)tZ−λγ (t).
Now if γ ≥ c̃(θ), corresponding to λγ ∈ (λmin, λ̃(θ)] and having ∆(γ)≤ 0,
we know from Theorem 17 that Z−λγ (∞) = 0 almost surely. Then,
γ > c̃(θ) ⇒
1{Xu(t)≤−γt;Yu(t)∈C} = 0 eventually, a.s.
44 Y. GIT, J. W. HARRIS AND S. C. HARRIS
Otherwise, if γ ∈ (0, c̃(θ)), corresponding to λγ ∈ (λ̃(θ),0) and having ∆(γ)>
0, Theorem 17 tells us that Z−λγ (∞)> 0 almost surely, hence
lim sup
t−1 log
1{Xu(t)≤−γt;Yu(t)∈C} ≤∆(γ).
Now we bound the growth rate from below. Let ε > 0 be small, λ̃(θ) <
λ < 0, and µ = λ − ε. We recall now that E−λ is convex so
≥ 0 and
γµ > γλ. Then
eλXu(t)−E
1{−(γλ+ε)t≤Xu(t)≤−(γλ−ε)t;Yu(t)∈C}
eλ(−(γλ+ε)t)−E
1{−(γλ+ε)t≤Xu(t)≤−(γλ−ε)t;Yu(t)∈C}
= e(−λγλ−E
−λε)t ∑
1{−(γλ+ε)t≤Xu(t)≤−(γλ−ε)t;Yu(t)∈C}
≤ e(−λγλ−E
−λε)t ∑
1{Xu(t)≤−(γλ−ε)t;Yu(t)∈C}.
t−1 log
1{Yu(t)∈C}e
λXu(t)−E−λ t
1{|Xu(t)/t+γλ |<ε}
≤−λγλ−E−λ − λε+ t
−1 log
1{Xu(t)≤−(γλ−ε)t;Yu(t)∈C}.
Letting t→∞, using Corollary 19 and remembering that for λ̃(θ)< λ≤ 0
we have Z−λ (∞)> 0 a.s., we find
0≤−λγλ −E−λ − λε+ lim inft→∞ t
−1 log
1{Xu(t)≤−(γλ−ε)t;Yu(t)∈C}
and as ε > 0 can be arbitrarily small we have
lim inf
t−1 log
1{Xu(t)≤−γλt;Yu(t)∈C} ≥E
λ + λγλ.
Equivalently,
lim inf
t−1 log
1{Xu(t)≤−γt;Yu(t)∈C} ≥E
+ λγγ =∆(γ)
and hence the lim sup and lim inf agree as required.
We note that these proofs will easily adapt to cover a multi-type branching
Brownian motion where the types evolve as a finite state Markov chain,
such as found in [2], where it will also be possible to prove the analogous
A TYPED BRANCHING DIFFUSION 45
convergence theorem required when we have a finite type space by adapting
the proof of Theorem 18 found in [12].
In the standard branching Brownian motion case things are even simpler
to adapt (where, of course, there is no need for any convergence result akin to
Theorem 18). All the information necessary is contained in the martingales
u∈Nt exp(λXu(t)− (λ
2/2 + r)t) studied by Neveu [22] and, as first came
to our attention during discussions with J. Warren, the martingale with
parameter λ can only be capable of “counting” particles near γλt in space
at large times t, so when this martingale is uniformly integrable particles
must perpetually be found with the corresponding speed. Of course, in this
case more precise results, in the spirit of Watanabe [25], also exist.
Acknowledgments. We would like to thank two anonymous referees for
providing extremely helpful and thorough reviews of earlier incarnations
of this manuscript. Their numerous invaluable comments led to a much
improved presentation of this work.
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MR0237007
Y. Git
Statistical Laboratory
Cambridge University
22 Mill Street
Cambridge CB1 2HP
E-mail: [email protected]
J. W. Harris
Department of Mathematics
University of Bristol
University Walk
Bristol BS8 1TW
E-mail: [email protected]
S. C. Harris
Department of Mathematical Sciences
University of Bath
Bath BA2 7AY
E-mail: [email protected]
URL: http://people.bath.ac.uk/massch/
http://www.ams.org/mathscinet-getitem?mr=1417979
http://www.ams.org/mathscinet-getitem?mr=1991122
http://www.ams.org/mathscinet-getitem?mr=0199891
http://www.ams.org/mathscinet-getitem?mr=1601737
http://www.ams.org/mathscinet-getitem?mr=2037473
http://www.ams.org/mathscinet-getitem?mr=1601749
http://www.ams.org/mathscinet-getitem?mr=1349164
http://www.ams.org/mathscinet-getitem?mr=0400428
http://www.ams.org/mathscinet-getitem?mr=1952568
http://www.ams.org/mathscinet-getitem?mr=1046418
http://www.ams.org/mathscinet-getitem?mr=1659492
http://www.ams.org/mathscinet-getitem?mr=0758258
http://www.ams.org/mathscinet-getitem?mr=0237007
mailto:[email protected]
mailto:[email protected]
mailto:[email protected]
http://people.bath.ac.uk/massch/
Introduction
The branching model
Application to reaction–diffusion equations
Main results
Martingales
The asymptotic growth-rate of particles along spatial rays
The asymptotic shape and growth of the branching diffusion
Some expectation calculations
The expected rate of growth along spatial rays
The expected asymptotic shape
Short climb large deviation heuristics
A birth–death process
Finding the optimal path and probability
An important note on the optimal paths
Proof of Theorem 3. Lower bound
The ``spine'' setup and results
Proof of Theorem 7. The short climb probability
Martingale results
Proof of Theorem 3. Upper bound
Proof of Theorem 1. The spatial growth rate
Acknowledgments
References
Author's addresses
|
0704.0381 | The collision velocity of the bullet cluster in conventional and
modified dynamics | Mon. Not. R. Astron. Soc. 000, 1–8 (2007) Printed 27 October 2018 (MN LATEX style file v2.2)
The collision velocity of the bullet cluster in conventional
and modified dynamics
G. W. Angus1⋆, S. S. McGaugh2†
1SUPA, School of Physics and Astronomy, University of St. Andrews, Scotland KY16 9SS
2Department of Astronomy, University of Maryland, College Park, MD 20742-242 USA
Accepted ... Received ... ; in original form ...
ABSTRACT
We consider the orbit of the bullet cluster 1E 0657-56 in both CDM and MOND using
accurate mass models appropriate to each case in order to ascertain the maximum
plausible collision velocity. Impact velocities consistent with the shock velocity (∼
4700 km s
) occur naturally in MOND. CDM can generate collision velocities of at
most ∼ 3800 km s−1, and is only consistent with the data provided that the shock
velocity has been substantially enhanced by hydrodynamical effects.
Key words: gravitation - dark matter - galaxies: clusters: individual (1E 0657-56)
1 INTRODUCTION
Many lines of observational evidence now oblige us to be-
lieve that the universe is filled with a novel, invisible form of
mass that dominates gravitationally over normal baryonic
matter. In addition, a dark energy component which exerts
negative pressure to accelerate the expansion of the universe
is also necessary (Chernin et al. 2007). Though this ΛCDM
paradigm is well established, we still have only ideas about
what these dark components might be, and no laboratory
detections thereof.
One possible alternative to ΛCDM is the Modified New-
tonian Dynamics (MOND; Milgrom 1983a,b,c). This hy-
pothesis has been more successful than seems to be widely
appreciated (McGaugh & de Blok 1998; Sanders & Mc-
Gaugh 2002), and has received a theoretical boost from the
introduction of generally covariant formulations (Bekenstein
2004; Sanders 2005; Zlosnik, Ferreira, & Starkman 2006,7).
The dark matter and alternative gravity paradigms are rad-
ically different, so every observation that might distinguish
between them is valuable.
ΛCDM is known to work well on large scales (e.g.,
Spergel et al. 2006) while MOND is known to work well in in-
dividual galaxies (Sanders & McGaugh 2002). This success,
incorporating the tight correlation between dark and lumi-
nous mass in the DM framework (McGaugh 2005; Famaey
et al. 2007b) extends over five decades in mass (Fig.1) rang-
ing from tiny dwarfs (e.g., Milgrom & Sanders 2007) through
spirals of low surface brightness (de Blok & McGaugh 1998),
our own Milky Way (Famaey & Binney 2005) and other high
⋆ email: [email protected]
† email: [email protected]
surface brightness (Sanders 1996; Sanders & Noordermeer
2007) to massive ellipticals (Milgrom & Sanders 2003). The
recent observations of tidal dwarf galaxies by Bournaud et
al. (2007) provides a severe challenge to CDM but is natu-
rally explained in MOND with zero free parameters (Mil-
grom 2007; Gentile et al. 2007). Having said that, MOND
persistently fails to completely explain the mass discrepancy
in rich clusters of galaxies. Consequently, clusters require
substantial amounts of non-luminous matter in MOND.
That rich clusters contain more mass than meets the
eye in MOND goes back to Milgrom’s original papers (Mil-
grom 1983c). At the time, the discrepancy was very much
larger than it is today, as it was not then widely appreci-
ated how much baryonic mass resides in the intra-cluster
medium. Further work on the X-ray gas (e.g., Sanders 1994,
1999) and with velocity dispersions (McGaugh & de Blok
1998) showed that MOND was at least within a factor of a
few, but close inspection revealed a persistent discrepancy
of a factor of two or three in mass (e.g., Gerbal et al. 1992;
The & White 1998; Pointecouteau & Silk 2005, Buote &
Canizares 1994). Weak gravitational lensing (Angus et al.
2007a; Takahashi & Chiba 2007; Famaey, Angus et al. 2007)
provides a similar result.
To make matters worse, the distribution of the unseen
mass does not trace that of either the galaxies or the X-
ray gas (Aguirre et al. 2001; Sanders 2003; Angus et al.
2007b; Sanders 2007). In Fig. 1 we plot the baryonic mass
of many spiral galaxies and clusters against their circular
velocity together with the predictions of MOND and CDM.
MOND is missing mass at the cluster scale. CDM suffers an
analogous missing baryon problem on the scale of individual
galaxies.
The colliding bullet cluster 1E-0657-56 (Clowe et al.
c© 2007 RAS
http://arxiv.org/abs/0704.0381v2
2 G. W. Angus and S. S. McGaugh
2004,2006, Bradac et al. 2006, Markevitch et al. 2004,
Markevitch & Vikhlinin 2007) illustrates in a spectacular
way the residual mass discrepancy in MOND. While cer-
tainly problematic for MOND as a theory, it does not con-
stitute a falsification thereof. Indeed, given that the need for
extra mass in clusters was already well established, it would
have been surprising had this effect not also manifested it-
self in the bullet cluster. The new information the bullet
cluster provides is that the additional mass must be in some
collision-less form.
It is a logical fallacy to conclude that because extra
mass is required by MOND in clusters, that dark matter
is required throughout the entire universe. While undeni-
ably problematic, the residual mass discrepancy in MOND
is limited to groups and rich clusters of galaxies: these are
the only systems in which it systematically fails to remedy
the dynamical mass discrepancy (see discussion in Sanders
2003). Could we be absolutely certain that we had accounted
for all the baryons in clusters, then MOND would indeed
be falsified. But CDM suffers an analogous missing baryon
problem in galaxies (Fig.1) in addition to the usual dynam-
ical mass discrepancy, yet this is not widely perceived to be
problematic. In either case we are obliged to invoke the ex-
istence of some dark mass which is presumably baryonic (or
perhaps neutrinos) in the case of MOND. In neither case
is there any danger of violating big bang nucleosynthesis
constraints. The integrated baryonic mass density of rich
clusters is much less than that of all baryons; having the
required mass of baryons in clusters would be the proverbial
drop in the bucket with regards to the global missing baryon
problem.
A pressing question is the apparently high relative ve-
locity between the two clusters that comprise the bullet
cluster 1E 0657-56 (Clowe et al. 2006, Bradac et al. 2006,
Markevitch et al. 2004, Markevitch & Vikhlinin 2007). The
relative velocity derived from the gas shockwave is vrel =
4740+710
−550 kms
−1 (Clowe et al. 2006). Taken at face value,
this is very high, and seems difficult to reconcile with ΛCDM
(Hyashi & White 2006). The problem is sufficiently large
that it has been used to argue for an additional long range
force in the dark sector (Farrar & Rosen 2007). Here we ex-
amine the possibility of such a large velocity in both CDM
and MOND.
One critical point that has only very recently been ad-
dressed is how the shock velocity relates to the collision ve-
locity of the clusters. Naively, one might expect the dissipa-
tional collision of the gas clouds to slow things down so that
the shock speed would provide a lower limit on the colli-
sion speed. Recent hydrodynamical simulations (Springel &
Farrar 2007; Milosavljevic et al. 2007) suggest the opposite.
A combination of effects in the two hydrodynamical sim-
ulations show that the shock velocity may be higher than
the impact velocity. The results of the two independent hy-
drodynamical simulations do not seem to be in perfect con-
cordance, and the precise result seems to be rather model
dependent. Nevertheless, it seems that the actual relative
velocity lies somewhere in the range 3500-4500 kms−1.
The difficulties posed by a high collision velocity for
CDM have been discussed previously by Hayashi & White
(2006) and Farrar & Rosen (2007). Whereas Springel & Far-
rar (2007) and Milosavljevic et al. (2007) consider the com-
plex hydrodynamic response of the two gas clouds during
Figure 1. Shows baryonic mass against circular velocity. Rotat-
ing galaxies (blue circles) are from McGaugh (2005) and clusters
(green triangles) are from Sanders (2003) using the measured tem-
perature to estimate the circular velocity assuming isothermality.
The solid orange line is the CDM M-V relation (Steinmetz &
Navarro 1999) assuming Mb = fbMvir with fb = 0.17 (Spergel et
al. 2006) and the dashed red line is the MOND prediction. The
spirals lie directly on the MOND prediction, but the clusters are
generally 2-3 times in mass below it. The CDM expectation is
nicely consistent with clusters, but implies many dark baryons in
spirals in addition to the non-baryonic dark matter.
the ongoing collision, here we investigate the ability of two
clusters like those comprising the bullet cluster to accelerate
to such a high relative velocity in the case of both CDM and
MOND prior to the merger. We compute a simple free fall
model for the two clusters in an expanding universe with
realistic mass models, and ask whether the observed colli-
sion velocity can be generated within the time available. We
take care to match the mass models to the specific observed
properties of the system appropriate to each flavor of grav-
ity in order to realistically evaluate the orbit of the clusters
prior to their collision.
2 MODELING THE FREEFALL
We wish to address a simple question. Given the observed
masses of the two clusters, is it possible to account for the
measured relative velocity from their gravitational freefall?
The expansion of the universe mitigates against large veloc-
ities, since the clusters must decouple from the Hubble flow
before falling together. Presumably it takes some time to
form such massive objects, though this is expected to occur
earlier in MOND than in ΛCDM (Sanders 1998, 2001; Mc-
Gaugh 1999, 2004; Nusser 2002; Stachniewicz & Kutschera
2002; Knebe & Gibson 2004; Dodelson & Liguori 2006). The
clusters are observed at z=0.3, giving at most 9Gyrs for
them to accelerate towards each other. This imposes an up-
per limit to the velocity that can be generated gravitation-
ally. Without doing the calculation, it is not obvious whether
the larger masses of the clusters in CDM or the stronger long
range force in MOND will induce larger relative velocities.
Since we know the state of the system directly prior
to collision, it makes sense to begin our simulations from
the final state and work backwards in time towards when
the relative velocity was zero. This point, where the clus-
c© 2007 RAS, MNRAS 000, 1–8
The collision velocity of the bullet cluster in conventional and modified dynamics 3
ters have zero relative velocity, is when they turned around
from the Hubble flow and began their long journey gravi-
tating towards each other. Working backwards in time leads
to potentially counter-intuitive discussions (such as Hubble
contraction), which we try to limit.
We must account for the Hubble expansion in a manner
representing the universe before z=0.3. The detailed form of
the expansion history of the universe a(t) is not known in
the case of MOND, so we take the scale factor of ΛCDM in
both models
da(t)
+ ΩΛa
. (1)
Where we take Ho = 72 km s
−1Mpc−1, Ωm=0.27 and ΩΛ =
0.73.
The important aspect is the basic fact that the uni-
verse is expanding and the mutual attraction of the clusters
must overcome this before they can plunge together at high
velocity.
We implement the scale factor in the simulations
through the equation of motion
[a(t)v] = g. (2)
Computing this numerically, from time step to time step we
calculate the ratio of the scale factor in the previous time
step to the current time step (i.e. ǫ = a(ti−1)/a(ti); we use
negative time steps to move backwards in time from the
presently known configuration, so ǫ > 1 (higher i means
earlier universe). We then have v(ti) = v(ti−1)ǫ+ g∆t.
The right hand side of Eq. 2 differs in MOND and CDM
not only because the law of gravity is altered, but also be-
cause the gravitating masses are higher in CDM.
The initial conditions are the crux of the problem, with
at least 4 unknowns. These include the masses of the two
clusters, the relative velocity of the clusters, and the dis-
tance of separation between the two when they had this
relative velocity. The separation is the same in MOND and
Newtonian gravity, but the Newtonian mass is higher.
The relative velocity of the two clusters can be mea-
sured because, in the last few 100Myrs, the less massive
sub cluster has passed through the centre of the more mas-
sive main cluster. The ram pressure has imposed a smooth
bow shock (Markevitch et al. 2004, Markevitch & Vikhlinin
2007) on the gas of the sub cluster. Since the relative veloc-
ity is the foundation of the problem we leave it free and try
to estimate it by fixing other variables. In our simulation,
we think it sensible to consider the separation of the two
clusters (i.e. of the two centres of mass) when they had the
calculated relative velocity to be when the leading edge of
the sub cluster’s gas cloud began to pass through the dense
region of gas belonging to the main cluster and separate
from the dark matter. It appears that the centre of the sub
cluster’s gas cloud (the location of the bullet) is preceeded
by the bow shock by around 200kpc further in the direc-
tion of travel. We take 200kpc to also be the radius within
which the gas of cluster 1 was dense enough to imprint the
bow shock. Indeed, the gas mass of the main cluster (sub
cluster) is only measured out to 180kpc (100kpc) and could
not be found further detailed in the literature. However,
the uncertainty is large enough that we wished to clarify
the impact of different initial separations by always using
a range of initial separations of between 350-500kpc. This
Figure 2. The total enclosed masses for the main cluster (black)
and sub cluster (red) for CDM (solid), MOND with standard µ
(dashed) and simple µ (dotted).
separation is defined as when the two pre-collision clusters
had the relative velocity of vrel related to the shock velocity
4740+710
−550 kms
−1. Now of course, they are on the opposite
sides on the sky after having passed through each other and
the gas has been offset from the DM.
3 THE COLLISION IN CDM
In the CDM framework, it is no problem to generate two
clusters in an N-body simulation and calculate all gravi-
tational accelerations exactly. However, in MOND we are
dealing with non linear gravity and the tools for such pur-
poses are only now being developed (Nipoti et al. 2007a).
Additionally, since we begin our simulations with the over-
lap of the two clusters, it is not guaranteed that the clusters
preserve their shapes as they separate. Furthermore, it was
not possible to simply include accretion history (Wechsler
et al. 2002) in the N-body simulations or easily vary the
truncation radius of the DM halo as was necessary.
Therefore, a better method was to semi-analytically ac-
count for these aspects in a simulation where gravity of one
cluster acting on the other is just the mass enclosed by a
sphere around the gravitating body’s centre of mass with
radius equal to the separation of the two cluster’s centres of
mass. For this procedure, the only two unknowns are the sep-
aration, which is initially known and computed each time-
step; and the mass enclosed. The enclosed mass depends on
the density profile of the two clusters and was fitted by Clowe
et al. (2006) using NFW profiles of the form (see Angus &
Zhao 2007)
ρ(r) =
m200r
r + r200c
= ln(1+c)− c
1 + c
where c is the concentration. The enclosed mass goes as
c© 2007 RAS, MNRAS 000, 1–8
4 G. W. Angus and S. S. McGaugh
Figure 3. The total enclosed masses for the main cluster as a
function of time, where 9.3Gyr marks the collision of the two
clusters and 0Gyr represents the Big Bang. The three lines corre-
spong to different mass assembly rates α=0.5 (black), 1.0 (blue)
and 1.5 (red). The mass loss is halted at m200/50.
m(r) = Am200
− 1 +
For the main cluster they give m200 = 1.5 × 1015M⊙,
r200 = 2100kpc and concentration c=1.94. For the sub clus-
ter, m200 = 1.5 × 1014M⊙, r200 = 1000kpc and c=7.12. We
augment the DM with a baryon fraction of 17% (Spergel et
al. 2006) which was part of the total mass during freefall.
We ran dynamical time steps (negative) such that
∆t = 10
d[pc]Myr (5)
Where d is the separation of the two centres of mass and
∆t has a maximum value of 1Myr. Since initially d∼ 400kpc,
the starting time steps are ∼0.06Myr. The simulations were
run until 9Gyr had elapsed.
The mass distributions as functions of radius for the
two clusters in CDM and the ones used in the MOND sim-
ulations are shown in Fig.2. A subtle point about the total
masses of the two clusters is that we do not expect the mass
to remain constant as we go back in time. Presumably they
grew from a seed of negligible mass at high redshift (see
discussion in Cameron & Driver 2007). This tends to im-
pede their freefall, reducing the maximum collision velocity
to ∼ 2900 kms−1 by the estimate of Farrar & Rosen (2007).
To include this, without the impedance, we use the proce-
dure of Wechsler et al. (2002) who used the relation
M(z) = M(z = 0.3)e
−α(z−0.3)
where α obviously encodes the speed of the accretion or
assembly of the halo. Typical values used in their work are
0.5< α <2.0. In Fig.3 we plot the mass enclosed within r200
for the two clusters as functions of reshift for α=0.5,1.0 and
1.5. Note we always keep a floor value of cluster mass of
m200/50 so the halo is never completely disassembled.
Another important point is whether mass integrated out
beyond r200 should be included, since the actual virial ra-
dius depends on both cosmology and redshift (Bullock et al.
2001). Indeed, the internal gravity of the main cluster has
not reached ao by r200 = 2100kpc meaning the MOND dy-
namical mass has not yet saturated. However, recall that it
takes ∼ (r200−d)/vrel = (2100kpc−400kpc)/3400 kms−1 ∼
500Myr for the clusters to separate enough for this extra
matter to even begin to manifest itself. The cluster is also
losing mass (backwards in time) due to accretion coupled
with the fact there must be overdensities on the opposite side
of the universe countering the influence of these overdensi-
ties. However, using Eq.4 it is straight-forward to include all
the enclosed mass out to any radius because the parameters
r200 and m200 do not explicitly force the enclosed mass to
truncate at r200, they simply define the shape of the profile.
4 THE COLLISION IN MOND
In MOND, the basic modification of purely Newtonian dy-
namics is
µ(x)g = gN, (7)
where gN is the Newtonian acceleration computed in the
usual way from the baryonic mass distribution, g is the ac-
tual acceleration (including the effective amplification due
to MOND conventionally ascribed to DM), ao is the char-
acteristic acceleration at which the modification becomes
effective (∼ 10−10 m s−2), x = g/ao, and µ(x) is an inter-
polation function smoothly connecting the Newtonian and
MOND regimes. In the limit of large accelerations, g ≫ ao,
µ → 1 and the Newtonian limit is obtained: everything be-
haves normally. The MOND limit occurs only for exceed-
ingly low accelerations, with µ → x for g ≪ ao. We im-
plement two possible versions of the interpolation function:
the ‘standard’ function traditionally used in fitting rotation
curves:
1 + x2
(e.g., Sanders & McGaugh 2002), and the ‘simple’ function
found by Famaey & Binney (2005) to provide a good fit the
terminal velocity curve of the Galaxy:
1 + x
. (9)
A well known problem with implementing the MOND
force law in numerical computations is that the original for-
mulation (Eq 7) does not conserve momentum (Felten 1984;
Bekenstein 2007). This was corrected with the introduction
of a Lagrangian formulation of MOND (Bekenstein & Mil-
grom 1984; Milgrom 1986) which has the modified Poisson
equation
∇ · [µ(|∇Φ|/ao)∇Φ] = 4πGρ. (10)
This formulation has been shown to obey the necessary
conservation laws (Bekenstein & Milgrom 1984; Bekenstein
2007). With some rearrangement, it leads to
µ(x)g = gN +∇× h, (11)
c© 2007 RAS, MNRAS 000, 1–8
The collision velocity of the bullet cluster in conventional and modified dynamics 5
which we recognize as Eq 7 with the addition of a curl field.
Unfortunately, implementing a numerical formulation of
the modified Poisson equation is not a simple one-line change
to typical N-body codes: this fails to obey the conservation
laws. Instead, one needs an entirely different numerical ap-
proach than is commonly employed. Progress has been made
along these lines (e.g., Brada & Milgrom 1995, 1999; Ciotti,
Londrillo, & Nipoti 2006; Nipoti, Londrillo & Ciotti 2007ab;
Tiret & Combes 2007; see also Nusser 2002; Knebe & Gib-
son 2004), but we do not seek here a full N-body treatment
of complex systems. Rather, we wish to develop and apply
a simple tool (Angus & McGaugh, in preparation) that can
provide some physical insight into basic problems. For the
specific case of the large collision velocity of the bullet clus-
ter, it suffices to treat the curl field as a small correction
to the center of mass motion (Milgrom 1986). The external
field effect (see Milgrom 1983a, Bekenstein 2007) is crudely
approximated as a constant of appropriate magnitude (Mc-
Gaugh 2004). We checked the effect of varying the external
field, which is modest. It is not possible to do better without
complete knowledge of the mass distribution in the environ-
ment of the clusters.
When modeling the bullet cluster in MOND, Angus et
al. (2007a) fitted the convergence map of Clowe et al. (2006)
using spherical potential models for the four mass compo-
nents. Their best fit gives masses for all four components in
MOND and standard gravity. Unfortunately, the map is only
sensitive out to 250kpc from the respective centres which ne-
glects an over large portion of the dynamical mass. So, in or-
der to remain consistent with the CDM simulations, we take
the NFW profile and calculate what the MOND dynamical
masses for the two commonly used interpolating functions
(Eq.8 & 9) are, as shown in Fig.2. The Newtonian mass for
the main cluster is twice that of the MOND dynamical mass
with the standard µ and three times when the simple µ is
used.
The mutual gravity imposed upon the sub cluster by
the main cluster is
|gsub + gex|
gsub = gn,sub = −GMmain(d)/d2 (12)
and we simply swap the subscripts around to find the mutual
gravity of the sub cluster upon the main cluster. Following
on from above, d is the distance between the two centres of
mass and Mmain(d) is the enclosed mass within a radius d
from the centre of mass of the main cluster. The gex is the
external field limiting the MOND correction which comes
from large scale structure and is always assumed orthogonal
to the direction of gsub, making the argument of the µ func-
tion more easily expressed as
(g2sub+g
. The direction
and amplitude of gex is unknown at all times. The MON-
Dian additional acceleration becomes minor when the accel-
eration drops below gex. We use gex = ao/30 (Aguirre et
al. 2001; McGaugh 2004) which is roughly the external field
imposed on the Milky Way by M31 and vice versa (Famaey
et al. 2007a, Wu et al. 2007).
5 N-BODY COLLISION
Our first attempt at simulating the collision in Newtonian
gravity was using a standard N-body tree code. The benefit
it gives is that in principal, we can more accurately compute
the mutual gravity at the beginning of the simulation when
the two clusters overlap. However, this is fraught with diffi-
culties and inconsistencies. The first being that tidal effects
undoubtedly stretch the two clusters and 2-body interac-
tions may eject particles from the two halos. Therefore, it
makes better sense to begin such a simulation from high
redshift where the clusters are greatly separated and tidal
effects are negligible and let them freefall in the expanding
universe and when they collide, the tidal effects will be well
accounted for. Of course, the problem is that it is not triv-
ial to then sample collision velocities because the separation
and time at which the two clusters began their freefall is not
simply related. Moreover, the truncation of the two halos
and different mass models are not easily varied. Neverthe-
less, we did attempt a CDM N-body model with truncation
at r200 for both halos. We found a similar result to that from
the semi-analytical models of 3800 kms−1.
6 RESULTS
The ability of the two clusters that comprise the bullet to
bring each other to a halt at a finite time in the past is sen-
sitive to both the flavour of gravity at work and the true
relative velocity. For velocities larger than the maximum,
the relative velocity never reaches zero and increase sharply
at early times (large z). The two clusters do not gravitate
strongly enough to generate such high velocities and would
have to have had a huge relative velocity towards each other
in the early universe in order to overcome the Hubble expan-
sion and fall together with such a high relative velocity at
z=0.3. Fig.4 shows how the relative velocity of the two clus-
ters varies with time for a large sample of initial (meaning
collisional) relative velocities for a CDM and MOND sample
simulation. A difference of just 100 km s−1 can have a signif-
icant impact on the time required to generate such a large
velocity and by the same token, the longer the two clusters
free-fall, the larger a velocity they can generate. Sadly, there
is only a finite time (∼ 9Gyr)) since the Big Bang for this
to happen.
In Table 1 we’ve put the key results of the simulations
so as to give the reader a feel for what the maximum rela-
tive velocity that can be achieved is. Each velocity is that
achieved with an initial separation of 425kpc, where taking
350 or 500kpc induces an increase or decrease of 100 kms−1
which we take as the minimum error. The most extreme
CDM model is to have no truncation of the DM halos, ex-
tending them out to r1. This absurd extreme allows a maxi-
mum relative velocity of 4500 kms−1. Then, if we still allow
the halos to extend to r1, but account for some assembly of
the halos with α = 1 then the relative velocity reduces to
4200 kms−1.
More realistically, if we truncate the halos at r200 and
try four different halo assembly rates such that α =0.0, 0.5,
1.0 & 1.5 we get respective maximum relative velocities of
4000, 3900, 3800 and 3800 km s−1. These numbers repre-
sent the plausible maximum relative velocities in the CDM
framework.
For the MOND case we ran simulations with both the
simple (Eq.9) and standard (Eq.9) µ functions. The stan-
dard function leads to higher dynamical masses from the
c© 2007 RAS, MNRAS 000, 1–8
6 G. W. Angus and S. S. McGaugh
Figure 4. Shows the relative velocity of the two clusters plotted against time (a) CDM and (b) MOND. Time=0Myr is the current (z=0.3)
relative velocity of the two clusters with larger times corresponding to higher redshifts. Black lines correspond to relative velocities that
are achievable, whereas red lines are not. In (a) we use the simulation (CDM2c) for which α=1.0, d=425kpc and we truncate the halos at
r200. The relative velocities used are vrel=3500-4200 km s
−1 in intervals of 100 km s−1. In (b) we use the simulation (MONDst2) which
uses the standard µ function and α=0.5, d=425kpc. The relative velocities used are vrel=4100-4800 km s
−1 in intervals of 100 km s−1. The
4 dashed lines are the predicted relative velocities according to the mean and 1σ error of the original relative velocity from Markevitch &
Vikhlinin (2007) in blue, the simulations of Milosavljevic (2007) in green and Springel & Farrar (2007) in turquoise. The high observed
collision velocity is more readily obtained in MOND than CDM.
Model Max Vrel [ km s
−1] Truncation Radius α Gravity
CDM1a 4500 r1 0.0 Newtonian
CDM1b 4200 r1 1.0 Newtonian
CDM2a 4000 r200 0.0 Newtonian
CDM2b 3900 r200 0.5 Newtonian
CDM2c 3800 r200 1.0 Newtonian
CDM2d 3800 r200 1.5 Newtonian
MONDst1 4800 r200 0.0 MOND-standard µ
MONDst2 4500 r200 0.5 MOND-standard µ
MONDsi1 4600 r200 0.0 MOND-simple µ
MONDsi2 4500 r200 0.5 MOND-simple µ
Table 1. Shows the parameters used in the different models and gives the maximum attainable relative velocity for each.
NFW profile, but lower 2-body gravity. The standard (sim-
ple) function with no accretion and with α = 0.5 gener-
ate 4800 (4600) and 4600 (4500) kms−1 respectively and for
comparison, the maximum CDM velocity with those reduced
masses is just 2700 (2300) kms−1. This is a clear demonstra-
tion of the expectation in MOND for larger peculiar veloc-
ities. We use the lower assembly parameter α = 0.5 be-
cause structure is expected to form more swiftly in MOND
(Sanders 1998, 2001).
An important factor is that of the fitted NFW density
profile to the convergence map, in which matter is extrapo-
lated to 2100kpc and 1000kpc for the main and sub cluster
respectively. Presumably the significance of the detection of
this mass is negligible and the NFW fit has been made as-
suming if we know the details in the central 250kpc, then we
know the density out to r200. The mass sheet degeneracy is
broken by constraining the mass at the edges of the fit based
on the slope of the profile in the inner regions - but if the
mass profile is wrong then it could lead to the completely
wrong measurement for the value of the mass sheet (Clowe,
de Lucia and King 2004).
All of this means that the density profiles of the two
clusters could be moderately different in reality. However,
the actual shape of any profile is less important to the rela-
tive velocity than simply the normalisation of the total mass.
To this end we have simulated the collision with 10% more
and 10% less mass for both clusters (with assembly param-
eter α = 1). The effect is to increase (10% more mass) or
decrease (10% less mass) the relative velocity by 200 kms−1
from 4800 kms−1 for model MONDst1.
Another concern is that the clusters are unlikely to be
spherically symmetric (Buote & Canizares 1996) and are
presumably elongated in the direction of motion. Again this
could lead to an incorrect density profile, whereas ellipticity
itself would have little effect on our results.
c© 2007 RAS, MNRAS 000, 1–8
The collision velocity of the bullet cluster in conventional and modified dynamics 7
7 SUMMARY
We have constructed specific mass models for the bullet
cluster in both CDM and MOND. We integrate backwards
from the observed conditions to check whether the large
(∼ 4700 km s−1) apparent transverse velocity can be at-
tained in either context. We find that it is difficult to
achieve vrel > 4500 kms
−1 under any conditions. Never-
theless, within the range of the uncertainties, the appropri-
ate velocity occurs fairly naturally in MOND. In contrast,
ΛCDM models can at most attain ∼ 3800 kms−1 and are
more comfortable with considerably smaller velocities.
Taken at face value, a collision velocity of 4700 km s−1
constitutes a direct contradiction to ΛCDM. Ironically, this
cluster, widely advertised as a fatal observation to MOND
because of the residual mass discrepancy it shows, seems to
pose a comparably serious problem for ΛCDM. It has often
been the case that observations which are claimed to falsify
MOND turn out to make no more sense in terms of dark
matter.
Two critical outstanding issues remain to be clarified.
The first is the exact density profiles and virial masses of the
two clusters and the second is how the observed shock ve-
locity relates to the actual collision velocity of the two grav-
itating masses. The recent simulations of Springel & Farrar
(2007) and Milosavljevic et al. (2007) seem to suggest that,
contrary to naive expectations, hydrodynamic effects reduce
the relative velocity of the mass with respect to the shock.
A combination of effects is responsible, being just barely
sufficient to reconcile the data with ΛCDM. Hydrodynam-
ical simulations are notoriously difficult, and indeed these
two recent ones do not agree in detail. It would be excellent
to see a fully self-consistent simulation including both hy-
drodynamical effects and a proper mass model and orbital
computation like that presented here.
There are a number of puzzling aspects to the hydrody-
namical simulations. First of all, Springel & Farrar (2007)
use Hernquist profiles for the DM distribution in the clus-
ters and not NFW halos. Furthermore, they find that the
morphology of the bullet is reproduced only for a remark-
ably dead head-on collision. If the impact parameter is even
12kpc — a target smaller than the diameter of the Milky
Way — quite noticeable morphological differences ensue.
This can be avoided if the separation of mass centres hap-
pens to be along our line of sight — quite a coincidence
in a system already remarkable for having the vector of its
collision velocity almost entirely in the plane of the sky. Fur-
thermore, the mass models require significant tweaking from
that infered from the convergence map and are unable to re-
produce the currently observed, post merger positions of the
gas and DM. It appears to us that only the first rather than
the last chapter has been written on this subject. Getting
this right is of the utmost importance, as the validity of both
paradigms rests on the edge of a knife, separated by just a
few hundred km s−1.
More generally, the frequency of bullet-like clusters may
provide an additional test. The probability of high collision
velocities drops with dramatic rapidity in ΛCDM (Hyashi &
White 2006). In contrast, somewhat higher velocities seem
natural to MOND. Naively it would seem that high impact
velocity systems like the bullet would be part and parcel
of what might be expected of a MOND universe. With this
in mind, it is quite intriguing that many bullet cluster like
systems have been detected (although none quite as unique).
The dark ring around Cl0024+17 tentatively observed by
Jee et al. (2007; see also Famaey et al. 2007c), the dark core
created by the “train wreck” in Abell 520 by Mahdavi et al.
(2007), Cl0152+1357 (Jee et al. 2005a), MS1054+0321 (Jee
et al. 2005b) and the line of sight merger with > 3000 kms−1
relative velocity observed by Dupke et al. (2007) for Abell
576 may all provide examples and potential tests.
ACKNOWLEDGEMENTS
We acknowledge discussions with Benoit Famaey, Tom Zlos-
nik, Douglas Clowe, HongSheng Zhao, Greg Bothun, Moti
Milgrom, Bob Sanders and Ewan Cameron. GWA thanks
Steve Vine for his N-body tree code. GWA is supported by
a PPARC scholarship. The work of SSM is supported in part
by NSF grant AST0505956.
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Introduction
Modeling the freefall
The collision in CDM
The collision in MOND
N-body collision
Results
Summary
|
0704.0382 | On Some Subgroup Chains Related to Kneser's Theorem | On some subgroup chains related to Kneser’s theorem
Yahya O. Hamidoune∗ Oriol Serra† Gilles Zémor‡
March 28, 2007
Abstract
A recent result of Balandraud shows that for every subset S of an
abelian group G there exists a non trivial subgroup H such that |TS| ≤
|T |+ |S| − 2 holds only if H ⊂ Stab(TS). Notice that Kneser’s Theorem
only gives {0} 6= Stab(TS).
This strong form of Kneser’s theorem follows from some nice properties
of a certain poset investigated by Balandraud. We consider an analogous
poset for nonabelian groups and, by using classical tools from Additive
Number Theory, extend some of the above results. In particular we obtain
short proofs of Balandraud’s results in the abelian case.
1 Introduction
In order to avoid switching from multiplicative to additive notation, all groups
will be written multiplicatively.
Kneser’s addition theorem states that if S, T are finite subsets of an abelian
group G then |ST | ≤ |S| + |T | − 2 holds only if ST is periodic (i.e, there
is a non trivial subgroup H such that HST = ST .) Kneser’s Theorem is a
fundamental tool in Additive number Theory. Proofs of this result may be
found in [4, 5, 6, 7, 9].
In all previously known proofs of Kneser’s Theorem, the subgroup H depends
crucially on both sets S and T . With the goal of breaking this double depen-
dence in S and T , Balandraud investigated in recent work [1, 2] the properties
of a combinatorial poset that we now present.
Let S be a finite subset containing 1 of a group G. Following Balandraud, let
us define a cell of S as a finite subset X such that, for all z /∈ X, it holds that
zS 6⊂ XS. This notion is defined in [1, 2] and it is equivalent to the notion of
Université Pierre et Marie Curie, Paris 6, Combinatoire et Optimisation - case 189, 4
place Jussieu, 75252 Paris Cedex 05. [email protected]
Universitat Politècnica de Catalunya, Matemàtica Aplicada IV, Campus Nord - Edif. C3,
C. Jordi Girona, 1-3, 08034 Barcelona, Spain. [email protected]
Université de Bordeaux 1, Institut de Mathématiques de Bordeaux, 351 cours de la
Libération, 33405 Talence. [email protected]
http://arxiv.org/abs/0704.0382v1
nonextendible subset used in [3]. Throughout the paper, by a cell we always
mean a cell of S.
A cell X is called a u-cell if |XS| − |X| = u. A u-cell with minimal cardinality
is called a u-kernel (of S).
Balandraud showed that, for a finite set S in an abelian group G, in the poset
of j–cells containing the unity ordered by inclusion with 1 ≤ j ≤ |S| − 2, the
set of kernels form a chain of subgroups. Moreover, if there exists a u–cell, then
there is a unique u–kernel containing the unit element which is contained in all
u–cells containing the unit element.
One of the consequences of this work is a new proof and the following strength-
ening of Kneser’s Theorem:
Theorem 1 (Balandraud) For any non-empty finite subset S of an abelian
group G, there exists a finite subgroup H of G such that for any finite subset T
of G one of the following conditions hold :
• |TS| ≥ |T |+ |S| − 1
• HTS = TS and |TS| ≤ |HS|+ |HT | − |H|
As far as the authors are aware this is a surprising and strong formulation that
was not observed before and does not follow straightforwardly from the classical
forms of Kneser’s Theorem.
The purpose of the present note is to give a short proof for the nonabelian
case that, in the poset of j–cells that are subgroups ordered by inclusion with
0 ≤ j ≤ |S| − 1, the set of kernels form a chain of subgroups. Moreover, each
u-kernel of this poset is unique and contained in all u–cells of this poset.
From this statement Kneser’s theorem allows one to deduce Balandraud’s re-
sults for the abelian case, and in particular Theorem 1. Kneser’s Theorem has
several equivalent forms. We use the following one; see e.g [4, 7]:
Theorem 2 (Kneser [5]) Let G be an abelian group and X,Y ⊂ G be finite
subsets such that |XY | ≤ |X| + |Y | − 2. Then
|XY | = |HX|+ |HY | − |H|,
where H = stab(XY ) = {x : xXY = XY }.
Our main tool is the following Theorem of Olson[8, Theorem 2]. We give an
equivalent formulation here where we use left–cosets instead of right–cosets.
Theorem 3 (Olson [8]) Let X,Y be finite subsets of a group G, and let H
and K be subgroups such that HX = X, KY = Y and KX 6= X, HY 6= Y .
|X \ Y |+ |Y \X| ≥ |H|+ |K| − 2|H ∩K|.
In particular either |X \ Y | ≥ |H| − |H ∩K| or |Y \X| ≥ |K| − |H ∩K|.
We shall use the following lemma.
Lemma 4 ([1, 2]) Let G be a group and 1 ∈ S ⊂ G be a finite subset. Then
the intersection of two cells M1,M2 of S is a cell of S.
Proof. Let x /∈ M1 ∩M2. There is i with x /∈ Mi. Then xS 6⊂ MiS. Hence
xS 6⊂ (M1 ∩M2)S.
We can now state our main result, namely Theorem 5 below.
2 An application of Olson’s Theorem
Balandraud [1, 2] proved that, in the abelian case, the set of kernels containing
the unit element and ordered by inclusion is a chain of subgroups. In the non
abelian case we can prove only that the set of kernels that are subgroups forms
a chain. The abelian case can then be easily recovered, since Kneser’s Theorem
implies (as we shall see below) that a kernel containing the unit element is a
subgroup.
Theorem 5 Let S be a finite subset containing 1 of a group G. Let M be a
u–kernel of S which is a subgroup. Let N be a subgroup which is a v–cell and
suppose u, v ≤ |S| − 1.
(i) If either N is a v–kernel or u = v then M ⊂ N or N ⊂ M .
(ii) If N is a v–kernel and v ≤ u then M ⊂ N .
Proof. Suppose that M 6⊂ N and N 6⊂ M . Note that, since M is a cell, if
NMS = MS then NM = M , thus N ⊂ M against our assumption. Hence we
may assume NMS 6= MS and similarly MNS 6= NS. By Theorem 3 we have
one of the two following cases.
Case 1: |MS| − |(MS) ∩ (NS)| = |(MS) \ (NS)| ≥ |M | − |M ∩N |. It follows
that |(M ∩N)S| − |M ∩N | ≤ |(MS)∩ (NS)| − |M ∩N | ≤ |MS| − |M |. On the
other hand we have u = |MS| − |M | < |S| ≤ |(M ∩N)S|. Since |MS| − |M | is
a multiple of |M ∩N | we have
u = |MS| − |M | = |(M ∩N)S| − |M ∩N |.
By Lemma 4, M ∩N is a cell. Since M is a u–kernel, we have M ∩N = M, a
contradiction.
Case 2: |NS| − |(NS) ∩ (MS)| = |(NS) \ (MS)| ≥ |N | − |N ∩M |. It follows
that |(N ∩ M)S| − |N ∩ M | ≤ |(NS) ∩ (MS)| − |N ∩ M | ≤ |NS| − |N |. On
the other hand we have |NS| − |N | < |S| ≤ |(N ∩M)S|. Since |NS| − |N | is a
multiple of |N ∩M | we have
|NS| − |N | = |(N ∩M)S| − |N ∩M |. (1)
Assume first u = v. Then u = |MS|−|M | = |NS|−|N | = |(N∩M)S|−|N∩M |.
Since M is a u–kernel, we have M ∩N = M, a contradiction.
Assume that N is a v–kernel. Then (1) implies N ∩M = N, a contradiction.
This proves (i).
Assume now that v ≤ u. Suppose M 6⊂ N . By (i) we have N ⊂ M , which
implies in particular that |MS| − |M | is a multiple of N . Therefore, from
u = |MS| − |M | < |S| ≤ |NS| we have u = |MS| − |M | ≤ |NS| − |N | = v
which gives u = v. But then M 6⊂ N and N ⊂ M imply |N | < |M |, and since
N is now a u–cell, this contradicts M being a u–kernel.
We can now deduce Balandraud’s description for kernels and cells :
Corollary 6 (Balandraud [1, 2]) Let G be an abelian group and S ⊂ G be
a finite subset with 1 ∈ S. Let M be a u–kernel of S containing 1 with 1 ≤ u ≤
|S| − 2. Then,
(i) M is a subgroup.
(ii) Each u-cell is M–periodic.
(iii) Each v–kernel with u < v ≤ |S| − 2 is a proper subgroup of M .
Proof. Let X be a u-cell with u ≤ |S| − 2. By Kneser’s Theorem, the
inequality |XS| − |X| = u ≤ |S| − 2 implies
u = |XS| − |HX| = |HS| − |H|, (2)
where H is the stabilizer of XS. Since X is a cell and HXS = XS, we
have X = HX. Note that, since G is abelian, ({y} ∪ H)S = HS implies
y ∈ Stab(HS) ⊂ Stab(XS), so that y ∈ H. This observation and (2) imply
that H is an u–cell. In particular, by taking X = M , the period K of MS is a
u–cell. Since KMS = MS and M is a u–cell, we have K ⊂ KM ⊂ M . Since
M is a u–kernel we have M = K. This proves (i).
Now let H be the stabilizer of XS, where X is a u–cell. As shown in the
preceding paragraph H is also a u–cell. By Theorem 5 we have M ⊂ H and
thus MH = H. Since X is a cell and HXS = XS, we have X = HX = MHX.
Hence X ⊂ MX ⊂ MHX = X implies X = MX. This proves (ii).
Finally, by (i), a v–kernel N is a subgroup. By Theorem 5 we have N ⊂ M .
From Corollary 6, one can deduce Theorem 1.
Proof of Theorem 1: We may assume without loss of generality that 1 ∈ S.
Case 1: There is no m–cell for any 1 ≤ m ≤ |S| − 2.
• either we have |TS| ≥ |S|+ |T | − 1 for any non-empty finite T , in which
case the theorem clearly holds with H = {1}.
• or there exists some non-empty finite T such that |TS| ≤ |S| + |T | − 2.
Without loss of generality, we may also suppose 1 ∈ T . Now T must be
contained in an m–cell with m ≤ |S| − 2, but since no such cell exists for
1 ≤ m, we have that T itself must be a cell (a 0-cell) i.e. |TS| = |T |. We
therefore have HT = TH = T = TS = HTS where H is the (necessarily
finite) subgroup generated by S. We have just proved that the theorem
holds in this case with H = 〈S〉.
Case 2: There exists an m–cell with 1 ≤ m ≤ |S| − 2. We may therefore
consider the largest integer u ≤ |S| − 2 for which S admits a u–cell. Let H be
the u–kernel containing 1. Note that u ≤ |S|−2 implies that H is different from
{1}. Now let T be any finite non-empty subset such that |TS| − |T | ≤ |S| − 2.
We shall prove that HTS = TS.
By adding elements to T as long as necessary, we can find a cell X that contains
T and such that XS = TS. Note that we then have |XS| − |X| ≤ |TS| − |T | ≤
|S| − 2, so that X is a v–cell for some v ≤ u. By Corollary 6 (ii) we have
TS = XS = MXS = MTS where M is the v-kernel containing 1. By part (i)
of Corollary 6, H is a subgroup of M so that TS = XS = HTS as well.
Finally, |ST | ≤ |HS|+ |HT | − |H| follows from |ST | being a multiple of |H|.
References
[1] E. Balandraud, Une variante de la méthode isopérimetrique de Hamidoune,
appliquée au theoreme de Kneser, Preprint, december 2005.
[2] E. Balandraud, Quelques résultats combinatoires en Théorie Additive des
Nombres, Thèse de Doctorat de l’Université de Bordeaux I, May 2006.
[3] D. Grynkiewicz, A step beyond Kempermann structure Theorem, Preprint
May 2006.
[4] J. H. B. Kemperman, On small sumsets in Abelian groups, Acta Math. 103
(1960), 66–88.
[5] M. Kneser, Summenmengen in lokalkompakten abelesche Gruppen, Math.
Zeit. 66 (1956), 88–110.
[6] H.B. Mann, Addition Theorems, R.E. Krieger, New York, 1976.
[7] M. B. Nathanson, Additive Number Theory. Inverse problems and the ge-
ometry of sumsets, Grad. Texts in Math. 165, Springer, 1996.
[8] J.E. Olson, On the symmetric difference of two sets in a group. European
J. Combin. 7 (1986), no. 1, 43–54.
[9] T. Tao and V.H. Vu, Additive Combinatorics, Cambridge Studies in Ad-
vanced Mathematics 105 (2006), Cambridge University Press.
Introduction
An application of Olson's Theorem
|
0704.0383 | The Exact Boundary Condition to Solve the Schrodinger Equation of Many
Electron System | Microsoft Word - arxiv_prasadtext.doc
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 1 of 21)
The Exact Boundary Condition to Solve the Schrödinger
Equation of Many Electron System
Rajendra Prasad
Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India
E-mail: [email protected]
In an attempt to bypass the sign problem in quantum Monte Carlo simulation of
electronic systems within the framework of fixed node approach, we derive the exclusion
principle “Two electrons can’t be at the same external isopotential surface
simultaneously” using the first postulate of quantum mechanics. We propose the exact
Coulomb-Exchange nodal surface i.e. the exact boundary condition to solve the non-
relativistic Schrödinger equation for the non-degenerate ground state of atoms and
molecules. This boundary condition was applied to compute the ground state energies of
N, Ne, Li2, Be2, B2, C2, N2, O2, F2, and H2O systems using diffusion Monte Carlo
method. The ground state energies thus obtained agree well with the exact estimate of
non-relativistic energies.
INTRODUCTION
An ideal target of a quantum chemist/physicist is to solve the non-relativistic
Schrödinger equation exactly as it describes much of the world of chemistry. If we can
solve this equation at a realistic cost, we can make very precise predictions. At present,
only the full-CI method is available for obtaining the exact wave function within a given
basis set, but this method is too demanding computationally and therefore not affordable
even for a small system.
In recent years increasing attention has been drawn to the random walk approach
called diffusion Monte Carlo (DMC) method1 2 3 4 for solving Schrödinger equation. The
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 2 of 21)
attractiveness of DMC method lies in that it can treat many body problems exactly. The
DMC method is a projection method based on the combination of the imaginary time
Schrödinger equation, generalized stochastic diffusion process, and Monte Carlo
integration. The solution, it yields has only statistical error, which can be properly
estimated and in principle, made as small as desired. Since in the DMC method the wave
function has to be a population density, therefore, the DMC method can only describe the
constant sign solution of the Schrödinger equation. This poses a serious problem if one is
interested in the solution of a many electron system where the wave function is known to
be antisymmetric (i.e. both positive and negative) with respect to interchange of two
electrons. This situation is known as fermion sign problem in the quantum Monte Carlo
literature1-4. The solution of this problem is one of the most outstanding in all of the
computational physics/chemistry. This problem is often (mis)understood as a technical
detail that defeating the numerical simulators. To the best of our knowledge no
methodology is available to handle this problem in a systematic and controlled fashion.
However, we think that it is essentially a problem of exact boundary, which is not known
for many electron systems for obtaining well-behaved solutions of non-relativistic
Schrödinger equation. It is our understanding that the boundary must be derived from the
link between the formal mathematics and the physics of the real world.
In this article, we will derive the boundary condition for atomic and molecular
systems to obtain well-behaved solutions (i.e. bound state solution is single valued,
continuous, quadratically integrable, and differentiable) of non-relativistic electronic
Schrödinger equation. To start with, we are dealing with situations in which the ground
state is non-degenerate only.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 3 of 21)
THE EXACT BOUNDARY CONDITION
We have the time independent Schrödinger equation:
ˆ Ψ=Ψ EH …………..(1)
where Ĥ is the time independent non-relativistic electronic Hamiltonian operator in the
Born-Oppenheimer approximation, E0 is the eigenvalue of the full many-electron ground
state 0Ψ . The Ĥ is defined in atomic units as follows:
+−∇−=
electrons
electrons
electrons
2 1)(
1ˆ r ………..(2)
where the external potential, ∑=
Nuclei
rV )(
, ………….(3)
2∇ is Laplacian, ZI denotes nuclear charge, and rIi and rij symbolize the electron-nucleus
and electron-electron distance, respectively.
Following Hohenberg-Kohn theorem I, a proof only of existence5, the electron
density )(0 r
ρ in the ground state 0Ψ is a functional of )(rV
, i.e.
)(0 r
ρ = )]([0 rV
ρ . ………………………. (4)
Further, the full many electron ground state 0Ψ is unique functional of )(0 r
ρ , i.e.
0Ψ = )]([ 00 r
ρΨ . ……………………(5)
Evidently we can say that 0Ψ is a functional of )(rV
i.e.
0Ψ = )]([0 rV
Ψ . …………………. (6)
We have a choice to express the exact density:
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 4 of 21)
)]([0 rV
ρ = ∑
2 )]([
φ ………………..(7)
Where N denotes number of electrons. The functionals, { }NirVi .....1)],([ =
φ are exact
ortho-normal one electron functions of the function )(rV
, which give exact )]([0 rV
(Caution to reader!! At the moment, here is nothing to do with so-called s, p, d, f, ..etc.
type orbitals . The functionals { }NirVi .....1)],([ =
φ are entirely different from those
orbitals obtained from Kohn-Sham6 or similar formalisms.)
Now we can write the exact N electron ground state wave function as a functional of N
exact one electron functionals { }NirVi .....1)],([ =
0Ψ = )]]([)],......,([)],([[
10 NN rVrVrV
φφφΨ ……(8)
or 0Ψ = )]]([)],......,([)],([[ 22110 NN rVrVrV
φφφΨ ……..(9)
Since each one electron functional in { }NirVi .....1)],([ =
φ is a function of external
potential )(rV
, we can also write the exact N electron ground state wave function in
functional form as follows:
0Ψ = )](),......,(),([ 210 NrVrVrV
Ψ …………………...(10)
Thus the exact N electron non-degenerate ground state wave function is a unique
functional of external potential experienced by each electron i.e. functional of
)(),......,(),( 21 NrVrVrV
So far, it is not clear:
• Whether the wave function is symmetric or antisymmetric with respect to
interchange of any two electrons.
• What are the analytical forms of { }NirVi .....1)],([ =
• What is the analytical form of the exact wave function?
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 5 of 21)
• How to get the exact wave function from the exact density.
However, we get some idea about the topology of a well behaved ground state non-
degenerate wave function and distribution functions in a given external potential )(rV
In particular: “The probability of n-electrons (where n = 2..N) being found
simultaneously on the isopotential surface of an external potential )(rV
is same
irrespective of positions of the electrons on the surface.”
Now we proceed to decide the nature (symmetric or antisymmetric with respect to
interchange of any two electrons) of a well behaved many electron wave function.
Defining the local energy, EL:
∑∑ ∑ ∑
electrons
electrons
electrons
Nuclei
11 10
………….(11)
The terms IiI rZ and ijr1 in the equation (11) will blow up if 0→Iir and 0→ijr unless
so-called cusp conditions are obeyed by 0Ψ . The 0Ψ is exact and obeys electron nucleus
(e-N) and electron-electron (e-e) cusp conditions.
The wave function for a system of N identical particles must be symmetric or
antisymmetric with regard to interchange of any two of the identical particles, i and j.
Since the N particles are all identical, we could not have the wave function symmetric
with regard to some interchanges and antisymmetric with regard to other interchanges.
Thus the wave function of N identical particles must be either symmetric or
antisymmetric with regard to every possible interchange of any two particles.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 6 of 21)
Let us assume 0Ψ is symmetric with regard to interchange of electrons i and j.
There is a cusp in 0Ψ at rij = 0. This implies that 0Ψ is not differentiable at rij = 0.
Therefore, 0Ψ , symmetric with respect to interchange of any two electrons is not a well-
behaved solution. To make 0Ψ a well-behaved wave function, 0Ψ must be zero when rij
= 0 and also it must change sign with respect to the interchange of two electrons, i.e. if
ji rr
= then 0Ψ = 0. This condition is universal and independent of kind of external
potential. However, we are interested in a well-behaved solution of a bound state in a
given external potential )(rV
. From the previous argument, we know that the
simultaneous probability of finding two electrons is same everywhere at the isopotential
surface. Therefore, if )( irV
- )( jrV
= 0 then 0Ψ = 0.
Extending to N electron system, we have
If ( ) 0)()(
=−= Π
rVrVf
then 0Ψ = 0.
We can also express f as Vandermonde determinant:
)()(....)()(
)()(....)()(
)()(....)()(
11....11
rVrVrVrV
rVrVrVrV
rVrVrVrV
rVrVf
=−= Π …(12)
Consequently we have exclusion principle in the following form:
“Two electrons can’t be at the same external isopotential surface simultaneously.”
We see that if we are interested in a well behaved solution of the time
independent Schrödinger equation, the boundary condition (12) (i.e. antisymmetric wave
function) is obtained naturally due to singularity in the e-e interaction potential, which
respects Pauli’s exclusion principle. If electrons i and j are of opposite spin then we say
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 7 of 21)
that )( irV
- )( jrV
= 0 represents Coulomb (nodal) surface. If electrons i and j are of same
spin then )( irV
- )( jrV
= 0 represents Coulomb-Exchange nodal surface. All together,
the ( ) 0)()(
=−= Π
rVrVf
represents the Coulomb-Exchange nodal surface of N
electron system. Hereafter, we will call f as ExchangeCoulombf − nodal surface. However, the
solution obtained for the Hamiltonian (2) within the boundary ExchangeCoulombf − =0 does not
tell us about the spin multiplicity of the N electron system.
Further, we can rewrite the functional f in terms of Hermite polynomials, )]([ rVH k
)]([)]([....)]([)]([
)]([)]([....)]([)]([
)]([)]([....)]([)]([
)]([)]([....)]([)]([
1112111
2122212
1112111
0102010
NNNNNN
rVHrVHrVHrVH
rVHrVHrVHrVH
rVHrVHrVHrVH
rVHrVHrVHrVH
−−−−−
= …..(13)
In particular, if we multiply an optimizable one electron functional )]([ rVψ to the
equation (13) and we obtain an N electron wave function:
)]([)]([)]......([)]([)]([ 21 rVfrVrVrVNormrV N
rrrrr
ψψψ=Ψ ……………………………(14)
The one-electron density functional corresponds to the wave function (14):
′′=′ ∑
1111 )]([)]([)]([)]([)]();([
kkk rVHrVHArVrVrVrV
rrrrrr
ψψρ ……………….(15)
where Ak is normalization constant of )]([)]([ rVHrV k
The two-electron density functional in terms of one-electron density functional:
)]();([)]();([)]();([)]();([
)](),();(),([
21122211
rVrVrVrVrVrVrVrV
rVrVrVrV
′′−′′=
rrrrrrrr
……(16)
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 8 of 21)
Here it appears that we can obtain exact ground state energy by optimizing only a one-
electron functional )]([ rVψ in the equation (14).
A very interesting and new physics obtained from the equation (13) is that each
row in the determinant represents different level (k) of Kamalpur breathing (anharmonic
quantum breathing) of electron cloud in a given external potential )(rV
and each level, k
is occupied by one electron (the elementary particle).
BYPASSING THE SIGN PROBLEM
We can bypass the fermion sign problem in the electronic structure diffusion
Monte Carlo (DMC) method using fixed node approach. Here one assumes a prior
knowledge of the nodal surface i.e.
0(R) = 0. Due to tiling property
7 of the exact ground
state wave function, the Schrödinger equation is solved in the volume embraced by the
nodal surface, where the wave function has a constant sign and in this way the fermion
sign problem is bypassed. The exact knowledge of Coulomb Exchange nodal surface
allows us for an exact stochastic solution of the Schrödinger equation. The restriction in
the random walk RR ′→ during the electronic structure diffusion Monte Carlo
simulation is as follows:
reject
accept
RfRf ExchangeCoulombExchangeCoulomb
)()( ………….(17)
We have applied the boundary condition (17) for the ground state electronic
structure diffusion Monte Carlo simulation of N, Ne, Li2, Be2, B2, C2, N2, O2, F2, and
H2O systems.
Monte Carlo calculations can be carried out using sets of random points picked
from any arbitrary probability distribution. The choice of distribution obviously makes a
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 9 of 21)
difference to the efficiency of the method. If the Monte Carlo calculations are carried out
using uniform probability distributions, very poor estimates of high-dimensional integrals
are obtained, which is not a useful method of approximation. These problems are handled
by introducing the importance sampling approach8 9. In this approach the sampling
points are chosen from a trial distribution, which concentrates on points where the trial
function, ΦT(R) is large.
In the present DMC calculations, we have chosen the trial function, ΦT in the
form:
ΦT = Φ.F , ….(18)
where Φ denotes the Hartree Fock (HF) or multi configuration self consistent field
(MCSCF) wave function and F is a correlation function that depends on inter-particle
distances. The HF and MCSCF wave functions were obtained using the GAMESS
package10 and employed Dunning’s cc-VTZ atomic basis set 11. In order to satisfy the
electron nucleus (e-N) cusp condition, all Gaussian type s basis functions were replaced
with eight Slater-type s basis functions. The exponents of Slater-type s functions were
taken from Koga et al. 12 and satisfy the e-N cusp condition.
We have chosen the Schmidt, Moskowitz, Boys, and Handy (SMBH) correlation
function FSMBH
13. For the SMBH correlation function, Eqn. (19), we have included terms
up to 2nd order, where order, s is defined as s = l + m + n.
( )
+−= ∑∑ ∑
µµµµµ
atoms
electrons
iAASMBH rrrrrcF exp ………….(19)
where
………….(20)
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 10 of 21)
and r denotes inter-particle distance. Six non-redundant parameters out of the total ten
were optimized keeping b = 1.0 as follows:
1) First we obtained optimal parameters by minimizing the energy and variance at
the variational Monte Carlo (VMC) level.
2) Using this VMC optimal trial function, the trial function fixed node DMC
calculation was carried out and walkers were collected after each 2000 steps.
Further, correlation parameters were reoptimized to minimize the variance with
~100,000 walkers. Here reference energy was set to the trial function fixed node
DMC energy.
These optimized trial functions were used for importance sampling in the DMC
simulation and a random walk RR ′→ was accepted if
0)()( >′−− RfRf ExchangeCoulombExchangeCoulomb .
The DMC calculations were performed using the open source quantum Monte
Carlo program, ZORI14. Around 10,000 walkers were used for the systems studied. The
Umrigar et al.15 algorithm was used for DMC walks and Caffarel Assaraf et al.16
algorithm for population control. We have allowed only one electron walk at a time. The
DMC calculations were done at several time steps. We report only those energies
extrapolated to zero time step.
We present the ground state DMC energies of N, Ne, Li2, Be2, B2, C2, N2, O2, F2,
and H2O systems in Table I. The DMC energies obtained using our newly derived
boundary ExchangeCoulombf − = 0 are far better than the trial function fixed node DMC
energies17 and compare well with the experimental counterpart. However, present
simulations were noisy and unpleasant compared to conventional trial function fixed
node DMC simulations. It is worth noting that we have obtained DMC energy even
lower than the exact value at smaller time steps for the atoms of relatively larger atomic
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 11 of 21)
radius perhaps due to failure of the distributions to reach the steady state or equilibrium
distributions in a finite number of steps. This problem can be handled in the Green’s
function quantum Monte Carlo (GFQMC) method as it takes the advantage of the
properties of Green’s functions in eliminating time-step entirely in treating the steady
state equation. The GFQMC is well suited if boundaries are exactly known18. If the trial
function boundary and the ExchangeCoulombf − = 0 does not coincide and also non-zero values
of trial function are very much different from the exact solution, which could lead to
large statistical fluctuations from poor sampling and possibly to an effective non-ergodic
diffusion process due to the finite projection time in practical calculations. Therefore, we
are looking for an alternative well behaved trial function whose boundary coincides with
those of ExchangeCoulombf − .
CONCLUSION
This article presents a progress of the author's research in order to get exact
solution of non-relativistic Schrödinger equation of many electron systems. A conclusion
of this on going research is that we have derived the exclusion principle “Two electrons
can’t be at the same external isopotential surface simultaneously” using the first postulate
of quantum mechanics. We propose the exact Coulomb-Exchange nodal surface i.e. the
exact boundary to solve the non-relativistic Schrödinger equation for non-degenerate
ground state of atoms and molecules. Using this newly derived boundary condition, one
can bypass the fermion sign problem in the electronic structure Quantum Monte Carlo
simulation and hence the exact ground state energy as well as the exact electron density.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 12 of 21)
REFERENCES
1. Anderson, J. B. A random-walk simulation of the Schrödinger equation: +3H . J. Chem.
Phys. 63, 1499-1503 (1975)
2. Ceperley, D. M. & Mitas, L. in New Methods in Quantum Mechanics, I. Prigogine, S.
A. Rice, Eds.,(John Wiley and Sons, New York, 1996), Vol. 93.
3. Hammond, B. L., Lester, Jr, W. A. & Reynolds, P. J. Monte Carlo Methods in Ab
Initio Quantum Chemistry; World Scientific: Singapore, 1994.
4. Ceperley, D. & Alder, B. Quantum Monte Carlo. Science 231, 555-560 (1986)
5. Hohenberg, P. & Kohn, W. Inhomogeneous Electron Gas. Phys. Rev. 136, B864-
B871 (1964)
6. Kohn, W. & Sham, L. J. Self-Consistent Equations Including Exchange and
Correlation Effects. Phys. Rev. 140, A1133-A1138 (1965)
7. Cepereley, D. M. Fermion nodes. J. Stat. Phys. 63, 1237-1267 (1991)
8. Metropolis, N. A., Rosenbluth, W., Rosenbluth, M. N., Teller, A. H. & Teller, E.
Equation of State Calculations by Fast Computing Machines. J. Chem. Phys. 21, 1087-
1092 (1953)
9. Grimm, R. C. & Storer, R. G. Monte-Carlo solution of Schrödinger's equation. J.
Comput. Phys. 7, 134-156 (1971)
10. Schmidt, M. W. et al. General atomic and molecular electronic structure system. J
Comput Chem 14, 1347-1363 (1993)
11. Dunning, Jr., T. H. Gaussian basis sets for use in correlated molecular calculations. I.
The atoms boron through neon and hydrogen. J. Chem. Phys. 90, 1007-1023 (1989)
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 13 of 21)
12. Koga, T., Kanayama, K., Watanabe, S. & Thakkar, A. J. Analytical Hartree-Fock
wave functions subject to cusp and asymptotic constraints: He to Xe, Li+ to Cs+, H- to I-.
Int. J. Quantum Chem. 71, 491-497 (1999)
13. Schmidt, K. E. & Moskowitz, J. W. Correlated Monte Carlo wave functions for the
atoms He through Ne. J. Chem. Phys. 93, 4172-4178 (1990)
14. Aspuru-Guzik , A. et al. Zori 1.0: A parallel quantum Monte Carlo electronic
structure package. J. Comp. Chem. 26, 856-862 (2005)
15. Umrigar, C. J., Nightingale, M.P. & Runge, K.J. A diffusion Monte Carlo algorithm
with very small time-step errors. J. Chem. Phys. 99,2865-2890 (1993)
16. Assaraf, R., Caffarel, M. & Khelif, A. Diffusion Monte Carlo methods with a fixed
number of walkers. Phys. Rev. E. 61, 4566-4575 (2000)
17. Filippi, C. & Umrigar, C. J. Multiconfiguration wave functions for quantum Monte
Carlo calculations of first-row diatomic molecules. J. Chem. Phys. 105, 213-226 (1996)
18. Kalos, M. H., Monte Carlo Calculations of the Ground State of Three- and Four-
Body Nuclei. Phys. Rev. 128, 1791-1795 (1962)
ACKNOWLEDGEMENTS
The QMC calculations were carried out at the Lawrence Berkeley National Laboratory,
Berkeley. The author gratefully acknowledges Professor W. A. Lester for his support
during the stay at Berkeley. The author is indebted to Professor P. Chandra of Banaras
Hindu University, Varanasi for his interest and helpful discussion during the preparation
of the manuscript. Professor S. K. Sengupta of Banaras Hindu University, Varanasi is
acknowledged for careful reading of the manuscript.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 14 of 21)
TABLE I. The total ground state energies obtained from fixed-node DMC calculation.
Atom /
Molecule
Bond
length
CSF,D ETFN-DMC
(Ref. 17)
ECEN-DMC
(Extrapolated to τ=0) E0
-, 111
-54.5753(3)
-54.5841(5)
-54.5902(11)
-54.5892
Ne 1,1 -128.9216(15) -128.938(1) -128.9375
Li2 5.051
-14.9911(1)
-14.9938(1)
-14.9955(5)
-14.9954
Be2 4.63
5,16
-29.3176(4)
-29.3301(2)
-29.3378(15)
-29.33854(5)
B2 3.005
6,11
-49.3778(8)
-49.3979(6)
-49.41655(45)
-49.415(2)
C2 2.3481
4,16
77, 314
-75.8613(8)
-75.8901(7)
-75.9035(9)
-75.9229(19)
-75.923(5)
2.068
4,17
-, 552
-109.487(1)
-109.505(1)
-109.520(3)
-109.5424(15)
-109.5423
O2 2.282
-150.268(1)
-150.277(1)
-150.3274(15)
-150.3268
F2 2.68
-199.478(2)
-199.487(1)
-199.5289(25)
-199.5299
-199.52891(4)
H2O
-, 300
-76.4175(4)
-76.429(1)
-76.4376(11)
-76.438(3)
-76.4376
Bond lengths and energies are in atomic units. In the third column, we list the number of
configuration state functions (CSFs) and number of determinants (D) in the trial function
(ΦT). ETFN-DMC denotes the DMC energy with ΦT =0. ECEN-DMC denotes the DMC energy
with ExchangeCoulombf − = 0. E0 denotes the exact, non-relativistic, infinite nuclear mass
energy. The numbers shown in bracket are error bar.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 15 of 21)
Supplementary note for reviewers:
1. The proposed theory is to deal only real interacting many electron systems. To
start with only non-degenerate ground state of many electron systems are
considered. Author is neither interested nor intended to deal any kind of non-
interacting systems such as free fermion, free electron gases, or free particles
because author think that none of the real system belong to either of these classes.
Author has chosen to construct the boundary condition from the link between the
formal mathematics and the physics of many electron systems.
2. Difference between spatial nodes and Coulomb-Exchange nodal surface:
I hope that people can distinguish spatial nodes and Coulomb Exchange nodes
and the physics behind the different kind of nodes. Whatever I have discussed in
this paper is only about Coulomb-Exchange nodal surfaces. There is no analogy
with a particle in a box node and Coulomb-Exchange nodal surfaces. For
example: The function f(r1,r2)=(r1-1)(r2-1)(r1-r2)exp(-2 r1-2 r2) is antisymmetric
with respect to interchange of two electrons. However, the node (r1-1)(r2-1) is
symmetric with respect to interchange of two electrons and fixed and this node
can be compared with nodes of particle in a box. The Coulomb-Exchange node
(r1-r2) is antisymmetric with respect to interchange of two electrons and
responsible for removal of singularity in e-e interaction potential. The Coulomb-
Exchange nodal surfaces only occur in a system of more than one electron
systems. Author understands that the Coulomb-Exchange nodal surfaces are
directly responsible for the existence of real many electron systems.
3. A consequence of proposed solution of the sign problem is that the ground
state of Helium atom in the non-relativistic limit has a nodal surface.
However, it is understood that He ground state wave function is symmetric
and has no such node.
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 16 of 21)
A consequence of proposed solution for the sign problem is that the ground state
of the Helium atom in the non-relativistic limit has Coulomb-Exchange nodal
surface, r1-r2=0.
An understanding that He atom ground state wave function is symmetric and has
no such node is an illusion only. This illusion arises due to a practice that the
QMC people using phi(1)*phi(2)*Jastrow trial function, where phi(r) is 1s orbital.
The trial function is symmetric with respect to exchange of two electrons. The
trial function also satisfies electron nucleus cusp condition. We also expect that
the final solution will satisfy e-e cusp condition. Since the trial function is
symmetric, people got accurate energy and assumed that the final solution is also
symmetric and it does not have any node also wave function is non-zero at the
point of coincidence of two electrons. Where is Coulomb hole? However, it can
be proven that a symmetric solution is not acceptable. Proofs are as follows:
“Proof for the existence of Coulomb-Exchange node in He ground
state exact wave function”
A.) Let assume
),( 21 rrsym
Ψ is an exact symmetric wave function.
i.e. ),( 21 rrsym
Ψ = ),( 12 rrsym
Since ),( 21 rrsym
Ψ is exact, it must satisfy the cusp condition at 21 rr
= . Clearly
there is a cusp at 21 rr
Since there is a cusp at 21 rr
= in ),( 21 rrsym
Ψ , the second derivative
2 ),( xrrsym ∂Ψ∂
is not defined at 21 rr
Therefore ),( 21 rrsym
Ψ is not a well behaved solution and hence it is not an
acceptable wave function. Only option left is antisymmetric solution.
B.) Another proof:
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 17 of 21)
),( 21 rr
Ψ = ∑∑
νµµν φφ )()( 21 rrc
………...(S-1)
Where { })(rrµφ forms a real infinite one-electron basis.
Since ),( 21 rr
Ψ is expanded over infinite basis set and hence it is exact.
This implies that ),( 21 rr
Ψ satisfies the cusp condition at 21 rr
),( 21
Ψ∇ =∑∑
νµµν φφ )()( 21
1 rrc
. ………...(S-2)
The second derivative ),( 21
Ψ∇ is continuous at 21 rr
= because each term in
the expansion is continuous (the rules of continuity for algebraic combinations).
This implies that there is no cusp in ),( 21 rr
Ψ at 21 rr
BUT ),( 21 rr
Ψ has to satisfy the cusp condition at 21 rr
This is only possible if ),( 21 rr
Ψ changes the sign at 21 rr
And hence the exact ),( 21 rr
Ψ has exchange node irrespective of its spin
multiplicity.
C.) More illustrative example:
Hamiltonian for He atom:
H +−−∇−∇−= ………...(S-3)
and )2,1()2,1( Ψ=Ψ EH ………...(S-4)
Let expand
νµµν φφ )2()1()2,1( c ………...(S-5)
Where { })(rµφ is complete set of eigen functions of the Hamiltonian
1ˆ 2 −∇−= with eigenvalue equation )()(ˆ rrh µµµ φεφ = .
Rewriting the He Hamiltonian:
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 18 of 21)
H +−−∇−∇−= =
hh ++ ………...(S-6)
)2,1(
)2()1()ˆˆ(
)2,1(
)2,1()
)2,1(
)2,1(ˆ
EL +Ψ
νµµν φφ
………...(S-7)
Since )(rµφ is an eigen function of ĥ .
We can write
)2,1(
)2()1()(
EL +Ψ
νµνµµν φφεε
………...(S-8)
µννµ εε dE ++=
µννµ εε dE −=+ ………...(S-9)
)2,1(
)2()1()(
EL +Ψ
νµµνµν φφ
………...(S-10)
)2,1(
)2()1(
)2,1(
)2()1(
EL +Ψ
∞ ∞∞ ∞
νµµνµν
νµµν φφφφ
………...(S-11)
)2,1(
)2()1(
)2,1(
)2,1(
EL +Ψ
νµµνµν φφ
………...(S-12)
)2,1(
)2()1(
EEL +Ψ
νµµνµν φφ
………...(S-13)
If )2,1(Ψ is exact then
EEL =
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 19 of 21)
and
)2,1(
)2()1(
νµµνµν φφ
………...(S-14)
Now let assume that the exact solution is symmetric with respect to interchange of
two electrons. )2,1(Ψ and )(rµφ are well behaved and differentiable. From the
rules of continuity for algebraic combinations,
the term in equation (S-14),
)2,1(
)2()1(
νµµνµν φφdc
is continuous and finite
and it should not diverge when 012 →r . Therefore, symmetric solution is not
acceptable.
However, if )2,1(Ψ = 0 at r12=0 then
)2,1(
)2()1(
νµµνµν φφdc
will also diverge
and can compensate the divergence in 1/r12 term.
Thus the only acceptable solution is antisymmetric (with respect to interchange
of two electrons) solution.
D.) Another example:
Almost all QMC people believe (their believe is based on some assumptions and
approximations) that He atom ground state wave function is symmetric. This is an
illusion. This can be understood as follows:
Let us take trial functions of two electron system:
( )221
1 xxb
xxg 21
211 ),(
−−=
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 20 of 21)
( ) eeexx Uxxxxg 21 22212 221),(
( ) eeexx Uxxxxg 21 22213 21),( −
−−=
If someone claims that He ground state is symmetric, what kind of exact
symmetric wave functions they are getting finally? The functions g1and g2 are
symmetric with respect to interchange of two electrons. The functions like g1, g2,
and g3 can satisfy cusp condition. The functions g1 and g2 are not differentiable at
x1=x2 and therefore these are not acceptable. The antisymmetric function g3 are
differentiable at x1=x2.
However, people have got very accurate ground state energy for He atom using
HF*Jastrow trial function and they concluded that He atom has no node without
examining the simultaneous probability of finding two electrons at exactly same
place. I think, they got good results due to inherent beauty of DMC technique.
The functions g2 is symmetric and g3 is antisymmetric with respect to interchange
of two electrons. However, g2*g2 and g3*g3 give the exactly same probability
distribution i.e. same physics. Functions g2 and g3 vanish when x1=x2. Further, the
VMC calculation for g2 and g3 will give exactly same energy. Can anyone predict
that the VMC energies obtained from g2 and g3 represent singlet or triplet state? I
am sure it is not possible.
An antisymmetric wave function can satisfy cusp condition as well as it’s
derivative will be continuous simultaneously at the point of coincidence. Here
symmetric and antisymmetric wave functions serve same distribution. Why I
should not prefer antisymmetric wave function for which a boundary condition
can be imposed easily?
4. Further, if I assume the argument “He ground state wave function is
symmetric and has no such node” is correct. The end will be a nonsense,
which is as follows:
____________________________________________________________________________________________________________
The Exact Boundary Condition to Solve the Schrödinger Equation of Many Electron System
By Rajendra Prasad, Village + Post: Kamalpur, District: Chandauli, Uttar Pradesh, PIN: 232106, India E-mail: [email protected]
4/3/2007 (Page 21 of 21)
It is very much common practice in QMC calculation to take Hartree-Fock trial
function as a product of alpha-beta determinants. For example N atom:
PSIT(1,2,3,4,5,6,7) = Det(1,2,3,4,5,6,7).
PSIT(1,2,3,4,5,6,7) = Detα (1,2,3,4,5)*Detβ(6,7).
PSIT(1,2,3,4,5,6,7) = Detα(1,2,3,4,5)*Detβ(6,7).
Detα(1,2,3,4,5)*Detβ(6,7) ≠ Detα(1,2,3,4,7)*Detβ(6,5)
The trial function PSIT is clearly neither symmetric nor antisymmetric with
respect to interchange of alpha-beta electrons. It is not clear to me what kind of
final solution (i.e. symmetric or antisymmetric) we will obtain with this trial
function fixed node DMC? The fact that we can not write PSIexact = Phiα*Phiβ.
5. It is natural to ask what is node of 3S He atom and why people are getting
very accurate energy with exchange node r1-r2=0?
At the moment, I can only say that this is due to artifact of importance sampling
because people have used HF*Jastrow trial function. The correlation energy for
He(3S) atom is around 2mH and the overlap of HF trial function with exact wave
function can be anticipated to be more than 99%. Perhaps due to technical reasons
final DMC solution may have converged to He(3S). I have seen some recent
papers on the node of He(3S) system. It is widely claimed that the node r1-r2=0
belongs to He(3S) system and it is exact. I differ with their argument and I proved
that the exchange node r1-r2=0 belongs to He ground state. I do not know if
anyone have performed DMC calculation with a trial function like psit(r1,r2)=(r1-
r2)*exp(-2*r1)*exp(-2*r2) and reported the energy for He(
3S). Anyway, at present
I am interested only in the non-degenerate ground state of atoms and molecules.
Author welcomes further comments, questions, and suggestions if any.
[email protected]
|
0704.0384 | Clustering features of $^9$Be, $^{14}$N, $^7$Be, and $^8$B nuclei in
relativistic fragmentation | Clustering features of 9Be, 14N, 7Be, and 8B nuclei in relativistic
fragmentation
D. A. Artemenkov,∗ T. V. Shchedrina, R. Stanoeva, and P. I. Zarubin†
Joint Insitute for Nuclear Research, Dubna, Russia
(Dated: November 10, 2021)
Abstract
Recent studies of clustering in light nuclei with an initial energy above 1 A GeV in nuclear treack
emulsion are overviewed. The results of investigations of the relativistic 9Be nuclei fragmentation
in emulsion, which entails the production of He fragments, are presented. It is shown that most
precise angular measurements provided by this technique play a crucial role in the restoration of the
excitation spectrum of the α particle sysytem. In peripheral interactions 9Be nuclei are dissociated
practically totally through the 0+ and 2+ states of the 8Be nucleus.
The results of investigations of the dissociation of a 14N nucleus of momentum 2.86 A GeV/c in
emulsion are presented as example of more complicated system. The momentum and correlation
characteristics of α particles for the 14N→3α +X channel in the laboratory system and the rest
systems of 3α particles were considered in detail.
Topology of charged fragments produced in peripheral relativistic dissociation of radioactive 8B,
7Be nuclei in emulsion is studied.
PACS numbers: 21.45.+v, 23.60+e, 25.10.+s
∗Electronic address: [email protected]
†Electronic address: [email protected]; URL: http://becquerel.lhe.jinr.ru
http://arxiv.org/abs/0704.0384v1
mailto:[email protected]
mailto:[email protected]
http://becquerel.lhe.jinr.ru
I. INTRODUCTION
The peripheral fragmentation of light relativistic nuclei can serve as a source of infor-
mation about their exitations above particle decay thresholds including many-body final
states. The interactions of this type are provoked either in electromagnetic and diffraction
processes, or in nucleon collisions at small overlapping of the colliding nucleus densities. A
fragmenting nucleus gains an excitation spectrum near the cluster dissociation thresholds.
In the kinetic region of fragmentation of a relativistic nucleus there are produced nuclear
fragment systems the total charge of it is close to the parent-nucleus charge. A relative
intensity of formation of fragments of various configurations makes it possible to estimate
the importance of different cluster modes.
The opening angle of the relativistic fragmentation cone is determined by the Fermi-
momenta of the nucleon clusters in a nucleus. Being normalized to the mass numbers
they are concentrated with a few percent dispersion near the normalized momentum of
the primary nucleus. When selecting events with dissociation of a projectile into a narrow
fragmentation cone we see that target-nucleus non-relativistic fragments either are absent
(“white”stars in Ref.[1]), or their number is insignificant. The target fragments are easily
separated from the fragments of a relativistic projectile since their fraction in the angular
relativistic fragmentation cone is small and they possess non-relativistic momentum values.
In the peripheral fragmentation of a relativistic nucleus with charge Z the ionization
induced by the fragments can decrease down to a factor Z, while the ionization per one
track – down to Z2. Therefore experiment should provide an adequate detection range. In
order to reconstruct an event, a complete kinematic information about the particles in the
relativistic fragmentation cone is needed which, e.g., allows one to calculate the invariant
mass of the system. The accuracy of its estimation decisively depends on the exactness of
the track angular resolution. To ensure the best angular resolution, it is necessary that the
detection of relativistic fragments should be performed with the best spacial resolution.
The nuclear emulsion technique, which underlies the BECQUEREL project at the JINR
Nuclotron [2], well satisfies the above-mentioned requirements. It is aimed at a systematic
search for peripheral fragmentation modes with statistical provision at a level of dozens
events, their classification and angular metrology. Emulsions provide the best spacial res-
olution (about 0.5 µm) which allows one to separate the charged particle tracks in the
FIG. 1: An event of the type of “white”star from the fragmentation of a relativistic 9Be nucleus into
two He fragments in emulsion. The photograph was obtained on the PAVIKOM(FIAN) complex.
three-dimensional image of an event within one-layer thickness (600 µm) and ensure a high
accuracy of angle measurements. The tracks of relativistic H and He nuclei are separated
by sight. As a rule, in the peripheral fragmentation of a light nucleus its charge can be
determined by the sum of the charges of relativistic fragments. Multiple-particle scattering
measurements on the light fragment tracks enable one to separate the H and He isotopes.
The analysis of the products of the relativistic fragmentation of neutron-deficient isotopes
has some additional advantages owing to a larger fraction of observable nucleons and min-
imal Coulomb distortions. Irradiation details and a special analysis of interactions in the
BR-2 emulsion are presented in Ref. [3, 4]. In what follows, we give the first results of
the study of the 9Be,8B, 7Be 14N nuclei fragmentation with a few A GeV energy which are
obtained with the use of a part of the material analyzed.
II. FRAGMENTATION OF 9BE NUCLEI
The 9Be nucleus is a loosely bound n+α+α system. The energy threshold of the
9Be→n+α+α dissociation channel is 1.57 MeV. The study of the 9Be fragmentation at
relativistic energies gives the possibility of observing the reaction fragments, which are the
decay products of unbound 8Be and 5He nuclei.
The method of nuclear emulsions used in the present paper allows one to observe the
charged component of the relativistic 9Be→2He+n fragmentation channel. Owing to a
good angular resolution of this method it is possible to separate the 9Be fragmentation
events, which accompanied by the production of an unstable 8Be nucleus with its subsequent
breakup to two a particles. In this case, the absence of a combinatorial background (of three
and more α particles) for 9Be, which is typical for heavier Nα nuclei 12C and 16O makes it
possible to observe distinctly this picture.
Nuclear emulsions were exposed to relativistic 9Be nuclei at the JINR Nuclotron. A
beam of relativistic 9Be nuclei was obtained in the 10B→9Be fragmentation reaction using a
polyethylene target. The 9Be nuclei constituted about 80% of the beam, the remaining 20%
fell on Li and He nuclei.[5]
Events were sought by microscope scanning over the emulsion plates. In total 362 events
of the 9Be fragmentation involving the two He fragment production in the forward fragmen-
tation cone within a polar angle of 6◦(0.1 rad) were found. The requirement of conservation
of the fragment charge in the fragmentation cone was fulfilled for the detected events. In
event selection 5 - 7 tracks of various types were allowed in a wide (larger than 6◦) cone to
increase statistics. An example of the 9Be→2He fragmentation event in emulsion is given in
Fig. 1 [2]. This event belongs to the class of “white”stars as far as it contains neither target
nucleus fragments, nor produced mesons. This event sample includes 144 “white ”stars. The
angles of the tracks in emulsion for the detected events were obtained using a fine measuring
microscope. Angular measurements for the 362 events were carried out with an accuracy
not worse then 4.5×10−3 rad.
In analyzing the data both He fragments observed in the 9Be→2He+n channel were
supposed to be a particles. This assumption is motivated by the fact that at small angles the
9Be→ 24He+n fragmentation channel with an energy threshold of 1.57 MeV must dominate
the 9Be→3He+4He+n channel whose energy threshold is 22.15 MeV. The 3He fraction will
not exceed a few percent in this energy range [6] and all the He fragments in the detected
events may be thought of as α particles.
In Fig. 2a the PT transverse momentum distribution of α particles in the laboratory
frame of reference is calculated without the account of particle energy losses in emulsion by
the equation
PT = p0 · A · sinθ (1)
where p0, A and θ are the momentum per nucleon, the fragment mass and the polar emission
angle, respectively. The outer contour corresponds to all events. The inner histogram is
obtained for events accompanied by protons recoil of emulsion target (dashed area). The
mean value of the transverse momentum for the total event sample in the laboratory system
is equal to < PT >≈103 MeV/c with FWHM σ ≈72 MeV/c. This may be an indication
of the fact that the experimental data are not of the same kind which can be pronounced
when going over to the c.m.s. of two α particles.
The P∗T transverse momentum distribution of α particles in the c.m.s. of two α particles
, MeV/cTP
0 50 100 150 200 250 300 350
, MeV/c
0 50 100 150 200 250 300 350
FIG. 2: The PT transverse momentum distribution of α particles in the laboratory system (a),
and the P∗T momentum distribution in the c.m.s. of an α particle pair (b). The outer contour
corresponds to all events. The inner histogram is obtained for events, which accompanied by
protons recoil of emulsion target (dashed area).
described by the equation
∼= PT i −
where PT i is the transverse momentum of an i-th α particle in the laboratory system nα=2
is given in Fig. 2b. There is observed a grouping of events around two peaks with the values
< P ∗T i >≈24 MeV/c and < P
T i >≈101 MeV/c. In Ref [7] the appropriate mean values of the
α fragment transverse momenta are < P ∗T i >≈121 MeV/c for
16O→4α,< P ∗T i >≈141 MeV/c
[8] for 12C→3α and < P ∗T i >≈200 MeV/ for
22Ne→5α (processing of the available data).
There by we clearly see a tendency toward an increase of the mean α particle momentum
with increasing their multiplicity. This implies a growth of the total Coulomb interaction of
alpha clusters arising in nuclei.
In the opening angle Θ distribution (Fig. 3) one can also see two peaks with mean values
4.6×10−3rad. and 26.8×10−3rad. The ratio of the numbers of the events in the peaks is
close to unity.
The Θ distribution entails the invariant energy Q2α distribution, which is calculated as a
difference between the effective invariant mass M2α of an α fragment pair and the doubled
α particle mass by the equations
M22α = −(
Q2α = M2α − 2 ·mα (3)
rad.-310×, Θ
0 10 20 30 40 50 60 70 80
FIG. 3: The opening Θ angle distribution of α particles in the 9Be→2α fragmentation reaction at
1.2 A GeV energy. The outer contour corresponds to all events. The inner histogram is obtained
for events, which accompanied by protons recoil of emulsion target (dashed area).
where Pj is the α particle 4-momentum.
In the invariant energy Q2α distribution (Fig. 4) there are two peaks in the ranges 0 to
1 MeV and 2 to 4 MeV. The shape of the distribution does not contradict the suggestion
about the 9Be fragmentation involving the production of an unstable 8Be nucleus which
decays in the 0+ and 2+ states. The values of the peaks of the invariant energy Q2α and
the transverse momenta P∗T in the c.m.s. relate to each other. To the Q2α range from 0 to
1 MeV with a peak at 100 keV there corresponds a peak P∗T with < P
T i >≈24 MeV/c , and
to the Q2α range from 2 to 4 MeV there corresponds a peak with < P
T i >≈101 MeV/c.
III. FRAGMENTATION OF 14N NUCLEI
A stack of layers of BR-2 emulsion was exposed to a beam of 14N nuclei accelerated
[9] to a momentum of 2.86 A GeV/c at the Nuclotron of the Laboratory of High Energy
Physics (JINR). Already been found amoung 950 inelastic events in which the total fragment
charge was equal to the Z0=7 fragment charge and there were no produced particles. Events
were sought by viewing over the track length which provided the accumulation of statistics
, MeVα2Q
0 2 4 6 8 10 12
, keV
0 100 200 300 400 500 600 700 800 900 10000
FIG. 4: The invariant energy Q2α distribution of α particle pairs in the
9Be→2α fragmentation
reaction at 1.2 A GeV energy. On the intersection: the Q2α range from 0 to 1 MeV. Arrows
mark the 8Be nucleus levels (0+ and 2+). The outer contour corresponds to all events. The inner
histogram is obtained for events, which accompanied by protons recoil of emulsion target (dashed
area).
without selection. The selected events are divided in two classes. The events of the type
of “white”star and the interactions involving the production of one or a few target-nucleus
fragments belong to the first class.
Table I shows the charge multi-fragmentation topology which was studied for the events
satisfying the above-mentioned conditions. The upper line is the Z>2 fragment charge, the
second line is the number of single-charged fragments , the third one the number of two-
charged fragments, and the fourth and fifth lines are the number of the detected events
with a given topology for “white”stars and events with target-nucleus excitation for each
channel, respectively. The two last lines present the total number of interactions calculated
in absolute values and in percent.
The analysis of the data of Table I shows that the number of channels involving Z>3
fragments for the “white”stars is larger by about a factor of 1.5 than that for the events
accompanied by a target breakup. On the contrary, for the 2+2+2+1 charge configuration
TABLE I: The charge topology distribution of the “white”stars and the interactions involving the
target-nucleus fragment production in the 14N dissociation at 2.86 A GeV/c momentum.
Zfr 6 5 5 4 3 3 – –
NZ=1 1 – 2 1 4 2 3 1
NZ=2 – 1 – 1 – 1 2 3
NW.S. 13 4 3 1 1 1 6 17
Nt.f. 15 1 3 3 – 2 5 32
NP 28 5 6 4 1 3 11 49
NP,% 26 5 5 4 1 3 10 46
channel this number is smaller by about a factor of 1.5. Thus, in the events with target
breakup, the projectile fragments more strongly than in the “white”stars. The data of Table
I points to the predominance of the channel with the 2+2+2+1 charge configuration (49
events) which has been studied in more detail. The obtained results show that the 14N
nucleus constitutes a very effective source for the production of 3α system.
In order to estimate the energy scale of production of 3α particle systems in the
14N→3α+X channel, we present the invariant excitation energy Q3α distribution with re-
spect to the 12C ground state:
M23α = −(
Q3α = M
3α −M(
12C) (4)
where M(12C) is the mass of the ground state corresponding to the charge and the weight
of the system being analyzed, M∗3α the invariant mass of the system of fragments. Statistics
was increased to 132 events 14N→3α+X including 50 “white ”stars by scanning over the
emulsion plates.
The main part of the events is concentrated in the Q3α area from 10 to 14 MeV, covering
the known 12C levels (Fig. 5). Softening of the conditions of the 3He + H selection, for which
the target fragment production is allowed, does not result in a shift of the 3α excitation peak.
This fact suggests the universality of the 3α state population mechanism.
To estimate the fraction of the events involving the production of an intermediate 8Be
nucleus in the reactions 14N→8Be+X→3α+X we present the invariant excitation energy
Q3α, MeV
0 10 20 30 40 50
FIG. 5: The invariant excitation energy Q3α distribution of three α particles with respect to the
12C ground state for the process 14N→3α+X. The following notation is used: 1) all the events of
the given dissociation, 2) “white”stars.
distribution for an α particle pair with respect to the 8Be ground state (Fig. 6). The first
distribution peak relates to the value to be expected for the decay products of an unstable
8Be nucleus in the ground state 0+. The distribution centre is seen to coincide well with the
decay energy of the 8Be ground state. The fraction of the α particles originating from the
8Be decay is 25-30%.
IV. FRAGMENTATION OF 7BE, AND 8B NUCLEI
The results of investigations dealing with the charge topology of the fragments produced
in peripheral dissociation of relativistic 8B, 7Be nuclei in emulsion are presented in Ref [2,
10, 11, 12].
Table II presents the numbers of the events detected in various channels of the 7Be
, MeVα2Q
0 2 4 6 8 10 12 14 16 18 20
, keVα2Q
0 50 100 150 200 250 300 350 400 450 500
FIG. 6: The invariant excitation energy Q2α distribution of α particle pairs for the process
14N→3α+X. In the inset: a fraction of the distribution at 0-500 keV.
fragmentation. Of them, the 3He+4He channel noticeably dominates, the channels 4He+d+p
and 6Li+p constitute 10% each. Two events involving no emission of neutrons in the three-
body channels 3He+t+p and 3He+d+d were registered. The reaction of charge-exchange
of 7Be nuclei to 7Li nuclei was not detected among the events not accompanied by other
secondary charged particles.The events involving no target fragments (nb=0) are separated
from the events involving one or a few fragments (nb >0).
For the first time, nuclear emulsions were exposed to a beam of relativistic 8B nuclei.
We have obtained data on the probabilities of the 8B fragmentation channels in peripheral
interactions. 55 events of the peripheral 8B dissociation which do not involve the produc-
tion of the target-nucleus fragments and mesons (“white” stars ) were selected. A leading
contribution of the 8B→7Be+p mode having the lowest energy threshold was revealed on
the basis of these events. Information about a relative probability of dissociation modes
with larger multiplicity have been obtained. Among the found events there are 320 stars in
which the total charge of the relativistic fragments in a 8◦ fragmentation cone ΣZfr satisfies
the condition ΣZfr >3. These stars were attributed to the number of peripheral dissocia-
tion events Npf . The Npf relativistic fragment distribution of over charges NZ is given in
TABLE II: 7Be fragmentation channel (number of events)
Channel 2He 2He He+2H He+2H 4H 4H Li+H Li+H Sum
nb =0 nb >0 nb =0 nb >0 nb =0 nb >0 nb =0 nb >0
3He+4He 30 11 41
3He+3He 11 7 18
4He+2p 13 9 22
4He+d+p 10 5 15
3He+2p 9 9 18
3He+d+p 8 10 18
3He+2d 1 1
3He+t+p 1 1
3p+d 2 2
2p+2d 1 1
6Li+p 9 3 12
Sum 41 18 42 33 2 1 9 3 149
Table III. There are given the data for 256 events containing the target-nucleus fragments
- Ntf , as well as for 64 events which contain no target-nucleus fragments (“white” stars )–
Npf . The role of the channels with multiple relativistic fragments NZ >2 is revealed to be
dominant for the N“white” stars. Of peripheral events, the “white” stars Nws (Table III)
are of very particular interest. They are not accompanied by the target-nucleus fragment
tracks and makes it possible to clarify the role of different cluster degrees of freedom at a
minimal excitation of the nuclear structure.
Table IV gives the relativistic fragment charge distribution in the “white” stars for 7Be
and 8B nuclei. The 8B events are presented without one single-charged relativistic fragment,
that is a supposed proton halo. The identical fraction of the two main 2He and He+2H
dissociation channels is observed for 7Be and 8B nuclei which points out that the 8Be core
excitation is independent of the presence of an additional loosely bound proton in the 8B
nucleus.
TABLE III: The charge topology distribution of the number of interactions of the peripheral Npf
type (Npf=Ntf+Nws), which were detected in an emulsion exposed to a second
8B nucleus beam.
Here Zfr is the total charge of relativistic fragments in a 8
◦ angular cone in an event, NZ the
number of fragments with charge Z in an event, Nws the number of “white”stars, Ntf the number
of events involving the target fragments, Nws the number of “white” stars.
Zfr N5 N4 N3 N2 N1 Ntf Nws
7 - - - 1 5 1 -
6 - - - 2 2 8 2
6 - - - 1 4 6 4
6 - - - - 6 1 -
5 - - - 1 3 61 14
5 - - - 2 1 44 12
5 - - 1 - 2 8 -
5 - - 1 1 - 1 -
5 - 1 - - 1 17 24
5 1 - - - - 17 1
5 - - - - 5 21 4
4 - - - - 4 5 1
4 - - - 2 - 24 1
4 - - - 1 2 42 -
TABLE IV: The charged dissociation mode distribution of the “white” stars produced by the 7Be
and 8B nuclei. To make the comparison more convenient, for the 8B nucleus one H nucleus is
eliminated from the charged mode and the channel fractions are indicated.
ΣZfr=4
7Be % 8B (+H) %
2He 41 43 12 40
He+2H 42 45 14 47
4H 2 2 4 13
V. CONCLUSIONS
The degree of the dissociation of the relativistic nuclei in peripheral interactions can reach
a total destruction into nucleons and singly and doubly charged fragments. The emulsion
technique allows one to observe these systems to the smallest details and gives the possibility
of studying them experimentally.
New experimental observations are reported from the emulsion exposures to 14N, 9Be,
8B, 7Be nuclei with energy above 1 A GeV. The main features of 9Be→2He relativistic frag-
mentation are presented. For the particular case of the relativistic 9Be nucleus dissociation
it is shown that precise angular measurements play a crucial role in the restoration of the
excitation spectrum of the alpha particle fragments. This nucleus is dissociated practically
totally through the 0+ and 2+ states of the 8Be nucleus. The data obtained from 9Be angular
measurements can be employed for the estimation of the role of 8Be in more complicated
Nα systems.
The results of the study of the dissociation of 14N nuclei of a primary momentum of
2.86 A GeV/c in their interactions with the emulsion nuclei are also presented. The present
investigation indicates the leading role of the 2+2+2+1 charge configuration channel. The
energy scale of the 3α system production has been estimated. According to the avail-
able statistics 80% of interactions are concentrated at 10-14 MeV. The fraction of the
14N→8Be+X→3α+X channel involving the production of an intermediate 8Be nucleus is
about 25%.
Advantages of emulsion technique are exploited most completely in the study of peripheral
fragmentation of light stable and neutron deficient nuclei. The results of investigations
dealing with the charge topology of the fragments produced in peripheral dissociation of
relativistic 7Be, 8B nuclei in emulsion are presented. Information on the relative probability
of dissociation modes with a larger multiplicity was obtained. The dissociation of a 7Be core
in 8B indicates an analogy with that of the free 7Be nucleus.
[1] N. P. Andreeva, et al., Phys. At. Nucl. 68, 455–465 (2005).
[2] Web site of the BECQUEREL Project: http://becquerel.jinr.ru (2006).
[3] M. I. Adamovich, et al., Phys. At. Nucl. 62, 1378–1387 (1999).
[4] M. I. Adamovich, et al., Phys. At. Nucl. 62, 514–517” (2004).
[5] D. A. Artemenkov, arXiv:nucl-ex/0605018 (2006).
[6] V. V. Belaga, et al., Phys. At. Nucl. 59, 869–877 (1996).
[7] F. A. Avetyan, et al., Phys. At. Nucl. 59, 110–116 (1996).
[8] V. V. Belaga, et al., Phys. At. Nucl. 58, 2014–2020 (1995).
[9] T. V. Shchedrina, et al., arXiv:nucl-ex/0605022 (2006).
[10] R. Stanoeva, et al., arXiv:nucl-ex/0605013 (2006).
[11] N. G.Peresadko, et al., arXiv:nucl-ex/0605014 (2006).
[12] N. P. Andreeva, et al., arXiv:nucl-ex/0604003 (2006).
http://becquerel.jinr.ru
Introduction
Fragmentation of 9Be nuclei
Fragmentation of 14N nuclei
Fragmentation of 7Be, and 8B nuclei
Conclusions
References
|
0704.0385 | Super-shell structures and pairing in ultracold trapped Fermi gases | Super-shell stru
tures and pairing in ultra
old trapped Fermi gases
Magnus Ögren
and Henning Heiselberg
Mathemati
al Physi
s, Lund Institute of Te
hnology, P.O. Box 118, SE-22100 Lund, Sweden
University of Southern Denmark, Campusvej 55, DK-5230 Odense M, Denmark
(Dated: April 3, 2007)
We
al
ulate level densities and pairing gaps for an ultra
old dilute gas of fermioni
atoms in
harmoni
traps under the in�uen
e of mean �eld and anharmoni
quarti
trap potentials. Super-
shell stru
tures, whi
h were found in Hartree-Fo
k
al
ulations, are
al
ulated analyti
ally within
periodi
orbit theory as well as from WKB
al
ulations. For attra
tive intera
tions, the underlying
level densities are
ru
ial for pairing and super-shell stru
tures in gaps are predi
ted.
PACS numbers: 03.75.Ss, 05.30.Fk
Ultra
old atomi
gases have re
ently been used to
re-
ate novel quantum many-body systems su
h as strongly
intera
ting high temperature super�uids of fermions,
Bose-Einstein
ondensates, Mott insulators in opti
al lat-
ti
es, et
. These lab phenomena have a strong over-
lap with
ondensed matter [1℄, nu
lear [2℄ and neutron
star physi
s [3℄. Finite fermion systems su
h as atoms in
traps, nu
lei, helium and metal
lusters, semi
ondu
tor
quantum dots, super
ondu
ting grains, et
., have addi-
tional interesting quantum stru
tures su
h as level spe
-
tra, densities and pairing. These will be observable as
temperatures are further lowered in atomi
trap experi-
ments. The high degree of
ontrol over physi
al param-
eters, in
luding intera
tion strength and density, makes
the atomi
traps marvelous model systems for general
quantum phenomena.
The purpose here is to
al
ulate the level spe
tra,
densities and pairing for zero-temperature Fermi gases
in harmoni
os
illator (HO) traps with anharmoni
and
mean �eld perturbations, and to show that novel super-
shell stru
tures appear in both level densities and pairing.
In
al
ulating level spe
tra by analyti
al periodi
orbit
theory and WKB as well as numeri
al Hartree-Fo
k, we
also relate these di�erent theoreti
al approa
hes to one
another.
We treat a gas ofN fermioni
atoms of massm in a HO
potential at zero temperature, intera
ting via a two-body
intera
tion with s-wave s
attering length a. We shall
mainly dis
uss a spheri
ally symmetri
trap and a dilute
gas (i.e. where the density ρ obeys the
ondition ρ|a|3 ≪
1) of parti
les with two spin states of equal population.
The Hamiltonian is then given by
mω2r2i + U(ri)
, (1)
We will
onsider both external anharmoni
potentials of
the form U = εr4 and parti
le intera
tions: U(ri) =
(2π~2a/m)
j 6=i δ
3(ri−rj). When intera
tions are weak,
the latter
an be approximated by the mean �eld poten-
U(r) =
2π~2a
ρ(r) . (2)
For a large number of parti
les and U = 0 the Fermi
energy is EF = ñF~ω where nF = ñF − 3/2 ≃ (3N)1/3
is the HO quantum number at the Fermi surfa
e. The
HO shells are highly degenerate with states having an-
gular momenta l = nF , nF −2, ...,mod(nF , 2), due to the
U(3) symmetry of the 3D spheri
ally symmetri
HO po-
tential. However, intera
tions split this degenera
y. In
the Thomas-Fermi (TF) approximation (see, e.g., [4℄) the
Fermi energy is
2k2F (r)
mω2r2 + U(r) . (3)
The density ρ(r) = k3F (r)/3π
ρ(r) = ρ0
1− r2/R2TF
, (4)
inside the
loud r ≤ RTF = aosc
2ñF , where ρ0 =
(2ñF )
3/2/3π2a3osc is the
entral density [5℄. For
onve-
nien
e we set the os
illator length aosc =
~/mω = 1 in
the following.
Taylor expanding the density and thereby also the
mean �eld of Eq. (2) around the
enter gives
ρ(r) ≃ ρ0
r2/R2TF +
r4/R4TF + ...
, (5)
the �rst term will simply in
orporate a
onstant shift
in energies whereas the term quadrati
in radius renor-
malizes the HO frequen
y as ωeff = ω
1− 6πaρ0/R2TF .
The third term is quarti
in radius and is therefore also
of the same form as the external potential
U(r) ≃ εr4 , (6)
with ε = (3π~2a/4m)ρ0/R
TF . Both the pure quarti
potential and the mean �eld potential of Eq. (2) are an-
harmoni
and
hange the level density by splitting the l
degenera
y of the HO shell nF at the Fermi surfa
e.
We will now
al
ulate analyti
ally the level spe
tra
from perturbative periodi
orbit theory for the quarti
potential and subsequently within semi
lassi
al WKB
wavefun
tions for both the quarti
and the mean �eld
potential of Eq. (2). We will start with repulsive inter-
a
tions where pairing is not present.
http://arxiv.org/abs/0704.0385v1
In periodi
orbit theory [6℄, the level density
an be
written (to leading order in ~
) in terms of a perturba-
tive HO tra
e formula [7, 8℄
g(E) =
1 +Re
(−1)k M ei2πkE/~ω
. (7)
For the unperturbed HO (U = 0) the modulation fa
tor is
M = 1. For a quarti
perturbed potential, as in Eq. (6),
the modulation fa
tor was
al
ulated in [8℄
e−i2kσ/~−iπ/2 + e−i3kσ/~+iπ/2
, (8)
with σ = επE2/~2ω3, being a small
lassi
al a
tion. The
two terms arise from the
hange in a
tions for the
ir
le
and diameter orbits respe
tively due to the quarti
po-
tential [8℄. The resulting level density
an be written in
the fa
torised form [9℄
g(E) =
(~ω)3
(−1)k
. (9)
Here, the �rst term is the average level density, the
osine fa
tor gives the rapid HO shell os
illations (mod-
i�ed by the perturbation) whi
h, however, are slowly
modulated by the sine fa
tor resulting in a beating pat-
tern. Moreover, the non-perturbed HO limit, equivalent
to M = 1 in Eq. (7), is re
overed in the limit of |ε| → 0,
where the U (3) symmetry is restored. The k = 1 term
in Eq. (9) gives the major os
illations in the level density
and is shown in Fig. 1 (a). The beating pattern or super-
shells is
learly observed. The shell os
illations vanish
when the argument of the sine in Eq. (9) is an integer
S = 1, 2, 3, ... times π, i.e. |ε|E2/2(~ω)3 = S. This gives
the supernode
ondition
nF = E/~ω =
2S~ω/|ε| . (10)
We now turn to an alternative
al
ulation of the level
density with WKB. The splitting of the HO shells degen-
erate levels l = nF , nF −2, ...,mod (nF , 2) in the shell nF
by the mean-�eld potential
an be
al
ulated perturba-
tively in the dilute limit. An ex
ellent approximation for
the radial HO wave fun
tion with angular momentum l
and (nF − l)/2 radial nodes in the HO shell when nF ≫ 1
is the WKB one [10, 11℄:
RnF l(r) ≃
sin(kl(r)r + θ)
l (r)r
, (11)
between turning points r2± = ñF ±
ñ2F − l(l + 1). Here,
ñF = nF + 3/2 and the WKB wave number kl(r) is
k2l (r) = 2ñF − r2 − l(l+ 1)/r2 . (12)
When nF ≫ 1 the wave fun
tion has many nodes 1 ≪
l ≪ nF and the os
illations in R2nl(r)
an be averaged
〈sin2(kl(r)r)〉 = 1/2 [10℄. The phase θ is then unim-
portant. The single-parti
le energies for the anharmoni
potential of Eq. (6) are simply
EnF ,l − ñF~ω =
U(r)|RnF l(r)|2r2dr (13)
πkl(r)
3ñ2F − l(l + 1)
.(14)
It is spe
ial for the quarti
perturbation that the level
energies are linear in l(l+1). The resulting level spa
ing
in
reases as (2l+1) just as the level degenera
y for SO(3)
symmetry. Therefore the level density is
onstant within
the bandwidth
D ≡ |EnF ,l≡0 − EnF ,l=nF | = εn2F /2 (15)
on energy s
ales larger than 2D/nF but smaller than D.
The level density vanishes between the bandwidths of
two neighbouring n shells and therefore it generally has
a strong os
illatory behavior as shown in Fig. 1 (a). Its
amplitude is largest when D ∼ ~ω/2. However, when
D ≃ ~ω the level density is
onstant and the os
illa-
tory behavior vanishes. This phenomenon repeats when
D = S~ω sin
e the level spe
tra then overlap S times.
With the bandwidth of Eq. (15) under this
ondition, we
obtain exa
tly the same supernode
ondition as for pe-
riodi
orbit theory, Eq. (10). We
on
lude that Craig's
perturbative periodi
orbit theory [7℄ is in exa
t agree-
ment with perturbative WKB for a quarti
ally perturbed
spheri
al symmetri
HO in three dimensions.
We now turn to the slightly more
ompli
ated mean
�eld potential of Eq. (2). Its level spe
trum
an also be
al
ulated from the WKB wave fun
tions of Eq. (11).
Inserting them in Eq. (13), we obtain
EnF ,l − ñF~ω = 2/
F ~ω I . (16)
Here, the integral I is
1− l(l+ 1)/ñ2F
1− x2
where x = (r2−ñF )/
ñ2F − l(l + 1). This integral is I =
π for l ≃ nF and I = 8
2/3 for l = 0. The bandwidth is
therefore
D = 2/
2/3− π
. (18)
Inserting this bandwidth in the supernode
ondition D =
S~ω gives
2/3− π
F = S . (19)
For example in the
ase 2πa = 1 the supernodes S =
1, 2, 3, .. should o
ur when nF ∼ 28, 44, 58, et
. The
Hartree-Fo
k (HF)
al
ulations of the os
illating part of
the total energy, whi
h is proportional to the level den-
sity at the Fermi level [6℄, result in slightly higher su-
pernodes, as in Fig 1 (b). The di�eren
es arise be
ause
the WKB
al
ulations are perturbative in the intera
tion
strength, whereas in the HF
al
ulation the MF poten-
tial U in
ludes a large s
attering length whi
h, e.g., leads
to
orre
tions for the e�e
tive os
illator frequen
y. Also
for the purely quarti
term the perturbative approa
h
underestimates the exa
t supernodes (see Fig. 3 of [8℄).
For weaker intera
tions 2πa = 0.1, the �rst supernode
S = 1 should o
ur at nF = 130 a
ording to the
ondi-
tion of Eq. (19), in
loser agreement with the HF result
of Fig. 1 (
).
For
omparison, the Taylor expansion of the mean �eld
potential leads to the supernode
ondition of Eq. (10)
with ε = (3π~2a/4m)ρ0/R
TF . It di�ers from Eq. (19)
by the prefa
tor, whi
h is ∼ 34% smaller. It is a better
approximation to expand e.g. around r = RTF /2
nF /2, where the
orresponding prefa
tor is only ∼ 8%
smaller, su
h that the supernode in Fig 1 (
) is predi
ted
to nF = 137. Now expanding I of Eq. (17) for small
l ≪ nF , one �nds
I = (8/3)
2− l2/
2n2F , (20)
resulting in the level spe
trum [10℄
EnF ,l − ñ~ω =
− l(l+ 1)
. (21)
This level density is
onstant at low l as for the potential
in Eq. (14). However, near l ∼ nF the density of lev-
els is slightly smaller as
an be seen from the bandwidth
orresponding to Eq. (21), whi
h is ∼ 12% larger, for a
given nF , than the bandwidth of Eq. (18). That the level
density is not
ompletely
onstant within the bandwidth
has the e�e
t that a small periodi
ity remains even at
the super-shell
ondition D = S~ω. Therefore the shell
os
illations do not disappear
ompletely at the supern-
odes, as
an be seen in Fig. 1 (b,
), whereas for the purely
quarti
ase (a) the os
illations disappear
ompletely at
the supernodes.
Most atomi
traps are not spheri
al but
igar shaped
(prolate) with ωz
<∼ω⊥. The unperturbed HO energies
E = nz~ωz + n⊥~ω⊥ will generally lead to a
onstant
level density for energy s
ales larger than ~ωz but smaller
than ~ω⊥. When the os
illator frequen
y ratio ω⊥/ωz is a
rational number, level degenera
ies and larger os
illations
will o
ur on the s
ale ~ωz. Intera
tions will, however,
smear this level density. In any
ase, super-shell stru
ture
is not expe
ted as in the spheri
al symmetri
ase. In
very oblate traps ωz ≫ ω⊥ the mean �eld potential is
e�e
tively two-dimensional and quadrati
, i.e. it does
not split the HO shells [10, 13℄. Thus we may expe
t
strong os
illations in the level density on the s
ale ~ω⊥,
but again no super-shell stru
ture.
0 20 40 60 80 100
−1000
1000 (a)
0 20 40 60 80 100 120 140
200 (c)
0 20 40 60 80 100
40 (b)
Figure 1: (
olor online) The upper �gure (a) shows the leading
term (k = 1) of the os
illating part of the perturbative level
density of Eq. (9) as a fun
tion of nF = E/~ω, for the
ase of
an external potential V = VHO+εr
with ε = 2/402 ≈ 0.0013.
The middle and lower �gures (b,
) show the os
illating part
of the total energy a
ording to a numeri
al HF
al
ulation
[12℄, with intera
tion strength 2πa = 1 and 2πa = 0.1, as a
fun
tion of the HO shell number (~ = ω = 1). This illustrates
qualitatively that a supernode, e.g. at nF = 40,
an be due
to intera
tion (b) and/or an additional quarti
term to the
HO potential (a).
Attra
tive intera
tions lead to pairing by an amount
that is exponentially sensitive to the underlying level den-
sity near the Fermi surfa
e [2, 10, 11, 14℄. The level den-
sity is the same for repulsive and attra
tive intera
tions
ex
ept that the levels are reversed when the sign of ε
(Eqs. (9)) and (13)) or a is
hanged (Eq. (16)). Therefore
we
an use the level densities and bandwidths
al
ulated
above for pairing
al
ulations. Pairing in �nite systems
is des
ribed by the Bogoliubov-de Gennes (BdG) equa-
tions [15℄ and take pla
e between time-reversed states.
As shown in [14℄ these states
an be approximated by HO
wave fun
tions in dilute HO traps as long as the gap does
not ex
eed the os
illator energy, ∆<∼~ω. Solving BdG for
su
h �nite systems is numeri
ally
ompli
ated and we
shall therefore apply further simplifying approximations,
namely that the pairing gap ∆nl and the wavefun
tion
overlap matrix elements vary slowly with level l in a shell
n. Both approximations are fair for the trapped atoms
as argued in [11℄ and deviations
an be understood. As
result we arrive at a mu
h simpli�ed gap equation
∫ ∼2nF
g(E) dE
(E − µ)2 +∆(µ)2
. (22)
Here, the supergap G = 32
2nF |a|~ω/15π2 was
al-
ulated in [10℄ as the pairing gap when all states in
a shell
an pair; this is the
ase for a region of in-
tera
tion strengths and parti
le number where the gap
is large as
ompared to the level splitting, yet small
0 10 20 30 40 50 60
=(3N)1/3
38 40 42 44
Figure 2: (
olor online) Multi-shell pairing gaps for a HO trap
with an additional quarti
term in the potential with ε =
2/402, i.e. for the level density of Fig. 1 (a) with supernodes
at nF ≃ 40, 40
2 ≈ 57, etc. The intera
tion strength is a =
−0.05 (top red
urve), a = −0.03 (middle blue
urve, with the
inset �gure around the �rst supernode) and a = −0.01 (lower
green
urve). In the inset plot it is
learly seen that the lo
al
minima for l ∼ nF and l ∼ 0 before the supernode turns
into lo
al maxima after the supernode, as a
onsequen
e of
overlapping shells. The dashed (red) line is the multi-shell gap
∆ = G/(1−2G ln(nF )/~ω) for a = −0.05 and the upper/lower
thin solid line (bla
k) are the single mid-/end-shell pairing for
a = −0.01 (see text).
ompared to the shell splitting ~ω. ∆(µ) = ∆nl is
the gap at the Fermi surfa
e. g(E) = n2F /D is
the level density within ea
h bandgap D around ev-
ery shell n = 0, 1, ...,∼2nF but vanishes between
the bandgaps. The gap equation thus redu
es to
1 = (G/D)
∑∼2nF
(E + n~ω − µ)2 +∆2. The
hemi
al potential µ
an be determined from the level
spe
trum; as we gradually �ll parti
les into the shell nF
at the Fermi surfa
e, µ in
reases from nF~ω to nF~ω+D.
The
ut-o� n<∼2nF in the sum of the gap equation models
as a �rst approximation the more rigorous regularization
pro
edure des
ribed in Ref. [16℄ that is required for a
delta-fun
tion pseudo-potential.
By solving this simpli�ed gap equation of Eq. (22),
we �nd that it still
ontains and displays the essential
interplay between the variation in level density and pair-
ing. To illustrate the super-shell stru
ture in pairing,
we take the strongly anharmoni
trap potential used for
the level spe
tra in Fig. 1 (a), and
al
ulate the pairing
arising from a weak attra
tive s
attering length a < 0.
For su�
iently weak intera
tions su
h that pairing only
takes pla
e in the shell at the Fermi surfa
e, we obtain
the expe
ted result from the gap equation: ∆ = G when
D ≪ ∆, whereas for D ≫ ∆ we get ∆ = D exp(−D/2G)
midshell (µ = nF~ω + D/2) and ∆ = 2D exp(−D/G)
endshell (µ = nF~ω or µ = nF~ω +D). Pairing is thus
stronger at midshell than at endshell, where there are
fewer states to pair [11℄, and strong shell os
illations fol-
low as shown in Fig. 2. For stronger intera
tions, pairing
also takes pla
e between states in shells around the Fermi
shell and Eq. (22) gives: ∆ = G/ (1− 2G ln (nF ) /~ω) for
small bandwidth [14℄. In Fig. 2 this
urve is
ompared
with the �nite bandwidth result, whi
h has strong os
il-
lations ex
ept at the supernodes where the level density
is
ontinuous. At a supernode D = ~ω and the gap equa-
tion (22) leads to a gap ∆ = 2nF~ω exp(−~ω/2G) [11℄.
In summary, level densities, shell-os
illations and
super-shell stru
tures in anharmoni
traps
al
ulated
from numeri
al Hartree-Fo
k and analyti
al periodi
or-
bit theory as well as WKB were found to mat
h to lead-
ing order. Analogous super-shell stru
tures were found
in pairing from an approximated BdG
al
ulation. The
mean �eld in atomi
nu
lei also have a large anharmoni
potential and the HO shells start to overlap (the �rst
supernode) already for heavy nu
lei with nF ∼ 5 − 6.
The interplay of level spe
tra and multishell pairing is,
however, di�
ult to disentangle in nu
lear pairing due to
strong spin-orbit e�e
t and small parti
le number. Ul-
tra
old atomi
traps, however, provide ideal systems for
observing the ri
h quantum stru
tures su
h as level den-
sities and pairing.
Dis
ussions with Matthias Bra
k on periodi
orbit the-
ory, Ben Mottelson on (nu
lear) shell theory and pairing,
and proof reading by Joel Corney, are gratefully a
knowl-
edged.
[1℄ J. Bardeen, L. N. Cooper, J. R. S
hrie�er, Phys. Rev.
108, 1175 (1957).
[2℄ A. Bohr and B. R. Mottelson, Nu
lear Stru
ture Vols.
I+II, Benjamin, New York 1969.
[3℄ A. Bohr, B. R. Mottelson, D. Pines, Phys. Rev. 110, 936
(1958).
[4℄ C. J. Pethi
k and H. Smith, Bose-Einstein Condensation
in Dilute Gases, Cambridge Univ. Press 2002.
[5℄ For a �nite number of parti
les the fa
tor ñ = nF + 3/2
in
ludes a
orre
tion to nF , whi
h has been
he
ked nu-
meri
ally to improve the TF approximation and slightly
hange the predi
tion of supernodes.
[6℄ M. Bra
k and R. K. Bhaduri, Semi
lassi
al Physi
s, re-
vised edn (Boulder, CO: Westview) (2003).
[7℄ S. C. Creagh, Ann. Phys., NY 248 60 (1996).
[8℄ M. Bra
k et al., J. Phys. A 38, 9941 (2005).
[9℄ M. Ögren, unpublished (2006):
www.magnus.ogren.se/notes/pot/derivationofgpert.pdf
[10℄ H. Heiselberg and B. R. Mottelson, Phys. Rev. Lett. 88,
190401 (2002).
[11℄ H. Heiselberg, Phys. Rev. A 68, 053616 (2003). Note that
the square root of kl was missing in Eq. (6) of this Ref.
as
ompared to Eq. (11).
[12℄ Y. Yu et al., Phys. Rev A 72, 051602(R) (2005).
[13℄ B. P. van Zyl et al., Phys. Rev. A 67, 023609 (2003).
[14℄ G. M. Bruun and H. Heiselberg, Phys. Rev. A 65, 053407
(2002).
[15℄ P. G. de Gennes, Super
ondu
tivity of Metals and Alloys
(Addison-Wesley, New York, 1989).
[16℄ G. M. Bruun et al., Eur. Phys. J. D9, 433 (1999).
|
0704.0386 | Quantum non-local effects with Bose-Einstein condensates | Quantum non-local effects with Bose-Einstein condensates
F. Laloë a and W. J. Mullin b
Laboratoire Kastler Brossel, ENS, UPMC, CNRS; 24 rue Lhomond, 75005 Paris, France
Department of Physics, University of Massachusetts, Amherst, Massachusetts 01003 USA
We study theoretically the properties of two Bose-Einstein condensates in different spin states,
represented by a double Fock state. Individual measurements of the spins of the particles are per-
formed in transverse directions, giving access to the relative phase of the condensates. Initially, this
phase is completely undefined, and the first measurements provide random results. But a fixed value
of this phase rapidly emerges under the effect of the successive quantum measurements, giving rise
to a quasi-classical situation where all spins have parallel transverse orientations. If the number of
measurements reaches its maximum (the number of particles), quantum effects show up again, giving
rise to violations of Bell type inequalities. The violation of BCHSH inequalities with an arbitrarily
large number of spins may be comparable (or even equal) to that obtained with two spins.
PACS numbers: 03.65.Ta,03.65.Ud,03.75.Gg,03.75.Mn
The notion of non-locality in quantum mechanics
(QM) takes its roots in a chain of two theorems, the
EPR (Einstein Podolsky Rosen) theorem [1] and its log-
ical continuation, the Bell theorem. The EPR theorem
starts from three assumptions (Einstein realism, locality,
the predictions of quantum mechanics concerning some
perfect correlations are correct) and proves that QM is
incomplete: additional quantities, traditionally named λ,
are necessary to complete the description of physical re-
ality. The Bell theorem [2, 3] then proves that, if λ exists,
the predictions of QM concerning other imperfect corre-
lations cannot always be correct. The ensemble of the
three assumptions: Einstein realism, locality, all predic-
tions of QM are correct, is therefore self-contradictory;
if Einstein realism is valid, QM is non-local. Bohr [4]
rejected Einstein realism because, in his view, the no-
tion of physical reality could not correctly be applied to
microscopic quantum systems, defined independently of
the measurement apparatuses. Indeed, since EPR con-
sider a system of two microscopic particles, which can be
“seen” only with the help of measurement apparatuses,
the notion of their independent physical reality is open
to discussion.
Nevertheless, it has been pointed out recently [5, 6]
that the EPR theorem also applies to macroscopic sys-
tems, namely Bose-Einstein (BE) condensates in two dif-
ferent internal states. The λ introduced by EPR then cor-
responds to the relative phase of the condensates, i.e. to
macroscopic transverse spin orientations, physical quan-
tities at a human scale; it then seems more difficult to
deny the existence of their reality, even in the absence
of measurement devices. This gives even more strength
to the EPR argument and weakens Bohr’s refutation. It
is then natural to ask whether the Bell theorem can be
transposed to this stronger case.
The purpose of this article is to show that it can. We
consider an ensemble of N+ particles in a state defined by
an orbital state u and a spin state +, and N− particles
in the same state with spin orientation −. The whole
system is described quantum mechanically by a double
Fock state, that is, a “double BE condensate”:
| Φ > =
(au,+)
]N+ [
(au,−)
| vac. > (1)
where au,+ and au,− are the destruction operators asso-
ciated with the two populated single-particle states and
|vac. > is the vacuum state. We introduce a sequence
of transverse spin measurements that leads to quantum
predictions violating the so called BCHSH [7, 8] Bell in-
equality. This is reminiscent of the work of Mermin [9],
who finds exponential violations of local realist inequal-
ities with N -particle spin states that are maximally en-
tangled. By contrast, here we consider the simplest way
in which many bosons can be put in two different in-
ternal levels, with a N -particle state containing only the
minimal possible correlations, those due to statistics. We
find violations of inequalities that are the same order of
magnitude as with the usual singlet spin state and may
actually saturate the Cirel’son bound [10].
Double Fock states are experimentally more accessi-
ble and much less sensitive to dissipation and decoher-
ence than maximally entangled states [11]. Considering
a system in a double Fock state, we assume that a se-
ries of rapid spin measurements can be performed and
described by the usual QM postulate of measurement,
without worrying about decoherence between the mea-
surements, thermal effects, etc.
The operators associated with the local density of par-
ticles and spins can be expressed as functions of the
two fields operators Ψ±(r) associated with the two in-
ternal states ± as: n(r) = Ψ†+(r)Ψ+(r) + Ψ
−(r)Ψ−(r),
σz(r) = Ψ
+(r)Ψ+(r)−Ψ
−(r)Ψ−(r), while the spin com-
ponent in the direction of plane xOy making an angle ϕ
with Ox is: σϕ(r) = e
+(r)Ψ−(r)+ e
−(r)Ψ+(r).
Consider now a measurement of this component per-
formed at point r and providing result η = ±1. The
http://arxiv.org/abs/0704.0386v4
corresponding projector is:
Pη=±1(r, ϕ) =
[n(r) + η σϕ(r)] (2)
and, because the measurements are supposed to be per-
formed at different points (ensuring that these projectors
all commute) the probability P(η1, η2, ...ηN ) for a series
of results ηi± 1 for spin measurements at points ri along
directions ϕi can be written as:
< Φ | Pη1(r1, ϕ1)× Pη2(r2, ϕ2)× ....PηN (rN , ϕN ) | Φ >
We now substitute the expression of σϕ(r) into (2) and
(3), exactly as in the calculation of ref. [5], but with one
difference: here we do not assume that the number of
measurements is much smaller than N±, but equal to
its maximum value N = N+ + N−. In the product of
projectors appearing in (3), because all r’s are different,
commutation allows us to push all the field operators to
the right, all their conjugates to the left; one can then
easily see that each Ψ±(r) acting on | Φ > can be re-
placed by u(r) × au,± , and similarly for the Hermitian
conjugates. With our initial state, a non-zero result can
be obtained only if exactly N+ operators au,+ appear in
the term considered, and N− operators au,−; a similar
condition exists for the Hermitian conjugate operators.
To express these conditions, we introduce two additional
variables. As in [5], the first variable λ ensures an equal
number of creation and destruction operators in the in-
ternal states ± through the mathematical identity:
einλ = δn,0 (4)
The second variable Λ expresses in a similar way that the
difference between the number of destruction operators in
states + and − is exactly N+−N−, through the integral:
e−inΛ ei(N+−N−)Λ = δn,N+−N− (5)
The introduction of the corresponding exponentials into
the product of projectors (2) in (3) provides the expres-
sion (c.c. means complex conjugate):
|u(rj)|2
eiΛ + e−iΛ + ηj
ei(λ−ϕj+Λ) + c.c.
where, after integration over λ and Λ, the only surviving
terms are all associated with the same matrix element in
state | Φ > (that of the product of N+ operators a†u,+
and N− operators a
u,− followed by the same sequence
of destruction operators, providing the constant result
N+!N−!). We can thus write the probability as:
P(η1, η2, ...ηN ) ∼
ei(N+−N−)Λ
|u(rj)|2
eiΛ + e−iΛ + ηj
ei(λ−ϕj+Λ) + c.c.
or, by using Λ parity and changing one integration variable (λ′ = λ+ Λ), as:
P(η1, η2, ...ηN ) =
cos [(N+ −N−)Λ]
{cos (Λ) + ηj cos (λ′ − ϕj)} (8)
The normalization coefficient CN is readily obtained by writing that the sum of probabilities of all possible sequences
of η’s is 1 (this step requires discussion; we come back to this point at the end of this article):
cos [(N+ −N−)Λ] [cos (Λ)]N (9)
Finally, we generalize (8) to any number of measurements M < N . A sequence of M measurements can always be
completed by additional N −M measurements, leading to probability (8). We can therefore take the sum of (8) over
all possible results of the additional N −M measurements to obtain the probability for any M as:
P(η1, η2, ...ηM ) =
cos [(N+ −N−)Λ] [cosΛ]N−M
{cos (Λ) + ηj cos (λ′ − ϕj)} (10)
The Λ integral can be replaced by twice the integral between ±π/2 (a change of Λ into π −Λ multiplies the function
by (−1)N+−N−+N−M+M = 1). If M ≪ N , the large power of cosΛ in the first integral concentrates its contribution
around Λ ≃ 0, so that a good approximation is Λ = 0. We then recover the results of refs [5, 6], with a single integral
over λ defining the relative phase of the condensates (Anderson phase), initially completely undetermined, so that
the first spin measurement provides a completely random result. But the phase rapidly emerges under the effect of a
few measurements, and remains constant [12, 13, 14]; it takes a different value for each realization of the experiment,
as if it was revealing the pre-existing value of a classical quantity. Moreover, when cosΛ is replaced by 1, each factor
of the product over j remains positive (or zero), leading to a result similar to that of stochastic local realist theories;
the Bell inequalities can then be obtained. However, when N − M is small or even vanishes, cosΛ can take values
that are smaller than 1 and the factors may become negative, opening the possibility of violations. In a sense, the
additional variable Λ controls the amount of quantum effects in the series of measurements.
We now discuss when these standard QM predictions violate Bell inequalities. We need the value of the quantum
average of the product of results, that is the sum of η1, η2, ...ηM × P(η1, η2, ...ηM ) over all possible values of the η’s,
which according to (10) is given by:
E(ϕ1, ϕ2, ..ϕM ) =
cos [(N+ −N−)Λ] [cosΛ]N−M
Consider a thought experiment where two condensates
in different spin states (two eigenstates of the Oz spin
component) overlap in two remote regions of space A
and B , with two experimentalists Alice and Bob; they
measure the spins of the particles in arbitrary transverse
directions (perpendicular to Oz) at points of space where
the orbital wave functions of the two condensates are
equal. All measurements performed by Alice are made
along a single direction ϕa, which plays here the usual
role of the “setting” a, while all those performed by Bob
are made along angle ϕb. We assume that Alice retains
just the product A of all her measurements, while Bob
retains only the product B of his; A and B are both ±1.
We now assume two possible orientations ϕa and ϕ
for Alice, two possible orientations ϕb and ϕ
b for Bob.
Within deterministic local realism, for each realization of
the experiment, it is possible to define two numbers A,
A′, both equal to ±1, associated with the two possible
products of results η that Alice will observe, depending
of her choice of orientation; the same is obviously true
for Bob, introducing B and B′. Within stochastic local
realism [8, 15], A and A′ are the difference of probabilities
associated with Alice observing +1 or −1, i.e. numbers
that have values between +1 and −1. In both cases, the
following inequalities (BCHSH) are obeyed:
− 2 ≤ AB +AB′ ± (A′B − A′B′) ≤ 2 (12)
In standard quantum mechanics, of course, “unper-
formed experiments have no results” [16], and several of
the numbers appearing in (12) are undefined; only two
of them can be defined after the experiment has been
performed with a given choice of the orientations. Thus,
while one can calculate from (11) the quantum average
value < Q > of the sum of products of results appearing
in (12), there is no special reason why < Q > should be
limited between +2 and −2. Situations where the limit
is exceeded are called “quantum non-local”.
We have seen that the most interesting situations oc-
cur when the cosines do not introduce their peaking effect
around Λ = 0, i.e. when N+ = N− and M has its maxi-
mum value N . Then, for a given N , the only remaining
choice is how the number of measurements is shared be-
tween Na measurements for Alice and Nb for Bob.
Assume first that Na = 1 (Alice makes one measure-
ment) and therefore Nb = N − 1 (Bob makes all the oth-
ers). Since we assume that N+ = N− and M = N , the
Λ integral in (11) disappears, and the λ integral contains
only the product of cos (λ′ − ϕa) by the (N − 1)th power
of cos (λ′ − ϕb), which is straightforward and provides
cos (ϕa − ϕb) times the normalization integral CN . The
quantum average associated with the product AB is thus
merely equal to cos (ϕa − ϕb), exactly as the usual case
of two spins in a singlet state. Then it is well-known that,
when the angles form a “fan” [17] spaced by χ = π/4,
a strong violation of (12) occurs, by a factor
2, sat-
urating the Cirel’son bound [10]. A similar calculation
can be performed when Alice makes 2 measurements and
Bob N − 2, and shows that the quantum average is
now equal to 1
1 + 1
+ (1 − 1
) cos 2 (ϕa − ϕb)
no longer independent of N. If N = 4, the maximum
of < Q > is 2.28 < 2
2, and rises to 2.41 as N → ∞.
An expression for the generalization of the quantum av-
erage to any number P and N − P of measurements by
Alice and Bob, respectively, is (with χ = ϕa − ϕb):
E(χ) =
{P/2}
P !(N − 2k)!
k!(P − 2k)!(N
− k)!
sin2k χ cosP−2k χ
where {P/2} is the integer part of P/2. The maximum of
< Q > can then be found using a numerical Mathematica
routine. Results are shown for several values of P in Fig.
1. The angles maximizing the quantum Bell quantity
always occur in the fan shape, although the basic angle
χ changes with P and N. All of the curves where P is
held fixed have a finite < Q > limit with increasing N ,
and the optimum values of the angles approach constants.
For the curve P = N/2, the limit is 2.32 when N → ∞,
and the fan opening decreases as 1/
10080604020
P
FIG. 1: The maximum of the quantum average < Q > for
Alice doing P experiments and Bob N − P , as a function of
the total number of particles N . The usual Bell situation is
obtained for N = 2, P = 1. Local realist theories predict an
upper limit of 2; large violations of this limit are obtained,
even with macroscopic systems (N → ∞). If P = 1, the
violation saturates the Cirel’son limit for any N .
We can also study cases where the number of measure-
ments is M < N : if Bob makes all his measurements,
but ignores one or two of them (independently of the or-
der of the measurements), when he correlates his results
with Alice, the BCHSH inequality is never violated. All
measurements have to be taken into account to obtain
violations. Furthermore, if the number of particles in the
two condensates are not equal, no violation occurs either.
Finally, it is possible to consider cases where we gener-
alize the angles considered: experimenter Carole makes
measurements at ϕc and ϕ
c, and David at ϕd and ϕ
We then find that a maximization of < Q > reduces to
the cases already studied, where the new angles collapse
to the previous angles ϕa, · · · , ϕ′b.
For the sake of simplicity, we have not yet discussed
some important issues that underlie our calculations. One
is related to the so called “sample bias loophole” (or
“detection/efficiency loophole”) and to the normalization
condition (9), which assumes that one spin is detected at
each point of measurement. A more detailed study (see
second ref. [5]) should include the integration of each
r in a small detection volume and the possibility that
no particle is detected in it. This is a well-known dif-
ficulty, which already appears in the usual two-photon
experiments [8], where most photons are missed by the
detectors. If this loophole still raises a real experimen-
tal challenge, the difficulty can be resolved in theory by
assuming the presence of additional spin-independent de-
tectors [2, 8], which ensure the detection of one particle
in each detector and create appropriate initial conditions
(see for instance [18] for a description of an experiment
with veto detectors). We postpone this discussion to an-
other article [19]. A second issue deals with the definition
of the local realist quantities A, B, etc. For two conden-
sates, we have a slightly different situation than in the
usual EPR situation: the local realist reasoning leads to
the existence of a well-defined phase λ between the con-
densates [5], not necessarily to deterministic properties
of the individual particles. Fortunately, Bell inequalities
can also be derived within stochastic local realist theories
[3, 8] (see also for instance [9] or appendix I of [15]), and
this difference is not a problem [19].
In conclusion, strong violations of local realism may
occur for large quantum systems, even if the state is a
simple double Fock state with equal populations; within
present experimental techniques, this seems reachable
with N ∼ 10 or 20. We have assumed that the mea-
sured quantity is the product of many microscopic mea-
surements, not their sum, which would be macroscopic; a
product of results remains sensitive to the last measure-
ment, even after a long sequence of others. Curiously, for
very few measurements only the results are quantum, for
many measurements they can be interpreted in terms of
a classical phase, but become again strongly quantum
when the maximum number of measurements is reached,
a sort of revival of quantum-ness of the system.
Laboratoire Kastler Brossel is “UMR 8552 du CNRS,
de l’ENS, et de l’Université Pierre et Marie Curie”.
[1] A. Einstein, B. Podolsky and N. Rosen, Phys. Rev. 47,
777 (1935).
[2] J.S. Bell, Physics 1, 195 (1964), reprinted in [3].
[3] J.S. Bell, “Speakable and unspeakable in quantum me-
chanics”, Cambridge University Press (1987).
[4] N. Bohr, Phys. Rev. 48, 696 (1935).
[5] F. Laloë, Europ. Phys. J. D, 33, 87 (2005); see also
cond-mat/0611043.
[6] W.J. Mullin, R. Krotkov and F. Laloë, Phys. Rev. A74,
023610 (2006).
[7] J.F. Clauser, M.A. Horne, A. Shimony and R.A. Holt,
Phys. Rev. Lett. 23, 880 (1969).
[8] J.F. Clauser and A. Shimony, Rep. on Progress in Phys.
41, 1883 (1978).
[9] N.D. Mermin, Phys. Rev. Lett. 65, 1838 (1990).
[10] B.S. Cirel’son, Letters in math. phys. 4, 93 (1980).
[11] J.A. Dunningham, K. Burnett and S.M. Barnett, Phys.
Rev. Lett. 89, 150401 (2002).
[12] J. Javanainen and Sung Mi Yoo, Phys. Rev. Lett. 76,
161 (1996).
[13] Y. Castin and J. Dalibard, Phys. Rev. A55, 4330 (1997).
[14] I. Cirac, C. Gardiner, M. Naraschewski and P. Zoller,
Phys. Rev. A54, R3714 (1996) and references in [6]
[15] F. Laloë, Am. J. Phys. 69, 655 (2001).
[16] A. Peres, Am. J. Phys. 46, 745 (1978).
[17] The term “fan” refers to the angles arranged as ϕba =
ϕa′b = ϕab′ and ϕa′b′ = 3χ where ϕab ≡ ϕa − ϕb.
[18] J.S. Bell, Comments on at. and mol. phys. 9, 121 (1979);
reprinted in [3].
[19] W.J. Mullin and F. Laloë, to be published
http://arxiv.org/abs/cond-mat/0611043
|
0704.0387 | Low mass visual binaries in the solar neighbourhood: The case of
HD141272 | Astron. Nachr. / AN Volume, No. Issue, 0 – 5 (Year of publication) / DOI DOI
Low mass visual binaries in the solar neighbourhood:
The case of HD141272⋆
T. Eisenbeiss1 ⋆⋆, A. Seifahrt1,2, M. Mugrauer1, T. O. B. Schmidt1, R. Neuhäuser1, and
T. Roell1
1 Astrophysikalisches Institut und Universitäts-Sternwarte Jena, Schillergässchen 2-3, 07745 Jena, Germany
2 European Southern Observatory, Karl-Schwarzschild-Str. 2, 85748, Garching, Germany
Received 15.08.06, accepted 29.03.07
Published online . . .
Key words binaries: visual – stars: late-type, low mass – astrometry
We search for stellar and substellar companions of young nearby stars to investigate stellar multiplicity and
formation of stellar and substellar companions. We detect common proper-motion companions of stars via
multi-epoch imaging. Their companionship is finally confirmed with photometry and spectroscopy. Here we
report the discovery of a new co-moving (13σ) stellar companion ∼ 17.8 arcsec (350AU in projected sepa-
ration) north of the nearby star HD141272 (21 pc). With EMMI/NTT optical spectroscopy we determined
the spectral type of the companion to be M3±0.5V. The derived spectral type as well as the near infrared
photometry of the companion are both fully consistent with a 0.26+0.07−0.06M⊙ dwarf located at the distance
of HD141272 (21 pc). Furthermore the photometry data rules out the pre-main sequence status, since the
system is consistent with the ZAMS of the Pleiades.
c© Year of publication WILEY-VCH Verlag GmbH&Co. KGaA, Weinheim
1 Introduction
HD141272 is a nearby G8 dwarf with a mass of
0.83+0.07
−0.03M⊙ (Nordström et al. 2004) located in the
constellation Serpens Caput (αJ2000.0 = 15
h 48m 09.4s,
δJ2000.0 = +01
◦ 34′ 18′′). Its proper motion (µα cos δ =
−176.19 ± 1.08mas/yr, µδ = −166.72 ± 1.13mas/yr)
and parallax (π = 46.84± 1.05mas, i.e. 21 pc) are both
well determined by the European astrometry satellite
Hipparcos (Perryman et al. 1997). While Montes et al.
(2001) list HD 141272 as a member of the Local asso-
ciation with an age of ∼ 120Myr (Mart́ın et al. 2001),
Fuhrmann (2004) suggested that this star belongs to
the young Her-Lyr moving group, according to its UV-
velocities. The age of some Her-Lyr members is esti-
mated by Fuhrmann (2004) to approximately 100Myr
(e.g. HR857, HD 82443, HD113449 and HR5829) which
recently reached their main sequence position, while
others seemed to be older than ∼ 200Myr
(Fuhrmann 2004). Also Fuhrmann (2004) argued that
HD141272, with an effective temperature of Teff =
(5270±80)K, an absolute bolometric magnitudeMbol =
(5.54±0.07)mag and metallicity of [Fe/H ] = (−0.08±
0.07) dex appears slightly too bright for its main se-
quence position, indicating that it might be non single
or young.
⋆ Based on observations obtained on La Silla in ESO programs
77.C-0572(A) and Calar Alto project number F06-3.5-016.
⋆⋆ E-mail: [email protected]
On the other hand Gaidos, Henry & Henry (2000)
measured a Fe corrected Li-equivalent width ofW6708 =
3.9 ± 1.9mÅ and a rotational velocity of
v sin i ≈ 4.0 km/s, which might be too small for a
100Myr old star. Furthermore Chen et al. (2005) ob-
served HD141272 using the infrared space telescope
Spitzer and did not find any IR-excess at 24µm and
70µm indicating that HD141272 is not surrounded by
an optically thick disk.
Finally López-Santiago et al. (2006) revised the list
of Her-Lyr members and candidates of Fuhrmann (2004)
and classified HD141272 as an doubtful member, due
to its lithium depletion.
In our program we search for companions to Her-
Lyr members and candidates and first results are pre-
sented here. We found a co-moving companion of
HD 141272 by a combination of archival first epoch im-
ages and recent observations. We present our imaging,
the astrometric data and reduction techniques in sec-
tion 2 and 3, followed by a description of the spectro-
scopic and photometric analysis of the new companion
in section 4. The results are discussed in section 5.
2 Archival first epoch data
Astrometry is an effective method to find companions
of stars, by comparing two images taken with suffi-
ciently long epoch difference. In order to find late-type
stellar and substellar objects, we concentrate our search
c© Year of publication WILEY-VCH Verlag GmbH&Co. KGaA, Weinheim
http://arxiv.org/abs/0704.0387v1
Astron. Nachr. / AN (Year of publication) 1
Fig. 1 POSS-I E image of HD141272 from 17 June
1950. The star is located at αJ2000.0 = 15
h 48m 09.4s,
δJ2000.0 = +01
◦ 34′ 18′′. A faint object is located in the
north of HD141272, which is hardly recognizable due
to the diffraction spikes of the primary star induced by
saturation. With a pixel size of ∼ 10microns the pixel
scale of the plate is ∼ 6.72 arcsec/pixel.
on companions of young stars. Young objects are still
in contraction and are brighter than older objects of
the same mass hence, low mass objects are easier to
detect.
We found HD141272 in three epochs of the Super-
COSMOS-Sky-Survey, namely a POSS-I (Palomar Ob-
servatory Sky Survey) plate from 1950, as well as in
UKST (United Kingdom Schmidt Telescope) infrared
and red observations from 1981 and 1992. On all three
plates we detected by eye inspection a faint object,
located approximately 18 arcsec north of HD141272,
which was not detected by the SuperCOSMOS ma-
chine due to its small angular separation to the much
brighter star and due to its overlap with the diffraction
spike (Fig. 1).
The diffraction spike of HD141272 intersects the
northern object on all three plates hence, the detec-
tion of this object would be inaccurate by means of
most common detection techniques. Nevertheless, we
obtained a position measurement of the companion can-
didate on the POSS-I plate, using the Source Extractor
package (Bertin & Arnouts 1996), included in the Star-
link application GAIA (Gray et al. 2004). The source
extractor uses thresholding and deblending of point-
spread functions hence the method is more accurate
than other detection techniques (e.g. Gaussian fitting)
under the circumstances in Fig. 1. However, an system-
atical error is possible, due to the perturbation of the
primary’s spike. This error is larger in right accession
than in declination and would affect the measurement
of the position angle rather than the separation (see
section 3, Fig. 4), due to the orientation of the system
(Fig. 1 and 3).
Due to its brightness HD141272 saturates the POSS-
I plate. Furthermore the PSF (point spread function)
is contaminated by the stray light of the companion
candidate hence, position measurement via PSF cen-
tering does not work sufficiently. We used the diffrac-
tion spikes of the saturated primary to determine its
position, since they are unaffected by the companion.
We determined the intensity center of a spike taking
∼ 30 measurements for each spike using the data re-
duction and analysis package ESO-MIDAS. The appli-
cation of a linear regression gives the position of the
star as intersection of the two spikes and leads to very
small astrometric uncertainties (∆αH = 0.047 arcsec
and ∆δH = 0.050 arcsec).
In addition to the detection on the POSS-I plate
HD 141272 and its companion-candidate are also de-
tected in 2MASS images from observing epoch 2000.
The 2MASS point source catalog (Cutri et al. 2003)
lists the position of both objects with accurate astro-
metric precision, see Tbl. 1.
Equipped with these data we determined the proper
motion of all stars in a 15 arcmin box around HD141272
which are detected at the POSS-I plate and listed in
the 2MASS point source catalog (see Fig. 2). We de-
rived the proper motion of all stars in the field by
comparing the positions of all detected objects. The
majority of sources only shows small proper motion
following a normal distribution, since these stars are
most probably at high distances. Using the Lilliefors
test for normal distribution we derived the subsample
of stars belonging to the background stars, since their
proper motion follows a normal distribution (non mov-
ing background stars). The standard deviation of the
background stars gives the statistically derived proper
motion error (σp.m., α = 8.8mas/yr, σp.m., δ = 6.8
mas/yr). Objects not belonging to the background stars
are considered as companion candidates, if they are ly-
ing within a 5-σ vicinity of HD141272 (ellipse in Fig. 2).
Other objects are omitted, since these are either false
detections or high-proper-motion stars moving in other
directions.
The proper motion of the nearby star HD141272
is clearly separated from the background stars. The
companion candidate clearly shares the proper motion
of HD 141272 and will be denoted HD141272B, here-
after. Fig. 2 shows with high confidence (∼ 13σ) that
HD141272A and B are co-moving over roughly 50 years.
Due the above discussed astrometric uncertainties of
HD 141272B this analysis gives a first indication of a
new nearby young double star system
www.an-journal.org c© Year of publication WILEY-VCH Verlag GmbH&Co. KGaA, Weinheim
2 T. Eisenbeiss et al.: Low mass visual binaries in the solar neighbourhood: The case of HD141272
−300−250−200−150−100−50050
∆ Ra/Year [mas/yr]
standard error
search radius
HD141272
companion
candidate
non moving background stars
Fig. 2 Proper motion plot of HD141272 (cross) and
its companion candidate (circle) and non moving back-
ground stars (upper left). X- and Y-axis show the
change of the positions (in mas/yr). The plot is based
on POSS-I Schmidt plate (17 June 1950) and 2MASS
catalog data (29 April 2000). Error estimates are taken
as 2-σ errors from the background stars. Data points
lying outside the background stars and outside a 5-σ
vicinity of HD141272 (large ellipse) are omitted, since
these are either false detections or high proper motions
stars moving in other directions. The statistical error
of all data points is shown by the thick error cross in
the lower left. The diagram shows the common proper
motion of HD141272 and its new companion with a
confidence of ∼ 13σ.
Moreover, we used the non-moving background stars
to estimate the positional error of the detections in the
POSS-I plate. The mean of the distribution shows the
systematic error of the POSS-I measurements (∆sys, α =
−4.5mas/yr and ∆sys, δ = −4.9mas/yr as offset to
(0, 0). The whole set of data points in Fig. 2 is shifted
by that offset to correct for calibration errors between
POSS-I and 2MASS data. The standard deviation shows
the statistical measurement error (∆stat = σp.m.) hence,
can be applied as standard detection error. The total
detection error derived for the POSS-I plate is ∆α =
0.29 arcsec and ∆δ = 0.25 arcsec. The additional sys-
tematic error for the companion candidate due to the
diffraction spike of HD141272 is not included in this
error analysis.
Fig. 3 H-band image of HD 141272 and its compan-
ion candidate taken with the near infrared camera Ω-
Cass at the 3.5m telescope of the Calar Alto obser-
vatory in Spain. The separation between HD141272
and its companion candidate is ∼ 17.8 arcsec at a
position angle of ∼ 352.62◦ with a pixel scale of ∼
0.2 arcsec/pixel. Note that HD141272 is slightly satu-
rated.
3 Follow-up observations
In order to get a third epoch on our astrometric re-
sult and to detect or rule out further companions we
observed HD141272 again in April 2006 (Fig. 3). We
carried out H-band as well as narrow-band observations
(1.644µm) with the near infrared camera Ω-Cass, in-
stalled at the Cassegrain focus of the 3.5m telescope of
the Calar Alto observatory in Spain. Ω-Cass is equipped
with a 1024 × 1024 HgTeCd-detector with a pixel scale
of ∼0.2 arcsec per pixel. We always used the short-
est possible detector integration time (0.84 s) to limit
strong saturation effects due to the bright star. For
background subtraction we applied the standard jitter
technique and chose 12 jitter-positions. On each jitter
position 49 integrations (0.84 s) were co-added, yield-
ing a total integration time in the H-band of 8.2min.
All images were flatfielded with a skyflat image taken
during twilight. The whole data reduction (background
subtraction, flatfielding, and shift+add) was carried out
with the ESO data-reduction package Eclipse (Devil-
lard 2001).
We calibrated our Ω-Cass image for relative astrom-
etry, using the well known binary systems HIP 63322
and HIP 82817, which we observed during the same
night and with the same instrumentation as our sci-
ence image. Using the Hipparcos astrometry (Perry-
c© Year of publication WILEY-VCH Verlag GmbH&Co.KGaA, Weinheim www.an-journal.org
Astron. Nachr. / AN (Year of publication) 3
Table 1 Separation and position angle of the co-moving companion HD141272B relative to its primary
HD141272A for all observing epochs. We also show the expected change of separation and position angle in
case that the companion is a non-moving background source, derived with the well known proper and parallactic
motion of the primary.
epoch telescope/ pixel scale band sepobs. sepifback PAobs. PAifback
[dd/mm/yyy] catalogs [arcsec] [arcsec] [arcsec] [◦] [◦]
17/07/1950 POSS-I 1.0 E (6442Å) 17.85±0.31 − 353.6±1.1 −
29/04/2000 2MASS 0.7 JHKS 17.83±0.150 26.92±0.33 352.42±0.48 14.61±0.75
20/04/2006 3.5m CA 0.2 H 17.851±0.041 28.12±0.31 352.62±0.18 16.48±0.68
man et al. 1997) and considering the maximal orbital
motion of the calibration binaries we estimated the
pixel scale (192 ± 0.43mas/pixel) and the orientation
(−1.86±0.18 ◦) of the Ω-Cass images. This yields to the
relative astrometric parameters of the system (Tbl. 1).
For the detection of both objects we used the Gaussian
centroiding technique, implemented in ESO-MIDAS.
Further co-moving companions could be ruled out
around HD141272 within an angular separation of ∼ 5
to 73 arcsec (1500AU of projected separation) with H-
band magnitudes down to 18.3mag (S/N= 3).
HD 141272A and B are separated by ∼ 17.8 arcsec
(Fig. 3), hence the projected separation of the system
is approximately 380AU and its orbital period can be
estimated with Kepler’s third law to be roughly 7000
years (we use 0.83M⊙ for HD 141272A and 0.26M⊙
for B). During 56 years of epoch difference between
the POSS-I and our H-band observation, this yields
maximal orbital motion as large as ∼0.5 arcsec in sep-
aration (edge-on orbit assumed) or ∼3◦ in position an-
gle (face-on orbit assumed). Therefore, we derived the
separation and the position angle of the companion for
all three observing epochs which are summarized in
Tab. 1. These results are also visualized in Fig. 4. Note
that absolute calibrated astrometric data, derived for
the POSS-I image as described in section 2, as well as
catalog data from the 2MASS catalog is used in Fig. 4,
while the third epoch data is based on relative astrom-
etry, hence the uncertainties of that data point are sig-
nificantly smaller.
While the separation between HD141272A and B
did not change during 56 years, we found a slight de-
crease of its position angle. This effect is most likely
due to the perturbation of the companions PSF by the
diffraction spike of the primary (see section 2 and Fig.
1). Nevertheless Fig. 4 ensures the companionship of
HD 141272B, since all data points are lying within the
given error bars of the first epoch.
4 Photometry and Spectroscopy
The infrared colors of both components of the new bi-
nary system HD141272AB are listed in the 2MASS
point source catalog, i.e. accurate J, H, and KS band
2.43 2.44 2.45 2.46
JD−2400000.5
2.43 2.44 2.45 2.46
JD−2400000.5
Fig. 4 Separation (sep) and position angle (PA) for
HD 141272A and B from 1950 to 2006 (three data
points). Upper lines show the changes of the proper-
ties under the assumption HD141272B was a back-
ground star (including parallactic motion of A) while
the straight, opening lines give the range of the bi-
nary movement, considering maximal orbital move-
ment. While the separation stays approximately con-
stant there is a change in the position angle, caused
by the perturbation of the companions PSF due to the
diffraction spikes of the primary.
Table 2 2MASS Photometry of HD141272A and B
Comp. J H KS
[mag] [mag] [mag]
A 5.991±0.021 5.610±0.027 5.501±0.018
B 9.298±0.020 8.725±0.055 8.456±0.023
photometry is available for the primary and its co-
moving companion, which is summarized in Tab. 2. Ad-
ditionally the I-band magnitude of both components
(mI = 8.59 ± 0.02mag for A and mI = 10.572 ±
0.02mag for B) is measured in the second release of the
DENIS database, while the accuracy for HD141272A
is limited due to saturation effects, hence the given er-
ror is probably underestimated.
In order to obtain also unsaturated images of the
primary we observed the binary system with Ω-Cass
www.an-journal.org c© Year of publication WILEY-VCH Verlag GmbH&Co. KGaA, Weinheim
4 T. Eisenbeiss et al.: Low mass visual binaries in the solar neighbourhood: The case of HD141272
in the FeII (1.644µm) narrow-band filter. Thereby, we
used again the 12 point jitter pattern but co-added 15
integrations (4 s) per jitter position, yielding a total in-
tegration time of 12min. The bright primary as well as
its fainter co-moving companion are both well detected
in this narrow-band image and their fluxes did not ex-
ceed the linearity level of the Ω-Cass detector. Hence,
we could use this image to derive the magnitude differ-
ence between the primary star and its companion and
obtained ∆HFeII = 3.166± 0.005mag, fully consistent
with the magnitude difference derived from the 2MASS
data in H-band (∆H = 3.115± 0.061mag.)
Furthermore we acquired a low-resolution optical
spectrum with EMMI at the NTT on La Silla to de-
termine the spectral type of HD141272B and prove its
common distance with HD141272A. The spectrum was
taken in RILD and REMD mode covering a wavelength
of 400-900nm with a resolution of R ≈ 3000 at 600nm.
The data reduction followed the standard procedure for
low-resolution optical spectra: After bias subtraction,
flat fielding and wavelength calibration with a HeAr arc
spectrum we corrected for the instrumental response
and for telluric features using a spectrum of HR5501
taken at the same airmass as HD141272B.
We determined the spectral type by comparing our
spectrum with a standard sequence of M dwarfs in the
same spectral range and with comparable spectral res-
olution (Bochanski et al. 2006), see Fig. 5. The best
fit resulted in a spectral type of M3.25 ± 0.25 which
is consistent with a spectral type of M3.0± 0.5 deter-
mined from the TiO5 spectral index of 0.49 following
Cruz & Reid (2002).
Adopting the latter spectral type as final we derived
a spectrophotometric distance of 24.4±4.2 pc from the
MJ relation given in Cruz & Reid (2002) and the J
magnitude from 2MASS, assuming that the companion
is on the Main sequence. The determined distance is in
excellent agreement with the HIPPARCOS measured
distance of 21.35±0.48 pc for HD 141272A, confirming
their common distance. Hence, we call the companion
HD141272B.
5 Conclusions
With the astrometric data reduction and analysis tech-
niques presented in this work, we could verify the com-
mon proper motion of both components of the binary
system HD141272AB during 56 years of epoch differ-
ence between the first successful observation of this sys-
tem on the POSS-I plates taken in July 1950 and our
H-band imaging obtained with Ω-Cass in April 2006.
Furthermore we obtained an optical spectrum of the
companion and derived its spectral type to range be-
tween M2.5V and M3.5V. The infrared apparent mag-
nitudes of the co-moving companion are fully consis-
tent with a M3 dwarf which is located at the distance
400 500 600 700 800 900
HD 141272 B
Wavelength [nm]
Fig. 5 Relative flux of the spectral sequence from
M1 to M5 (Bochanski et al. 2006) in comparison to the
EMMI spectrum of HD141272B, ranging from 400 to
900 nm. The resolutions are comparable (R ∼ 3000 for
the EMMI spectrum and R ∼ 6000 for the standard
spectra at 600 nm). HD 141272B shows good agree-
ment with an M3 star.
of HD141272A which finally confirms the companion-
ship of this new binary system. The companion is an
addition to the Catalog of Nearby Stars within 25 pc
(Gliese & Jahreiß 1991).
In order to get an estimation of the system age we
compared the infrared photometry of HD 141272A and
B with ≈ 1300 members of the Pleiades cluster which
are listed in the WEBDA database (Mermilliod 1998).
All objects are plotted in a J-K vs. MH color-magnitude
diagram (Fig. 6). The colors of all objects are obtained
from the 2MASS catalog and we derived the absolute
H-Band magnitudes of all comparison stars using their
2MASS H-band photometry and a mean distance mod-
ule of the Pleiades of 5.97mag (WEBDA database).
The expected distance uncertainty of the cluster mem-
bers which results in an uncertainty of their absolute H-
band magnitudes was approximated with the angular
diameter of the Pleiades cluster on the sky, assuming a
similar extension of the cluster also in the radial direc-
tion. The absolute H-band magnitudes of HD141272A
and B are derived with 2MASS photometry and the
Hipparcos parallax of the binary system. Compared to
the Pleiades of the same J-K color HD141272A and
B appear a little fainter, indicating that the system is
already on the ZAMS, which is similar to the results of
earlier works (Gaidos 1998; Wright et al. 2004).
If we assume that both components of the binary
system have already reached the ZAMS we can deter-
mine the mass of the secondary using equation (11)
from Kirkpatrick & McCarthy (1994) with the given
c© Year of publication WILEY-VCH Verlag GmbH&Co.KGaA, Weinheim www.an-journal.org
Astron. Nachr. / AN (Year of publication) 5
−0.2 0 0.2 0.4 0.6 0.8 1 1.2 1.4 1.6 1.8
J−K [mag]
0.75 0.8 0.85 0.9 0.95 1
J−K [mag]
HD 141272 A
HD 141272 B
Fig. 6 J-K vs. MH diagram for the Pleiades and
HD141272A and B (rectangles symbolize the error
boxes). The inserted plot shows HD141272B and the
surrounding Pleiades stars drawn to a larger scale. The
main sequence of the cluster can be seen although there
are some outliers due to the mean distance module
(5.97mag for Pleiades) applied. The mean error of the
Pleiades is shown by the error cross in the lower left.
HD 141272A and B appear a little fainter than Pleiades
stars of the same J-K color. This indicates, that the
system already reached the ZAMS.
errors for the constants a and b and the range of the
spectral type. We derived a mass of
M∗ = 0.26
+0.07
−0.06M⊙.
Future work should ascertain the age of the system
and derive more properties of the M dwarf, which en-
larges the list of nearby low mass stars bound in binary
systems.
Acknowledgements. We would like to thank the technical
staff of the ESO NTT at La Silla as well as of the Centro
Astronómico Hispano Alemán (CAHA) at Calar Alto for all
their help and assistance in carrying out the observations.
In addition we would like to thank John Bochanski, An-
drew West, Suzanne Hawley and Kevin Covey for providing
the electronic sequence of M-stars composite spectra.
T.O.B. Schmidt acknowledges support from a Thur-
ingian State Scholarship and from a Scholarship of the Evan-
gelisches Studienwerk e.V. Villigst.
This publication makes use of data products from the
Two Micron All Sky Survey, which is a joint project of
the University of Massachusetts and the Infrared Process-
ing and Analysis Center/California Institute of Technology,
funded by the National Aeronautics and Space Administra-
tion and the National Science Foundation.
We use imaging data from the SuperCOSMOS Sky Sur-
vey, prepared and hosted by the Wide Field Astronomy
Unit, Institute for Astronomy, University of Edinburgh,
which is funded by the UK Particle Physics and Astron-
omy Research Council.
This research has made use of the VizieR catalogue
access tool and the Simbad database, both operated at
the Observatoire Strasbourg, as well as of the WEBDA
database, operated at the Institute for Astronomy of the
University of Vienna.
The DENIS project has been partly funded by the SCI-
ENCE and the HCM plans of the European Commission
under grants CT920791 and CT940627. It is supported by
INSU, MEN and CNRS in France, by the State of Baden-
Württemberg in Germany, by DGICYT in Spain, by CNR
in Italy, by FFwFBWF in Austria, by FAPESP in Brazil,
by OTKA grants F-4239 and F-013990 in Hungary, and by
the ESO C&EE grant A-04-046.
Jean Claude Renault from IAP was the Project man-
ager. Observations were carried out thanks to the contri-
bution of numerous students and young scientists from all
involved institutes, under the supervision of P. Fouqué, sur-
vey astronomer resident in Chile.
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http://irsa.ipac.caltech.edu
Introduction
Archival first epoch data
Follow-up observations
Photometry and Spectroscopy
Conclusions
|
0704.0388 | Sterile neutrinos at the CNGS | IFT-UAM/CSIC-07-16
Sterile neutrinos at the CNGS
Andrea Donini, a Michele Maltoni, a Davide Meloni, b
Pasquale Migliozzi, c Francesco Terranova d
aInstituto F́ısica Teórica UAM/CSIC, Cantoblanco, E-28049 Madrid, Spain
bI.N.F.N., Sezione di Roma I and Dip. Fisica, Univ. Roma “La Sapienza”,
Pl. A. Moro 2, I-00185, Rome, Italy
cI.N.F.N., Sezione di Napoli, I-80126, Naples, Italy
dI.N.F.N., Laboratori Nazionali Frascati, Via E. Fermi 40, I-00044, Frascati, Italy
PACS: 14.60.Pq, 14.60.Lm
Abstract
We study the potential of the CNGS beam in constraining the parameter space of a
model with one sterile neutrino separated from three active ones by an O(eV2) mass-
squared difference, ∆m2
. We perform our analysis using the OPERA detector as
a reference (our analysis can be upgraded including a detailed simulation of the
ICARUS detector). We point out that the channel with the largest potential to
constrain the sterile neutrino parameter space at the CNGS beam is νµ → ντ . The
reason for that is twofold: first, the active-sterile mixing angle that governs this
oscillation is the less constrained by present experiments; second, this is the signal
for which both OPERA and ICARUS have been designed, and thus benefits from
an extremely low background. In our analysis we also took into account νµ → νe
oscillations. We find that the CNGS potential to look for sterile neutrinos is limited
with nominal intensity of the beam, but it is significantly enhanced with a factor 2
to 10 increase in the neutrino flux. Data from both channels allow us, in this case,
to constrain further the four-neutrino model parameter space. Our results hold for
any value of ∆m2
& 0.1 eV2, i.e. when oscillations driven by this mass-squared
difference are averaged. We have also checked that the bound on θ13 that can be
put at the CNGS is not affected by the possible existence of sterile neutrinos.
http://arxiv.org/abs/0704.0388v2
1 Introduction
The results of solar [1,2,3,4,5,6], atmospheric [7,8], reactor [9,10,11,12] and accelera-
tor [13,14,15] neutrino experiments show that flavour mixing occurs not only in the
hadronic sector, as it has been known for long, but in the leptonic sector as well. The
full understanding of the leptonic mixing matrix constitutes, together with the dis-
crimination of the Dirac/Majorana character of neutrinos and with the measurement
of their absolute mass scale, the main goal of neutrino physics for the next decade.
The experimental results point to two very distinct mass-squared differences, ∆m2
7.9 × 10−5 eV2 and |∆m2
| ≈ 2.4 × 10−3 eV2. On the other hand, only two out of
the four parameters of the three-family leptonic mixing matrix UPMNS [16,17,18,19] are
known: θ12 ≈ 34
◦ and θ23 ≈ 43
◦ [20]. The other two parameters, θ13 and δ, are still
unknown: for the mixing angle θ13 direct searches at reactors [9,10,11] and three-family
global analysis of the experimental data give the upper bound θ13 ≤ 11.5
◦, whereas for
the leptonic CP-violating phase δ we have no information whatsoever (see, however,
Ref. [20]).
The LSND data [21,22,23], on the other hand, would indicate a ν̄µ → ν̄e oscillation with
a third neutrino mass-squared difference: ∆m2
∼ 0.3 − 6 eV2, about two orders of
magnitude larger than ∆m2
. Given the strong hierarchy among the solar, atmospheric
and LSND mass-squared splittings, ∆m2
≪ ∆m2
≪ ∆m2
, it is not possible to
explain all these data with just three massive neutrinos, as it has been shown with
detailed calculations in Ref. [24]. A necessary condition to explain the whole ensemble
of data in terms of neutrino oscillations is therefore the introduction of at least a
fourth light neutrino state. This new light neutrino must be an electroweak singlet [18]
in order to comply with the strong bounds on the Z0 invisible decay width [25,26]. For
this reason, the LSND signal has often been considered as an evidence of the existence
of a sterile neutrino.
In recent years, global analyses of solar, atmospheric, short-baseline [27,28,29,30] exper-
iments and LSND data have been performed to establish whether four-neutrino models
can really reconcile the data and solve the puzzle [31,32,33,34,35,36,37,38]. The point
is that providing a suitable mass-squared difference to each class of experiments is not
enough: it is also necessary to show that the intrinsic structure of the neutrino mixing
matrix is compatible with all the data. This turned out to be very hard to accomplish.
In Ref. [39] it was shown that four-neutrino models were only marginally allowed, with
best fit around ∆m2
≃ 1 eV2 and sin2 2θLSND ≃ 10
−3. Generically speaking, the
global analysis indicated that a single sterile neutrino state was not enough to reconcile
LSND with the other experiments. For this reason, to improve the statistical compat-
ibility between the LSND results and the rest of the oscillation data, models with
two sterile neutrino states have been tested (see, for example, Ref. [40] and references
therein). Although a slightly better global fit was achieved, a strong tension between
the LSND data and the results from atmospheric and short-baseline experiments was
still present.
So far, the LSND signal has not been confirmed by any other experiment [41]. It is
therefore possible that the LSND anomaly arises from some some yet unknown problem
in the data set itself. To close the issue, the MiniBooNE collaboration [42] at FermiLab
has recently performed a search for νµ → νe appearance with a baseline of 540 m and a
mean neutrino energy of about 700 MeV. The primary purpose of this experiment was
to test the evidence for ν̄µ → ν̄e oscillation observed at LSND with a very similar L/E
range. No evidence of the expected signal has been found, hence ruling out once and for
all the four-neutrino interpretation of the LSND anomaly. However, MiniBooNE data
are themselves not conclusive: although no evidence for νµ → νe oscillation has been
reported in the spectrum region compatible with LSND results, an unexplained excess
has been observed for lower energy neutrinos. Furthermore, within a five-neutrino model
this excess can be easily explained, and even reconciled with LSND and all the other
appearance experiments [43]. On the other hand, a post-MiniBooNE global analysis
including also disappearance data show that five-neutrino models suffer from the same
problems as four-neutrino schemes, and in particular they are now only marginally
allowed – a situation very similar to that of four-neutrino models before MiniBooNE
data. Adding a third sterile neutrino 1 does not help [43], and in general global analyses
seem to indicate that sterile neutrinos alone are not enough to reconcile all the data.
Models with sterile neutrinos and exotic physics have been therefore proposed (see, for
example, Ref. [46]).
In summary, the present experimental situation is still confused. It is therefore worth-
wile to understand if, aside of MiniBooNE, new neutrino experiments currently running
or under construction can investigate the existence of sterile neutrinos separated from
the active ones by O(eV2) mass-squared differences. In this paper we explore in detail
the capability of the CNGS beam to perform this search. For definiteness we focus on
the simplest case with only one extra sterile neutrino. Note that this model is perfectly
viable once the LSND result is dropped, as it contains as a limiting case the usual
three-neutrino scenario. Furthermore, it is easily generalizable by adding new sterile
neutrino states, and it can be used as a basis for models with extra “sterile” states
strongly decoupled from active neutrinos (such as in extra-dimensions models with a
right-handed neutrino in the bulk [47]).
The CNGS beam [48] has been built to test the (supposedly) dominant oscillation in at-
mospheric neutrino data, νµ → ντ . In order to make possible τ production through CC
interactions, the mean neutrino energy, 〈Eν〉 = 17 GeV, is much above the atmospheric
oscillation peak for the CERN to Gran Sasso baseline, L = 732 Km. Two detectors are
illuminated by the CNGS beam: OPERA (see Ref. [49] and refs. therein) will start data
taking with the lead-emulsion target in 2007; ICARUS-T600 (see Ref. [50] and refs.
therein) will start operating in 2008. Both detectors have been especially designed to
look for τ ’s produced through νµ → ντ oscillation and to minimize the corresponding
1 A quite interesting scenario is, in our opinion, that in which three right-handed Majorana
neutrinos are added to the three weakly-interacting ones. If the Majorana mass term M is
O(eV), (3+3) light Majorana neutrinos are present at low-energy [44,45].
backgrounds. The expected number of τ events after signal selection in an experiment
such as OPERA (after five years of data taking with nominal CNGS luminosity) is
O(10) events with O(1) background event.
At the CNGS distance and energy, neutrino oscillations mediated by an O(eV2) mass
difference will appear as a constant term in the oscillation probability. In four-neutrino
models, fluctuations induced by this term over the atmospheric νµ → ντ oscillation can
be as large as 100% for specific points of the allowed parameter space. This is due to the
fact that the leading angle for this oscillation is the less constrained one. The νµ → ντ
channel, therefore, is extremely promising as a “sterile neutrino” smoking gun, as it has
been commented elsewhere (see, for example, Refs. [51,52] and refs. therein). To test
the model we will also make use of the νµ → νe channel. Notice that the background to
this signal coming from τ → e decay is modified in four-neutrino models with respect
to standard three-family oscillations. In fact, since νµ → ντ oscillations are depleted by
active-sterile mixing with respect to standard ones, the τ → e background to νµ → νe
oscillations gets depleted, too. A combined analysis of the two channels in four-neutrino
models at the OPERA detector has been performed, taking into account properly all
of the backgrounds. We stress, however, that the same analysis could be performed
at ICARUS, as well. The previous considerations hold for any facility operating well
beyond the kinematical threshold for τ production.
In the specific case of the CNGS beam, the limited flux implies a modest improvement
in the parameter space exclusion, see Sec. 6. An increase in the exposure of such
facilities, however, would permit to improve the present bounds on the parameters of
four-neutrino models and, in particular, to constrain the leading angle in νµ → ντ
oscillations at the level of the other mixing parameters.
The paper is organized as follows. In Sec. 2 we briefly review the main features of
four-neutrino models and we introduce our parametrization of the mixing matrix. In
Sec. 3 we compute the vacuum oscillation probabilities in the atmospheric regime and
we review the present bounds on the active-sterile mixing angles. In Sec. 4 we recall the
most relevant parameters of CNGS. In Sec. 5 we study theoretically the expectations
of the νµ → ντ and νµ → νe channels at the CNGS. In Sec. 6 we present our results
using these channels at the OPERA detector and the CNGS beam. Finally, in Sec. 7
we draw our conclusions.
2 Four neutrino mass schemes
In four-neutrino models, one extra sterile state is added to the three weakly interacting
ones. The relation between the flavor and the mass eigenstates is then described by a
4×4 unitary matrix U , which generalizes the usual 3×3 matrix UPMNS [16,17,18,19]. As
stated in the introduction, in this work we only consider the case when the fourth mass
eigenstate is separated by the other three by an O(eV2) mass-squared gap. There are
six possible four-neutrino schemes, shown in Fig. 1, that can accommodate the results
(3+1) (2+2)
Fig. 1. The two classes of four–neutrino mass spectra, (3+1) and (2+2).
from solar and atmospheric neutrino experiments and contain a third much larger ∆m2.
They can be divided in two classes: (3+1) and (2+2). In the (3+1) schemes, there is
a group of three close-by neutrino masses that is separated from the fourth one by
the larger gap. In (2+2) schemes, there are two pairs of close masses separated by the
large gap. While different schemes within the same class are presently indistinguishable,
schemes belonging to different classes lead to very different phenomenological scenarios.
A characteristic feature of (2+2) schemes is that the extra sterile state cannot be
simultaneously decoupled from both solar and atmospheric oscillations. To understand
why, let us define
i∈ sol
|Usi|
2 and cs =
j ∈ atm
|Usj|
2 (1)
where the sums in i and j run over mass eigenstates involved in solar and atmospheric
neutrino oscillations, respectively. Clearly, the quantities ηs and cs describe the fraction
of sterile neutrino relevant for each class of experiment. Results from atmospheric and
solar neutrino data imply that in both kind of experiments oscillation takes place
mainly between active neutrinos. Specifically, from Fig. 46 of Ref. [20] we get ηs ≤ 0.30
and cs ≤ 0.36 at the 3σ level. However, in (2+2) schemes unitarity implies ηs + cs = 1,
as can be easily understood by looking at Fig. 1. These models are therefore ruled out
at a very high confidence level [53], and in the rest of this work we will not consider
them anymore.
On the other hand, (3+1) schemes are not affected by this problem. Although the
experimental bounds on ηs and cs quoted above still hold, the condition ηs + cs = 1 no
longer applies. For what concerns neutrino oscillations, (3+1) models are essentially
unfalsifiable, since they reduce to the conventional three-neutrino scenario when the
mixing between active and sterile states are small enough.
The mixing matrix U can be conveniently parametrized in terms of six independent
rotation angles θij and three (if neutrinos are Dirac fermions) or six (if neutrinos are
Majorana fermions) phases δi. In oscillation experiments, only the so-called “Dirac
phases” can be measured, the effect of the “Majorana phases” being suppressed by
factors of mν/Eν . The Majorana or Dirac nature of neutrinos can thus be tested only
in ∆L = 2 transitions such as neutrino-less double β-decay [54] or lepton number
violating decays [25]. In the following analysis, with no loss in generality, we will restrict
ourselves to the case of 4 Dirac-type neutrinos only.
A generic rotation in a four-dimensional space can be obtained by performing six dif-
ferent rotations along the Euler axes. Since the ordering of the rotation matrices Rij
(where ij refers to the plane in which the rotation takes place) is arbitrary, plenty of
different parametrizations of the mixing matrix U are allowed. The large parameter
space (6 angles and 3 phases, to be compared with the standard three-family mixing
case of 3 angles and 1 phase) is however reduced to a subspace whenever some of
the mass differences become negligible. If the eigenstates i and j are degenerate, ro-
tations in the ij-plane become unphysical and the corresponding mixing angle should
drop from oscillation probabilities. If the matrix Rij is the rightmost one the angle
θij automatically disappears, since the matrix commutes with the vacuum hamilto-
nian. The parameter space gets therefore reduced to the physical angles and phases.
If a different ordering of the rotation matrices is taken, no angle explicitly disappears
from the oscillation formulas, but the physical parameter space is still smaller than the
original one. In this case, a parameter redefinition is needed to reduce the parameter
space to the observable sector. In Refs. [55,56] it was shown how the one-mass domi-
nance (∆sol → 0 and ∆atm → 0, where ∆ = ∆m
2L/4E [57]) and two-mass dominance
(∆sol → 0) approximations can be implemented in a transparent way (in the sense that
only the physical parameters appear in oscillation probabilities) using a parametriza-
tion in which rotations are performed in the planes corresponding to smallest mass
difference first:
USBL = R14(θ14) R24(θ24) R34(θ34) R23(θ23, δ3) R13(θ13, δ2) R12(θ12, δ1) . (2)
This parametrization was shown to be particularly useful when maximizing oscilla-
tions driven by a O(eV2) mass difference. The analytical expressions for the oscillation
probabilities in the (3+1) model in the one-mass dominance approximation in this
parametrization have been presented in Ref. [51].
In this paper, however, we are interested in a totally different regime: the “atmospheric
regime”, with oscillations driven by the atmospheric mass difference, ∆m2
L/E ∼
π/2. We will then make use of the following parametrization, adopted in Ref. [43]:
Uatm = R34(θ34) R24(θ24) R23(θ23, δ3) R14(θ14) R13(θ13, δ2) R12(θ12, δ1) . (3)
It is convenient to put phases in R12 (so that it automatically drops in the two-mass
dominance regime) and R13 (so that it reduces to the “standard” three-family Dirac
phase when sterile neutrinos are decoupled). The third phase can be put anywhere;
we will place it in R23. Note that in the one-mass dominance regime all the phases
disappear from the oscillation probabilities.
0 2 4 6 8 10 12 14 16 18 20
0 2 4 6 8 10 12 14 16 18 20
Fig. 2. Allowed regions at 90%, 95%, 99% and 3σ CL in the (θ13, θ14) plane (left) and in the
(θ24, θ34) plane (right) from the results of present atmospheric, reactor and LBL neutrino
experiments. The undisplayed parameters θ23 and δ3 are marginalized.
3 Oscillation probabilities and allowed parameter space
Let us first consider νe disappearance at L/E such that ∆sol can be safely neglected
with respect to ∆atm and ∆sbl. We get for this probability in vacuum:
Pee = 1− sin
2 2θ14 sin
∆sblL
sin2 2θ13 sin
∆atmL
, (4)
where cij = cos θij and sij = sin θij . It is clear from Eq. (4) that reactor experiments
such as Bugey and Chooz can put stringent bounds to θ13 and θ14, in this parametriza-
tion. This is depicted in Fig. 2(left), where 90%, 95%, 99% and 3σ CL contours in
the (θ13 − θ14)-plane are shown for ∆sol → 0 and ∆m
= 2.4 × 10−3 eV2. The third
mass difference, ∆m2
, is free to vary above 0.1 eV2. Notice that the νe disappearance
probability does not depend on θ23, θ24 and θ34. It can be clearly seen that the three-
family Chooz bound on θ13 is slightly modulated by θ14. Both angles, however, cannot
be much larger than 10◦. We will therefore expand in these two parameters to deduce
the other relevant oscillation probabilities.
At the CNGS beam atmospheric oscillations are large, solar oscillations can be ne-
glected and O(eV2) oscillations are extremely fast and can be averaged. It is useful to
write down the oscillation probability (in vacuum) at typical atmospheric L/E, in the
approximation ∆sol → 0, ∆sbl → ∞. In this regime:
P (να → νβ) = δαβ − 4ℜ
β3 (δαβ − U
α3Uβ3 − U
α4Uβ4)
δαβUα4U
β4 − |Uα4|
2|Uβ4|
α4Uβ4
sin∆23L ,
where + stands for neutrinos and − for antineutrinos, respectively. Up to second order
in θ13 and θ14 we get for the νµ disappearance oscillation probability:
Pµµ = 1− 2c
(1− c2
)− s2
− 2c3
s23(1− 2s
)s13s14s24 cos(δ2 − δ3)
∆atmL
A “negative” result in an atmospheric L/E νµ disappearance experiment (such as, for
example, K2K), in which νµ oscillations can be very well fitted in terms of three-family
oscillations, will put a stringent bound on the mixing angle θ24. The bound from such
experiments on θ24 can be seen in Fig. 2(right), where 90%, 95%, 99% and 3σ CL
contours in the (θ24 − θ34)-plane are shown for ∆sol → 0 and ∆m
= 2.4× 10−3 eV2.
The third mass difference, ∆m2
, is free to vary above 0.1 eV2. The mixing angles not
shown have been fixed to: θ23 = 45
◦; θ13 = θ14 = 0 (in this hypothesis, Pµµ does not
depend on phases). Notice that the νµ disappearance probability does not depend on
From the figure, we can see that θ24 cannot be much larger than 10
◦, either. We will
consider, therefore, the three mixing angles θ13, θ14 and θ24 being of the same order and
expand in powers of the three. At second order in θ13, θ14 and θ24, we get:
Pµµ = 1− 2s
− 4s2
(1− 2s2
)− s2
∆atmL
. (7)
Since both νe and νµ disappearance do not depend on θ34, we should ask which mea-
surements give the upper bound to this angle that can be observed in Fig. 2(right).
This is indeed a result of indirect searches for νµ → νs conversion in atmospheric exper-
iments, using the different interaction with matter of active and sterile neutrinos. This
can be understood from the (vacuum) νµ → νs oscillation probability at atmospheric
L/E for which, at second order in θ13, θ14 and θ24, we get:
Pµs = 2c
sin2 2θ23(c
+ 2c34 sin 2θ23s34
s24(1− 2s
) cos δ3 + 2s23s13s14 cos δ2
∆atmL
± c34 sin 2θ23s24s34 sin δ3 sin∆atmL .
As it can be seen, the bound on θ34 arises from a measurement of spectral distortion
(i.e., from the “atmospheric” term proportional to sin2∆atmL/2). On the other hand,
bounds on θ13, θ14 and θ24 are mainly drawn by a flux normalization measurement.
As a consequence, the bound on θ34 that we can draw by non-observation of νµ → νs
oscillation in atmospheric experiments is less stringent than those we have shown before.
For this reason, θ34 can be somewhat larger than θ13, θ14 and θ24, but still bounded to
be below 40◦. In the following, we will expand in powers of the four mixing angles
θ13, θ14, θ24 and θ34, that will be considered to be comparably small.
Up to fourth-order in θ13, θ14, θ24 and θ34, the νµ → νe appearance probability in the
atmospheric regime is:
Pµe = 4
[1− s2
] + s23s13s14s24 cos(δ2 − δ3)
∆atmL
± 2s23s13s14s24 sin(δ2 − δ3) sin∆atmL+ 2s
Eventually, the νµ → ντ appearance probability up to fourth-order in θ13, θ14, θ24 and
θ34 in the atmospheric regime is:
Pµτ = 2s
sin2 2θ23[c
− 4 sin 2θ23s13s14[s23s34 cos δ2 + c23s24 cos(δ2 − δ3)]
+ 2 sin 2θ23s24s34c
c34[c
− 2c2
] cos δ3
∆atmL
∓ sin 2θ23s24s34c
c34 sin δ3 sin∆atmL .
As it was shown in Refs. [51,52], the νµ → ντ appearance channel is a good place to
look for sterile neutrinos. This can be understood as follows: consider the νµ → ντ
three-family oscillation probability in the atmospheric regime, up to fourth-order in
P 3νµτ = Pµτ (θi4 = 0) ≃ c
sin2 2θ23 sin
∆atmL
. (11)
The genuine active-sterile neutrino mixing effects are:
∆Pµτ ≡ Pµτ − P
) sin2 2θ23 + 2 sin 2θ23(1− 2s
)s24s34 cos δ3
∆atmL
∓ sin 2θ23s24s34 sin δ3 sin∆atmL+ . . .
that is second-order in small angles θ13, θ14, θ24 and θ34. We would get a similar result
for νµ disappearance, also. On the other hand, computing the corresponding quantity
in the νµ → νe channel, we get:
∆Pµe ≡ Pµe − Pµe(θi4 = 0)
= s23s13s14s24 cos(δ2 − δ3) sin
∆atmL
± 2s23s13s14s24 sin(δ2 − δ3) sin∆atmL+ . . .
that is third-order in the same parameters.
Notice, eventually, that all oscillation probabilities start with an energy-independent
term and are, therefore, non-vanishing for L = 0, a result of our assumption that
∆sbl → ∞.
0 10 20 30 40 50
Eν(GeV)
anti−νµ
anti−νe
Fig. 3. CNGS neutrino fluxes (in arbitrary units) as a function of the neutrino energy. Both
muon and electron neutrino fluxes are illustrated.
4 The CNGS facility
The CNGS is a conventional neutrino beam in which neutrinos are produced by the
decay of secondary pions and kaons, obtained from collisions of 400 GeV protons from
the CERN-SPS onto a graphite target. The resulting neutrinos are aimed to the under-
ground Gran Sasso Laboratory (LNGS), located at 730 Km from CERN. This facility
provided the first neutrinos in August 2006 [49]. Differently from other long baseline
experiments, the neutrinos from CNGS can be exploited to search directly for νµ → ντ
oscillations, since they have a mean energy well beyond the kinematic threshold for τ
production. Moreover, the prompt ντ contamination (mainly fromDs decays) is negligi-
ble. The expected νe contamination is also relatively small compared to the dominant νµ
component, thus allowing for the search of sub-dominant νµ → νe oscillations through
an excess of νe CC events.
The energy spectra of the CNGS neutrino beam are shown (in arbitrary units) in
Fig. 3 [58]. In the present paper we assume the nominal intensity for the CNGS,
corresponding to 4.5× 1019 pot/year.
OPERA has been designed to search for τ appearance through identification of the
ντ CC interaction on an event-by-event basis. In particular, τ ’s are tagged identifying
explicitly their decay kink through high resolution nuclear emulsions interleaved with
lead sheets. For this detector, we can take advantage of the detailed studies of the
νµ → ντ signal (see Ref. [59]) and of the νµ → νe signal (see Ref. [60]).
The total non-oscillated CC event rates for a 1 Kton lead target with a neutrino flux
νµ ν̄µ νe ν̄e
669.0 13.7 5.9 0.3
Table 1
Nominal performance of the CNGS reference beam [58]. The total non-oscillated CC event
rates are calculated assuming 1019 pot and 1 Kton lead target mass.
normalized to 1019 pot are shown in Tab. 1 and are evaluated according to
dφνα(E)
σνα(E) dE , (14)
in which φνα is the flux of the neutrino flavour να and σνα the corresponding cross
section on lead.
5 Appearance channels at the CNGS
5.1 νµ → ντ oscillations
Since the CNGS experiments have been designed to search for νµ → ντ oscillation
in the parameter region indicated by the atmospheric neutrino data, we can take full
advantage of them in order to constrain (and, possibly, study) the four-family parameter
space.
The number of taus from νµ → ντ oscillations is given by the convolution of the νµ
flux dφνµ(E)/dE with the ντ charged-current cross-section on lead, σ
(E), weighted
by the νµ → ντ oscillation probability, Pµτ (E), times the efficiency for the OPERA
detector, εµτ :
Nµτ = A
dφνµ(E)
Pµτ (E)σ
(E) εµτ dE . (15)
A is a normalization factor which takes into account the target mass and the nor-
malization of the νµ flux in physical units. Specializing our analysis for the OPERA
detector, we have considered an overall efficiency εµτ ∼ 13%, [59]. This efficiency takes
into account that OPERA is able to exploit several decay modes of the final state τ ,
using both so-called short and long decays.
The dominant sources of background for the νµ → ντ signal are charm decays and
hadronic reinteractions. Both of them only depend on the total neutrino flux and not
on the oscillation probabilities. The OPERA experiment at the CNGS beam has been
designed precisely to measure this channel, and thus the corresponding backgrounds
are extremely low.
In Tab. 2 we report the expected number of τ events in the OPERA detector, according
to Eq. (15), for different values of θ13, θ14, θ24 and θ34. Input points have been chosen
according to the allowed regions in the parameter space shown in Sec. 3. The other
(θ13; θ14; θ24; θ34) Nτ background (θ13; θ14; θ24; θ34) Nτ background
(5◦; 5◦; 5◦; 20◦) 8.9 1.0 (10◦; 5◦; 5◦; 20◦) 8.5 1.0
(5◦; 5◦; 5◦; 30◦) 6.9 1.0 (10◦; 5◦; 5◦; 30◦) 6.5 1.0
(5◦; 5◦; 10◦; 20◦) 8.3 1.0 (10◦; 5◦; 10◦; 20◦) 7.9 1.0
(5◦; 5◦; 10◦; 30◦) 10.5 1.0 (10◦; 5◦; 10◦; 30◦) 10.3 1.0
3 families 15.1 1.0 3 families 14.4 1.0
Table 2
Event rates and expected background for the νµ → ντ channel in the OPERA detector, for
different values of θ14, θ24 and θ34 in the (3+1) scheme. The other unknown angle, θ13 has
been fixed to: θ13 = 5
◦, 10◦. The CP-violating phases are: δ1 = δ2 = 0; δ3 = 90
◦. As a
reference, the expected value in the case of standard three-family oscillation (i.e., for θi4 = 0)
is shown for maximal CP-violating phase δ. The rates are computed according to Eq. (15).
parameters are: θ12 = 34
◦; θ23 = 45
◦; ∆m2
= 7.9× 10−5 eV2; ∆m2
= 2.4× 10−3 eV2
and ∆m2
= 1 eV2 (all mass differences are taken to be positive). Eventually, phases
have been fixed to: δ1 = δ2 = 0; δ3 = 90
◦. The expected background is also shown. Rates
refer to a flux normalized to 4.5× 1019 pot/year (the nominal intensity of the CNGS),
an active lead target mass of 1.8 Kton and 5 years of data taking. For comparison, we
also report the expected number of events in the usual 3-family scenario.
As it can be seen, in most part of the parameter space we expect a significant depletion
of the signal with respect to standard three-neutrino oscillations. However, the differ-
ence between (3+1) model νµ → ντ oscillations and standard ones is much bigger than
the expected background. A good signal/noise separation can therefore be used to test
the model.
5.2 νµ → νe oscillations
The number of electrons from the νµ → νe oscillation is given by the convolution of
the νµ flux dφνµ(E)/dE with the νe charged-current cross-section on lead, σ
weighted by the νµ → νe oscillation probability, Pµe(E), times the efficiency for the
OPERA detector, εµe(E) [60]:
Nµe = A
dφνµ(E)
Pµe(E)σ
(E) εµe(E) dE , (16)
where A is defined as above. The overall signal efficiency εµe is the convolution of the
kinematic efficiency εkinµe (that ranges from 60% to 80% for neutrino energies between
5 to 20 GeV) and several (nearly factorizable) contributions. Among them, the most
relevant are trigger efficiencies, effects due to fiducial volume cuts, vertex and brick
finding efficiencies and the electron identification capability. They result in a global
constant factor εfactµe ∼ 48%.
(θ13; θ14; θ24; θ34) Ne ν
µ τ → e ν
(5◦; 5◦; 5◦; 20◦) 3.5 19.4 5.3 2.8 0.9
(5◦; 5◦; 5◦; 30◦) 3.5 19.4 5.3 2.1 0.9
(5◦; 5◦; 10◦; 20◦) 2.4 19.4 5.3 2.3 0.9
(5◦; 5◦; 10◦; 30◦) 2.4 19.4 5.3 2.4 0.9
3 families 3.7 19.7 5.3 4.6 0.9
(10◦; 5◦; 5◦; 20◦) 10.6 19.4 5.3 2.7 0.9
(10◦; 5◦; 5◦; 30◦) 10.4 19.4 5.3 2.0 0.9
(10◦; 5◦; 10◦; 20◦) 8.8 19.4 5.3 2.2 0.9
(10◦; 5◦; 10◦; 30◦) 8.6 19.4 5.3 2.4 0.9
3 families 15.1 19.7 5.3 4.8 0.9
Table 3
Event rates and expected background for the νµ → νe channel in the OPERA detector, for
different values of θ14, θ24 and θ34 in the (3+1) scheme. The other unknown angle, θ13, has
been fixed to: θ13 = 5
◦, 10◦. The CP-violating phases are: δ1 = δ2 = 0; δ3 = 90
◦. As a
reference, the expected value in the case of standard three-family oscillation(i.e., for θi4 = 0)
is shown for maximal CP-violating phase δ. The rates are computed according to Eq. (16).
Backgrounds have been computed following Ref. [60].
The dominant sources of background to the νµ → νe signal are, in order of importance:
(1) νe beam contamination;
(2) fake electrons due to π0 decays from νµ NC interactions;
(3) electrons produced through τ decay, where the τ comes from νµ → ντ oscillations;
(4) CC νµ events where the muon is lost and a track mimics an electron.
Backgrounds (1), (2) and (4) depend very little on the oscillation parameters. On the
other hand, the τ → e background depends strongly on the active-sterile mixing angles.
As we have seen in Sec. 5.1, in the allowed region of the parameter space νµ → ντ
oscillations are significantly depleted with respect to the standard three-neutrino ones.
As a consequence, this background gets depleted, too.
In Tab. 3 we report the expected number of electrons in the OPERA detector, according
to Eq. (16), for different values of θ13, θ14, θ24 and θ34. Input points have been chosen
according to the allowed regions in the parameter space shown in Sec. 3. The other
parameters are: θ12 = 34
◦; θ23 = 45
◦; ∆m2
= 7.9× 10−5 eV2; ∆m2
= 2.4× 10−3 eV2
and ∆m2
= 1 eV2 (all mass differences are taken to be positive). Eventually, phases
have been fixed to: δ1 = δ2 = 0; δ3 = 90
◦. Backgrounds have been computed accordingly
to Ref. [60]. Rates refer to a flux normalized to 4.5×1019 pot/year (the nominal intensity
of the CNGS), an active lead target mass of 1.8 Kton and 5 years of data taking. For
comparison, we also report the expected number of events in the usual 3-family scenario.
As it can be seen from Tab. 3, the difference between the (3+1) model and the stan-
dard three-neutrino oscillations are smaller in this channel than in the νµ → ντ one.
Moreover, they linearly depends on θ13, as it is clear from Eq. (13). For θ13 = 5
◦, this
channel will be of no help to test the allowed parameter space of the (3+1) model. On
the other hand, for θ13 saturating the Chooz-Bugey bound, both νµ → ντ and νµ → νe
might cooperate. However, notice that backgrounds to this signal are much larger than
the difference between (3+1) model and standard three-neutrino oscillations for any
value of θ13.
6 Sensitivity to (3 + 1) sterile neutrinos at OPERA
In this section we study the sensitivity to θ13 and to the active-sterile mixing angles
θ14, θ24 and θ34 at the CNGS beam, using both the νµ → ντ and νµ → νe appearance
channels at the OPERA detector. In the rest of this section, the known three-family
subspace angles have been fixed to: θ12 = 34
◦; θ23 = 45
◦. The mass differences have
been fixed to: ∆m2
= 7.9× 10−5 eV2 and ∆m2
= 2.4× 10−3 eV2. The CP-violating
phases δ1 and δ2 have been kept fixed to δ1 = δ2 = 0. On the contrary, the CP-violating
phase δ3 is fixed to two values: δ3 = 0 or 90
◦. Notice that this phase is still present in
the oscillation probabilities even when θ12 and θ13 vanish, see Eq. (10). At atmospheric
L/E, oscillations driven by an O(eV2) mass difference are averaged. We have checked
that our results apply for any value of ∆m2
≥ 0.1 eV2.
In Fig. 4 we show the sensitivity limit at 99% CL in the (θ13, θ14) plane (left) and
in the (θ24, θ34) plane (right) from a null result of the OPERA experiment, assuming
1, 2, 3, 5 and 10 times the nominal intensity of 4.5 × 1019 pot/year. The coloured
regions show the present bounds at 90% and 99% CL. We assume θ23 = 45
◦ and
δ3 = 0
◦ (top) or δ3 = 90
◦ (bottom). The sensitivity is defined as the region for which a
(poissonian) 2 d.o.f.’s χ2 is compatible with a “null result” at the 99% CL. We refer to
“null result” when θ13 and the three active-sterile mixing angles, θ14, θ24 and θ34 vanish
simultaneously. Both νµ → ντ and νµ → νe oscillations have been considered, with the
corresponding backgrounds treated properly as in Sec. 5. An overall systematic error
of 10% has been taken into account.
In the left panels of Fig. 4 we can see that OPERA can improve only a little the bound
on θ13 after 5 years of data taking working at nominal CNGS beam intensity, both
for δ3 = 0 (top panel) or δ3 = 90
◦ (bottom panel). Increasing the nominal intensity,
however, a significant improvement on the bound is achieved for any value of θ14. Notice
that the limit on θ14 is almost unaffected by the OPERA data. This is because for the
νµ → ντ and νµ → νe oscillation probabilities at atmospheric L/E, the θ14-dependence
always arises at third-order in the small parameters θ13, θ14, θ24 and θ34 (see Eqs. (9)
and (10) for the explicit expression in the adopted parametrization, Eq. (3)). On the
contrary, the θ13-, θ24- and θ34-dependences in the same oscillation probabilities are
quadratic in the small parameters. In case of vanishing active-sterile mixing angles,
0 2 4 6 8 10 12 14
= 45°, δ
= 0°
0 2 4 6 8 10 12 14
= 45°, δ
= 0°
0 2 4 6 8 10 12 14
= 45°, δ
= 90°
0 2 4 6 8 10 12 14
= 45°, δ
= 90°
× 5× 10
Fig. 4. Sensitivity limit at 99% CL in the (θ13, θ14) plane (left) and in the (θ24, θ34) plane
(right) from a null result of the OPERA experiment, assuming 1, 2, 3, 5 and 10 times the
nominal intensity of 4.5 × 1019 pot/year. The coloured regions show the present bounds at
90% and 99% CL. We assume θ23 = 45
◦ and δ3 = 0
◦ (top) or δ3 = 90
◦ (bottom).
θi4 = 0, see Ref. [60].
In the right panels of Fig. 4 the sensitivity of OPERA to θ24 and θ34 is shown. First
of all, notice that the sensitivity is strongly affected by the intensity of the beam.
No improvement on the existing bounds on these two parameters is achieved after 5
years of data taking at nominal CNGS beam intensity, for any of the considered value
of δ3. Already with a doubled flux intensity, some sensitivity to θ24, θ34 is achievable.
The sensitivity enhancement strongly depends on the value of the CP-violating phase
δ3, however. For δ3 = 0, OPERA can exclude a small part of the 99% CL allowed
region, only. On the other hand, for δ3 = 90
◦ twice the nominal CNGS flux suffices
0 2 4 6 8 10 12 14
= 45°, δ
= 0°
0 2 4 6 8 10 12 14
= 45°, δ
= 0°
0 2 4 6 8 10 12 14
= 45°, δ
= 90°
0 2 4 6 8 10 12 14
= 45°, δ
= 90°
Fig. 5. Sensitivity limit at 99% CL in the (θ13, θ14) plane (left) and in the (θ24, θ34) plane
(right) from the combined analysis of present data and a null result of the OPERA experiment,
assuming 1, 2, 3, 5 and 10 times the nominal intensity of 4.5 × 1019 pot/year. The coloured
regions show the present bounds at 90% and 99% CL. We assume θ23 = 45
◦ and δ3 = 0
(top) or δ3 = 90
◦ (bottom).
to put a bound on θ34 ≤ 25
◦ for θ24 ≥ 4
◦ at 99% CL. For maximal CP-violating
δ3, increasing further the CNGS flux can significantly constrain the (θ24, θ34) allowed
parameter space. Notice, eventually, the strong correlation between θ24 and θ34 in the
right panels of Fig. 4. This is an indication that the dominant channel that constrains
these angles is νµ → ντ . As it can be seen in Eq. (10), the two angles always appear in
combination, with an approximate exchange symmetry θ24 ↔ θ34.
The allowed regions at 99% CL in the (θ13, θ14) plane (left) and in the (θ24, θ34) plane
(right) from the combined analysis of present data and a null result of the OPERA
experiment after 5 years of data taking (assuming 1, 2, 3, 5 and 10 times the nominal
CNGS intensity of 4.5 × 1019 pot/year) are eventually shown in Fig. 5. The coloured
regions refer to the present bounds at 90% and 99% CL, for θ23 = 45
◦ and δ3 = 0
◦ (top)
or δ3 = 90
◦ (bottom). As it can be seen, the sensitivity of OPERA strongly benefits
from the complementary information on the neutrino parameters provided by other
experiments. In this case, even with the nominal beam intensity the extension of the
allowed regions is reduced by a moderate but non-negligible amount.
7 Conclusions
The results of atmospheric, solar, accelerator and reactor neutrino experiments show
that flavour mixing occurs not only in the quark sector, as it has been known for
long, but also in the leptonic sector. Experimental data well fit into a three-family
scenario. The existence of new “sterile” neutrino states with masses in the eV range is
not excluded, however, provided that their couplings with active neutrinos are small
enough.
In this paper, we have tried to test the potential of the OPERA experiment at the
CNGS beam to improve the present bounds on the parameters of the so-called four-
neutrino models. The model, in which only one sterile neutrino is added to the three
active ones responsible for solar and atmospheric oscillations, is the minimal extension
of the standard three-family oscillation scenario.
We have determined the presently allowed regions for all active-sterile mixing angles
and studied the OPERA capability to constrain them further using both the νµ → νe
and νµ → ντ channels. We have performed our analysis using the OPERA detector as
a reference. It can be extended including a detailed simulation of the ICARUS detector
at the CNGS beam.
Our conclusions are the following: if the OPERA detector is exposed to the nominal
CNGS beam intensity, a null result can improve a bit the present bound on θ13, but
not those on the active-sterile mixing angles, θ14, θ24 and θ34. If the beam intensity is
increased by a factor 2 or beyond, not only the sensitivity to θ13 increases accordingly,
but a significant sensitivity to θ24 and θ34 is achievable. The (θ24, θ34) sensitivity strongly
depends on the value of the CP-violating phase δ3, however, with stronger sensitivity for
values of δ3 approaching π/2. Only a marginal improvement is achievable on the bound
on θ14, that should be constrained by high-intensitiy νe disappearance experiments.
Notice that our results hold for any value of ∆m2
≥ 0.1 eV2, i.e. in the region of L/E
for which oscillations driven by this mass difference are effectively averaged.
Acknowledgements
We acknowledge E. Fernández-Mart́ınez, P. Hernández, J. López-Pavón, M. Sorel and
P. Strolin for useful discussions and comments. We thank T. Schwetz for pointing out
to us an error in the first version of the paper and for useful comments on it. The
work has been partially supported by the E.U. through the BENE-CARE networking
activity MRTN-CT-2004-506395. A.D. received partial support from CiCYT through
the project FPA2006-05423. M.M. received partial support from CiCYT through the
project FPA2006-01105 and the MCYT through the Ramón y Cajal program. A.D. and
M.M. acknowledge also financial support from the Comunidad Autónoma de Madrid
through the project P-ESP-00346. D.M. would like to thank CERN, where part of this
work has been accomplished.
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Introduction
Four neutrino mass schemes
Oscillation probabilities and allowed parameter space
The CNGS facility
Appearance channels at the CNGS
oscillations
e oscillations
Sensitivity to (3+1) sterile neutrinos at OPERA
Conclusions
References
|
0704.0389 | Evolution of the Carter constant for inspirals into a black hole: effect
of the black hole quadrupole | For reference, the following erratum corrects the published version of the paper. These errors have been fixed in
this arxiv-version (the article starting on page 2 has the corrected expressions).
Erratum: Evolution of the Carter constant for inspirals into a black hole: Effect of the
black hole quadrupole
[Phys. Rev. D 75, 124007 (2007)]
Éanna É. Flanagan, Tanja Hinderer
In Eqs. (3.16), (3.17), (3.18), (3.24), (3.25) and (3.26) of this paper, the variable r should be replaced everywhere
by the variable r̃, and the variable θ should be replaced everywhere by the variable θ̃. The definitions of r̃ and θ̃ are
given in Eq. (2.11). These replacements do not affect the any of the subsequent results in the paper.
Also, the right hand side of Eq. (B3) is missing a term −4SLzr̃ and Eq. (2.24) is missing a factor of dϕ/dt̃ in front
of Q.
Some terms are missing in Eqs. (3.18), (3.26) and (3.30) - (3.33). The additional terms in Eqs. (3.18) and (3.26)
15r̃7
−75K2 + 2Kr̃(51r̃E + 50) + 8r̃2(r̃E + 1)(3r̃E + 5)
15p2r̃7
25p3(3p− 4r̃) + p2r̃2
11− 51e2
+ 32pr̃3
1− e2
+ 6r̃4
1− e2
respectively. These result in additional fractional corrections to Eq. (3.30) given by
and the full expression replacing the O(Q) terms in Eq. (3.30) is then
〈K̇〉 = −
(1− e2)3/2
cos(2ι)
+O(S), O(S2)− terms.
Equations (3.31), (3.32) and (3.33) contain typos in the O(S) and O(Q) terms, the corrected expressions are given
below. We thank P. Komorowski for pointing this out. Equation (3.31) should be replaced by
〈ṗ〉 = −64
(1− e2)3/2
− S cos(ι)
96p3/2
1064 + 1516e2 + 475e4
149e2
469e2
227e4
cos(2ι)
+ e2 +
[13− cos(2ι)]
, (0.1)
Equation (3.32) should be replaced by
〈ė〉 = −304
e(1− e2)3/2
121e2
Se(1− e2)3/2 cos(ι)
5p11/2
1172 + 932e2 +
1313e4
Q(1− e2)3/2
785e2
− 219e
+ 13e6 +
2195e2
+ 251e4 +
218e6
cos(2ι)
2e(1− e2)3/2
2 + 3e2 +
[13− cos(2ι)] , (0.2)
and the corrected Eq. (3.33) is
〈ι̇〉 = S sin(ι)(1 − e
2)3/2
p11/2
1− e2
S2 sin(2ι)
240p6
8 + 3e2
8 + e2
Q cot(ι)(1 − e2)3/2
312 + 736e2 − 83e4 −
408 + 1268e2 + 599e4
cos(2ι)
. (0.3)
http://arxiv.org/abs/0704.0389v8
Evolution of the Carter constant for inspirals into a black hole: effect of the black hole
quadrupole
Éanna É. Flanagan1,2 and Tanja Hinderer1
Center for Radiophysics and Space Research, Cornell University, Ithaca, NY 14853, USA
Laboratory for Elementary Particle Physics, Cornell University, Ithaca, NY 14853, USA
(Dated: November 4, 2018)
We analyze the effect of gravitational radiation reaction on generic orbits around a body with an
axisymmetric mass quadrupole moment Q to linear order in Q, to the leading post-Newtonian order,
and to linear order in the mass ratio. This system admits three constants of the motion in absence
of radiation reaction: energy, angular momentum along the symmetry axis, and a third constant
analogous to the Carter constant. We compute instantaneous and time-averaged rates of change
of these three constants. For a point particle orbiting a black hole, Ryan [15] has computed the
leading order evolution of the orbit’s Carter constant, which is linear in the spin. Our result, when
combined with an interaction quadratic in the spin (the coupling of the black hole’s spin to its own
radiation reaction field), gives the next to leading order evolution. The effect of the quadrupole, like
that of the linear spin term, is to circularize eccentric orbits and to drive the orbital plane towards
antialignment with the symmetry axis.
In addition we consider a system of two point masses where one body has a single mass multipole
or current multipole of order l. To linear order in the mass ratio, to linear order in the multipole,
and to the leading post-Newtonian order, we show that there does not exist an analog of the Carter
constant for such a system (except for the cases of an l = 1 current moment and an l = 2 mass
moment). Thus, the existence of the Carter constant in Kerr depends on interaction effects between
the different multipoles. With mild additional assumptions, this result falsifies the conjecture that
all vacuum, axisymmetric spacetimes posess a third constant of the motion for geodesic motion.
PACS numbers: 04.25.Nx, 04.30.Db
I. INTRODUCTION AND SUMMARY
The inspiral of stellar mass compact objects with
masses µ in the range µ ∼ 1 − 100M⊙ into massive
black holes with masses M ∼ 105 − 107M⊙ is one of
the most important sources for the future space-based
gravitational wave detector LISA. Observing such events
will provide a variety of information: (i) the masses and
spins of black holes can be measured to high accuracy
(∼ 10−4); which can constrain the black hole’s growth
history [1]; (ii) the observations will give a precise test
of general relativity in the strong field regime and unam-
biguously identify whether the central object is a black
hole [2]; and (iii) the measured event rate will give in-
sight into the complex stellar dynamics in galactic nu-
clei [1]. Analogous inspirals may also be interesting for
the advanced stages of ground-based detectors: it has
been estimated that advanced LIGO could detect up to
∼ 10 − 30 inspirals per year of stellar mass compact
objects into intermediate mass black holes with masses
M ∼ 102 − 104M⊙ in globular clusters [3]. Detect-
ing these inspirals and extracting information from the
datastream will require accurate models of the gravita-
tional waveform as templates for matched filtering. For
computing templates, we therefore need a detailed un-
derstanding of the how radiation reaction influences the
evolution of bound orbits around Kerr black holes [4–7].
There are three dimensionless parameters characteriz-
ing inspirals of bodies into black holes:
• the dimensionless spin parameter a = |S|/M2 of
the black hole, where S is the spin.
• the strength of the interaction potential ǫ2 =
GM/rc2, i.e. the expansion parameter used in post-
Newtonian (PN) theory.
• the mass ratio µ/M .
For LISA data analysis we will need waveforms that are
accurate to all orders in a and ǫ2, and to leading order
in µ/M . However, it is useful to have analytic results in
the regimes a ≪ 1 and/or ǫ2 ≪ 1. Such approximate
results can be useful as a check of numerical schemes
that compute more accurate waveforms, for scoping out
LISA’s data analysis requirements [1, 6], and for assessing
the accuracy of the leading order in µ/M or adiabatic
approximation [8–10]. There is substantial literature on
such approximate analytic results, and in this paper we
will extend some of these results to higher order.
A long standing difficulty in computing the evolution
of generic orbits has been the evolution of the orbit’s
”Carter constant”, a constant of motion which governs
the orbital shape and inclination. A theoretical prescrip-
tion now exists for computing Carter constant evolution
to all orders in ǫ and a in the adiabatic limit µ ≪ M
[9, 11–13], but it has not yet been implemented numer-
ically. In this paper we focus on computing analytically
the evolution of the Carter constant in the regime a≪ 1,
ǫ ≪ 1, µ/M ≪ 1, extending earlier results by Ryan
[14, 15].
We next review existing analytical work on the effects
of multipole moments on inspiral waveforms. For non-
spinning point masses, the phase of the l = 2 piece of
the waveform is known to O(ǫ7) beyond leading order
[16], while spin corrections are not known to such high
order. To study the leading order effects of the central
body’s multipole moments on the inspiral waveform, in
the test mass limit µ ≪ M , one has to correct both the
conservative and dissipative pieces of the forces on the
bodies. For the conservative pieces, it suffices to use the
Newtonian action for a binary with an additional multi-
pole interaction potential. For the dissipative pieces, the
multipole corrections to the fluxes at infinity of the con-
served quantities can simply be added to the known PN
point mass results. The lowest order spin-orbit coupling
effects on the gravitational radiation were first derived by
Kidder [17], then extended by Ryan [14, 15], Gergely [18],
and Will [19]. Recently, the corrections of O(ǫ2) beyond
the leading order to the spin-orbit effects on the fluxes
were derived [20, 21]. Corrections to the waveform due
to the quadrupole - mass monopole interaction were first
considered by Poisson [22], who derived the effect on the
time averaged energy flux for circular equatorial orbits.
Gergely [23] extended this work to generic orbits and
computed the radiative instantaneous and time averaged
rates of change of energy E, magnitude of angular mo-
mentum |L|, and the angle κ = cos−1(S ·L) between the
spin S and orbital angular momentum L. Instead of the
Carter constant, Gergely identified the angular average
of the magnitude of the orbital angular momentum, L̄, as
a constant of motion. The fact that to post-2-Newtonian
(2PN) order there is no time averaged secular evolution of
the spin allowed Gergely to obtain expressions for L̇ and
κ̇ from the quadrupole formula for the evolution of the
total angular momentum J = L+S. In a different paper,
Gergely [18] showed that in addition to the quadrupole,
self-interaction spin effects also contribute at 2PN order,
which was seen previously in the black hole perturbation
calculations of Shibata et al. [24]. Gergely calculated
the effect of this interaction on the instantaneous and
time-averaged fluxes of E and |L| but did not derive the
evolution of the third constant of motion.
In this paper, we will re-examine the effects of the
quadrupole moment of the black hole and of the leading
order spin self interaction. For a black hole, our analysis
will thus contain all effects that are quadratic in spin to
the leading order in ǫ2 and in µ/M . Our work will extend
earlier work by
• Considering generic orbits.
• Using a natural generalization of the Carter-type
constant that can be defined for two point particles
when one of them has a quadrupole. This facilitates
applying our analysis to Kerr inspirals.
• Computing instantaneous as well as time-averaged
fluxes for all three constants of motion: energy
E, z-component of angular momentum Lz, and
Carter-type constant K. For most purposes, only
time-averaged fluxes are needed as only they are
gauge invariant and physically relevant. However,
there is one effect for which the time-averaged
fluxes are insufficient, namely transient resonances
that occur during an inspiral in Kerr in the vicin-
ity of geodesics for which the radial and azimuthal
frequencies are commensurate [10, 25]. The instan-
taneous fluxes derived in this paper will be used in
[10] for studying the effect of these resonances on
the gravitational wave phasing.
We will analyze the effect of gravitational radiation re-
action on orbits around a body with an axisymmetric
mass quadrupole moment Q to leading order in Q, to the
leading post-Newtonian order, and to leading order in the
mass ratio. With these approximations the adiabatic ap-
proximation holds: gravitational radiation reaction takes
place over a timescale much longer than the orbital pe-
riod, so the orbit looks geodesic on short timescales. We
follow Ryan’s method of computation [14]: First, we cal-
culate the orbital motion in the absence of radiation re-
action and the associated constants of motion. Next, we
use the leading order radiation reaction accelerations that
act on the particle (given by the Burke-Thorne formula
[26] augmented by the relevant spin corrections [14]) to
compute the evolution of the constants of motion. In the
adiabatic limit, the time-averaged rates of change of the
constants of motion can be used to infer the secular or-
bital evolution. Our results show that a mass quadrupole
has the same qualitative effect on the evolution as spin: it
tends to circularize eccentric orbits and drive the orbital
plane towards antialignment with the symmetry axis of
the quadrupole.
The relevance of our result to point particles inspi-
ralling into black holes is as follows. The vacuum space-
time geometry around any stationary body is completely
characterized by the body’s mass multipole moments
IL = Ia1,a2...al and current multipole moments SL =
Sa1,a2...al [27]. These moments are defined as coefficients
in a power series expansion of the metric in the body’s
local asymptotic rest frame [28]. For nearly Newtonian
sources, they are given by integrals over the source as
IL ≡ Ia1,...al =
ρx<a1 . . . xal>d
3x, (1.1)
SL ≡ Sa1,...al =
ρxpvqǫpq<a1xa2 . . . xal>d
3x.(1.2)
Here ρ is the mass density and vq is the velocity, and ”<
· · · >” means ”symmetrize and remove all traces”. For
axisymmetric situations, the tensor multipole moments
IL (SL) contain only a single independent component,
conventionally denoted by Il (Sl) [27]. For a Kerr black
hole of mass M and spin S, these moments are given by
Il + iSl =M
l+1(ia)l, (1.3)
where a is the dimensionless spin parameter defined by
a = |S|/M2. Note that Sl = 0 for even l and Il = 0 for
odd l.
Consider now inspirals into an axisymmetric body
which has some arbitrary mass and current multipoles
Il and Sl. Then we can consider effects that are linear in
Il and Sl for each l, effects that are quadratic in the mul-
tipoles proportional to IlIl′ , IlSl′ , SlSl′ , effects that are
cubic, etc. For a general body, all these effects can be sep-
arated using their scalings, but for a black hole, Il ∝ al
for even l and Sl ∝ al for odd l [see Eq.(1.3)], so the ef-
fects cannot be separated. For example, a physical effect
that scales as O(a2) could be an effect that is quadratic
in the spin or linear in the quadrupole; an analysis in
Kerr cannot distinguish these two possibilities. For this
reason, it is useful to analyze spacetimes that are more
general than Kerr, characterized by arbitrary Il and Sl,
as we do in this paper. For recent work on computing
exact metrics characterized by sets of moments Il and Sl,
see Refs. [29, 30] and references therein.
The leading order effect of the black hole’s multipoles
on the inspiral is the O(a) effect computed by Ryan [15].
This O(a) effect depends linearly on the spin S1 and is
independent of the higher multipoles Sl and Il since these
all scale as O(a2) or smaller. In this paper we compute
the O(a2) effect on the inspiral, which includes the lead-
ing order linear effect of the black hole’s quadrupole (lin-
ear in I2 ≡ Q) and the leading order spin self-interaction
(quadratic in S1).
We next discuss how these O(a2) effects scale with the
post-Newtonian expansion parameter ǫ. Consider first
the conservative orbital dynamics. Here it is easy to see
that fractional corrections that are linear in I2 scale as
O(a2ǫ4), while those quadratic in S1 scale as O(a
2ǫ6).
Thus, the two types of terms cleanly separate. We com-
pute only the leading order, O(a2ǫ4), term. For the dissi-
pative contributions to the orbital motion, however, the
scalings are different. There are corrections to the radi-
ation reaction acceleration whose fractional magnitudes
are O(a2ǫ4) from both types of effects linear in I2 and
quadratic in S1. The effects quadratic in S1 are due to
the backscattering of the radiation off the piece of space-
time curvature due to the black hole’s spin. This effect
was first pointed out by Shibata et al. [24], who com-
puted the time-averaged energy flux for circular orbits
and small inclination angles based on a PN expansion of
black hole perturbations. Later, Gergely [18] analyzed
this effect on the instantaneous and time-averaged fluxes
of energy and magnitude of orbital angular momentum
within the PN framework.
The organization of this paper is as follows. In Sec.
II, we study the conservative orbital dynamics of two
point particles when one particle is endowed with an ax-
isymmetric quadrupole, in the weak field regime, and to
leading order in the mass ratio. In Sec. III, we com-
pute the radiation reaction accelerations and the instan-
taneous and time-averaged fluxes. In order to have all
the contributions at O(a2ǫ4) for a black hole, we include
in our computations of radiation reaction acceleration
the interaction that is quadratic in the spin S1. The ap-
plication to black holes in Sec. IV briefly discusses the
qualitative predictions of our results and also compares
with previous results.
The methods used in this paper can be applied only
to the black hole spin (as analyzed by Ryan [14]) and
the black hole quadrupole (as analyzed here). We show
in Sec. V that for the higher order mass and current
multipole moments taken individually, an analog of the
Carter constant cannot be defined to the order of our
approximations. We then show that under mild assump-
tions, this non-existence result can be extended to exact
spacetimes, thus falsifying the conjecture that all vac-
uum axisymmetric spacetimes possess a third constant
of geodesic motion.
II. EFFECT OF AN AXISYMMETRIC MASS
QUADRUPOLE ON THE CONSERVATIVE
ORBITAL DYNAMICS
Consider two point particles m1 and m2 interacting in
Newtonian gravity, where m2 ≪ m1 and where the mass
m1 has a quadrupole moment Qij which is axisymmetric:
Qij =
d3xρ(r)
xixj −
r2δij
(2.1)
ninj −
. (2.2)
For a Kerr black hole of mass M and dimensionless spin
parameter a with spin axis along n, the quadrupole scalar
is Q = −M3a2.
The action describing this system, to leading order in
m2/m1, is
µv2 − µΦ(r)
, (2.3)
where v = ṙ is the velocity, the potential is
Φ(r) = −M
xixjQij , (2.4)
µ is the reduced mass and M the total mass of the bi-
nary, and we are using units with G = c = 1. We work to
linear order in Q, to linear order in m2/m1, and to lead-
ing order in M/r. In this regime, the action (2.3) also
describes the conservative effect of the black hole’s mass
quadrupole on bound test particles in Kerr, as discussed
in the introduction. We shall assume that the quadrupole
Qij is constant in time. In reality, the quadrupole will
evolve due to torques that act to change the orientation
of the central body. An estimate based on treating m1 as
a rigid body in the Newtonian field ofm2 gives the scaling
of the timescale for the quadrupole to evolve compared to
the radiation reaction time as (see Appendix I for details)
Tevol
(2.5)
Here, we have denoted the dimensionless spin and
quadrupole of the body by S̄ and Q̄ respectively, and
the last relation applies for a Kerr black hole. Since
µ/M ≪ 1, the first factor in Eq. (2.5) will be large, and
since 1/a ≥ 1 and for the relativistic regime M/r ∼ 1,
the evolution time is long compared to the radiation re-
action time. Therefore we can neglect the evolution of
the quadrupole at leading order.
This system admits three conserved quantities, the en-
µv2 + µΦ(r), (2.6)
the z-component of angular momentum
Lz = ez · (µr× v), (2.7)
and the Carter-type constant
K = µ2(r× v)2 − 2Qµ
(n · r)2
(n · v)2 −
. (2.8)
(See below for a derivation of this expression for K).
A. Conservative orbital dynamics in a
Boyer-Lindquist-like coordinate system
We next specialize to units where M = 1. We also
define the rescaled conserved quantities by Ẽ = E/µ,
L̃z = Lz/µ, K̃ = K/µ
2, and drop the tildes. These spe-
cializations and definitions have the effect of eliminating
all factors of µ and M from the analysis. In spherical
polar coordinates (r, θ, ϕ) the constants of motion E and
Lz become
(ṙ2 + r2θ̇2 + r2 sin2 θϕ̇2)−
(1− 3 cos2 θ), (2.9)
Lz = r
2 sin2 θϕ̇. (2.10)
In these coordinates, the Hamilton-Jacobi equation is not
separable, so a separation constant K cannot readily be
derived. For this reason we switch to a different coordi-
nate system (r̃, θ̃, ϕ) defined by
r cos θ = r̃ cos θ̃
r sin θ = r̃ sin θ̃
. (2.11)
We also define a new time variable t̃ by
cos(2θ̃)
dt̃. (2.12)
The action (2.3) in terms of the new variables to linear
order in Q is
r̃2 sin2 θ̃
sin2 θ̃
. (2.13)
However, a difficulty is that the action (2.13) does not
give the same dynamics as the original action (2.3). The
reason is that for solutions of the equations of motion for
the action (2.3), the variation of the action vanishes for
paths with fixed endpoints for which the time interval ∆t
is fixed. Similarly, for solutions of the equations of motion
for the action (2.13), the variation of the action vanishes
for paths with fixed endpoints for which the time interval
∆t̃ is fixed. The two sets of varied paths are not the
same, since ∆t 6= ∆t̃ in general. Therefore, solutions of
the Euler-Lagrange equations for the action (2.3) do not
correspond to solutions of the Euler-Lagrange equations
for the action (2.13). However, in the special case of zero-
energy motions, the extra terms in the variation of the
action vanish. Thus, a way around this difficulty is to
modify the original action to be
µv2 − µΦ(r) + E
. (2.14)
This action has the same extrema as the action (2.3),
and for motion with physical energy E, the energy com-
puted with this action is zero. Transforming to the new
variables yields, to linear order in Q:
r̃2 sin2 θ̃
sin2 θ̃
+ E − QE
cos(2θ̃)
. (2.15)
The zero-energy motions for this action coincide with the
zero energy motions for the action (2.14). We use this
action (2.15) as the foundation for the remainder of our
analysis in this section.
The z-component of angular momentum in terms of
the new variables (r̃, θ̃, ϕ, t̃) is
Lz = r̃
2 sin2 θ̃
sin2 θ̃
. (2.16)
We now transform to the Hamiltonian:
p2r̃ −
− E − Q
sin2 θ̃
+QE cos(2θ̃)
(2.17)
and solve the Hamiltonian Jacobi equation. Denoting
the separation constant by K we obtain the following
two equations for the r̃ and θ̃ motions:
= 2E +
, (2.18)
= K − L
sin2 θ̃
−QE cos(2θ̃). (2.19)
Note that the equations of motion (2.18) and (2.19) have
the same structure as the equations of motion for Kerr
geodesic motion. Using Eqs. (2.18), (2.19) and (2.16)
together with the inverse of the transformation (2.11)
to linear order in Q, we obtain the expression for K in
spherical polar coordinates:
K = r4(θ̇2 + sin2 θϕ̇2) +Q(ṙ cos θ − rθ̇ sin θ)2 + Q
(ṙ2 + r2 θ̇2 + r2 sin2 θϕ̇2)− 2Q
cos2 θ. (2.20)
This is equivalent to the formula (2.8) quoted earlier.
B. Effects linear in spin on the conservative orbital
dynamics
To include the linear in spin effects, we repeat Ryan’s
analysis [14, 15] (he only gives the final, time averaged
fluxes; we will also give the instantaneous fluxes). We
can simply add these linear in spin terms to our results
because any terms of order O(SQ) will be higher than
the order a2 to which we are working. The correction to
the action (2.3) due to spin-orbit coupling is
Sspin−orbit =
−2µSn
iǫijkxj ẋk
. (2.21)
We will restrict our analysis to the case when the unit
vectors ni corresponding to the axisymmetric quadrupole
Qij and to the spin Si coincide, as they do in Kerr.
Including the spin-orbit term in the action (2.3) results
in the following modified expressions for Lz and K:
Lz = n · (µr× v)−
[r2 − (n · r)2], (2.22)
K = (r× v)2 − 4S
n · (r× v)− 2Q
(n · r)2
(n · v)2 − 1
. (2.23)
In terms of the Boyer-Lindquist like coordinates, the con-
served quantities with the linear in spin terms included
Lz = r̃
2 sin2 θ̃
sin2 θ̃ −Q sin4 θ̃
(2.24)
K = r4(θ̇2 + sin2 θϕ̇2)− 4Sr sin2 θϕ̇
cos2 θ +Q(ṙ cos θ − rθ̇ sin θ)2 + QM
(ṙ2 + r2θ̇2 + r2 sin2 θϕ̇2). (2.25)
The equations of motion are
= 2E+
− 4SLz
, (2.26)
= K −
sin2 θ̃
−QE cos(2θ̃). (2.27)
III. EFFECTS LINEAR IN QUADRUPOLE AND
QUADRATIC IN SPIN ON THE EVOLUTION OF
THE CONSTANTS OF MOTION
A. Evaluation of the radiation reaction force
The relative acceleration of the two bodies can be writ-
ten as
a = −∇Φ(r) + arr, (3.1)
where arr is the radiation-reaction acceleration. Combin-
ing this with Eqs. (2.6), (2.22) and (2.23) for E, Lz and
K gives the following formulae for the time derivatives of
the conserved quantities:
Ė = v · arr, (3.2)
L̇z = n · (r× arr), (3.3)
K̇ = 2(r× v) · (r× arr)−
n · (r× arr)
+2Q(n · v) (n · arr)−Qv · arr. (3.4)
The standard expression for the leading order radiation
reaction acceleration acting on one of the bodies is [31]:
ajrr = −
jk xk +
ǫjpqS
pk xkxq +
ǫjpqS
pk xkvq
ǫpq[jS
xqvk. (3.5)
Here the superscripts in parentheses indicate the number
of time derivatives and square brackets on the indices
denote antisymmetrization.
The multipole moments Ijk(t) and Sjk(t) in Eq. (3.5)
are the total multipole moments of the spacetime, i.e.
approximately those of the black hole plus those due to
the orbital motion. The expression (3.5) is formulated
in asymptotically Cartesian mass centered (ACMC) co-
ordinates of the system, which are displaced from the
coordinates used in Sec. II by an amount [28]
δr(t) = − µ
r(t). (3.6)
This displacement contributes to the radiation reaction
acceleration in the following ways:
1. The black hole multipole moments Il and Sl, which
are time-independent in the coordinates used in
Sec. II, will be displaced by δr and thus will con-
tribute to the (l + 1)th ACMC radiative multipole
[28].
2. The constants of motion are defined in terms of the
black hole centered coordinates used in Sec. II, so
the acceleration arr we need in Eqs. (3.2) – (3.4)
is the relative acceleration. This requires calculat-
ing the acceleration of both the black hole and the
point mass in the ACMC coordinates using (3.5),
and then subtracting to find arr = a
rr − aMrr [14].
To leading order in µ, the only effect of the acceler-
ation of the black hole is via a backreaction of the
radiation field: the lth black hole moments couple
to the (l+1)th radiative moments, thus producing
an additional contribution to the acceleration.
For our calculations at O(S1ǫ
3), O(I2ǫ
4), O(S21ǫ
4), we
can make the following simplifications:
• quadrupole corrections: The fractional corrections
linear in I2 = Q that scale as O(a
2ǫ4) require only
the effect of I2 on the conservative orbital dynamics
as computed in Sec. IIA and the Burke-Thorne for-
mula for the radiation reaction acceleration [given
by the first term in Eq. (3.5)].
• spin-spin corrections: As discussed in the intro-
duction, the fractional corrections quadratic in S1
to the conservative dynamics scale as O(a2ǫ6) and
are subleading order effects which we neglect. At
O(a2ǫ4), the only effect quadratic in S1 is the
backscattering of the radiation off the spacetime
curvature due to the spin. As discussed in item 1.
above, the black hole’s current dipole Si = S1δi3
(taking the z-axis to be the symmetry axis) will
contribute to the radiative current quadrupole an
amount
ij = −
S1xiδj3. (3.7)
The black hole’s current dipole Si will couple to
the gravitomagnetic radiation field due to Sij as
discussed in item 2. above, and contribute to the
relative acceleration as [14]:
aj spinrr =
S1δi3S
ij . (3.8)
For our purposes of computing terms quadratic in
the spin, we substitute S
ij for Sij in Eq. (3.8).
Evaluating these quadratic in spin terms requires
only the Newtonian conservative dynamics, i.e. the
results of Sec. II and Eqs. (3.2) – (3.4) with the
quadrupole set to zero.
• linear in spin corrections: Contributions to these
effects are from Eq. (3.5) with the current
quadrupole replaced by just the spin contribution
(3.7), and from Eq. (3.8) evaluated using only the
orbital current quadrupole.
With these simplifications, we replace the expression
(3.5) for the radiation reaction acceleration with
ajrr = −
jk xk +
ǫjpqS
(6) spin
pk xkxq
ǫjpqS
(5) spin
pk xkvq +
ǫpq[jS
(5) spin
S1δi3
(5) orbit
ij + S
(5) spin
. (3.9)
To justify these approximations, consider the scaling of
the contribution of black hole’s acceleration to the orbital
dynamics. The mass and current multipoles of the black
hole contribute terms to the Hamiltonian that scale with
∆H ∼ Slǫ2l+3 & Ilǫ2l+2. (3.10)
Since the Newtonian energy scales as ǫ2, the fractional
correction to the orbital dynamics scale as
∆H/E ∼ Slǫ2l+1 & Ilǫ2l. (3.11)
To O(ǫ4), the only radiative multipole moments that con-
tribute to the acceleration (3.5) are the mass quadrupole
I2, the mass octupole I3, and the current quadrupole S2
(cf. [17]). Since we are focusing only on the leading or-
der terms quadratic in spin (these can simply be added
to the known 2PN point particle and 1.5PN linear in spin
results), the only terms in Eq. (3.5) relevant for our pur-
poses are those given in Eq. (3.9). The results from a
computation of the fully relativistic metric perturbation
for black hole inspirals [24] show that quadratic in spin
corrections to the l = 2 piece compared to the flat space
Burke-Thorne formula first appear at O(a2ǫ4), which is
consistent with the above arguments.
B. Instantaneous fluxes
We evaluate the radiation reaction force as follows.
The total mass and current quadrupole moment of the
system are
QTij = Qij + µxixj , (3.12)
STij = S
ij + xiǫjkmxkẋm, (3.13)
where from Eq. (2.11)
r̃ sin θ̃
cosϕ, r̃ sin θ̃
sinϕ,
r̃ cos θ̃
. (3.14)
Only the second term in Eq. (3.12) contributes to the
time derivative of the quadrupole. We differentiate five
times by using
cos(2θ̃)
, (3.15)
to the order we are working as discussed above. Af-
ter each differentiation, we eliminate any occurrences of
dϕ/dt̃ using Eq. (2.24), and we eliminate any occurrences
of the second order time derivatives d2r̃/dt̃2 and d2θ̃/dt̃2
in favor of first order time derivatives using (the time
derivatives of) Eqs. (2.26) and (2.27). For computing the
terms linear and quadratic in S1, we set the quadrupole
Q to zero in all the formulae. We insert the resulting ex-
pression into the formula (3.9) for the self-acceleration,
and then into Eqs. (3.2) – (3.4). We eliminate (dr̃/dt̃)2,
(dθ̃/dt̃)2, and (dϕ/dt̃) in favor of E, Lz, and K using
Eqs. (2.24) – (2.27). In the final expressions for the in-
stantaneous fluxes, we keep only terms that are of O(S),
O(Q) and O(S2) and obtain the following results:
15r̃4
− 40K
272KE
196K2 +
r̃2 − 3668
Kr̃ − 352KEr̃2 + 1024
Er̃3 +
E2r̃4
−49K2 − 169KL2z + r̃
+ 2r̃2
+ 47KE +
− 152
r̃3E − 16r̃4E2
−562K2 +
Kr̃ −
r̃2 +
KEr̃2 −
r̃3E − 160r̃4E2
cos(2θ̃)
sin(2θ̃)
439K − 926
r̃ − 1528
θ ˙̃r
−K2 + 22
Kr̃ − 28
r̃2 +
KEr̃2 − 236
r̃3E − 32
r̃4E2
cos(2θ̃)− r̃3 sin(2θ̃)
θ ˙̃r
−49K2 + 6KL2z + 2r̃
63K − 16
L2z −
+ r̃2
112KE − 48
− 1652
r̃3E − 224
r̃4E2
, (3.16)
L̇z =
144LzE
− 24KLz
−50K2 + 240KL2z +
Kr̃ − 7376
L2z r̃ +
r̃2 + 56KEr̃2 − 1824
EL2z r̃
Er̃3 +
E2r̃4
50K2 − 62
Kr̃ − 316
r̃2 − 56KEr̃2 − 624
Er̃3 − 128
E2r̃4
cos(2θ̃)
−104K + 64r̃ + 64Er̃2
sin(2θ̃) ˙̃r
660Er̃2 + 753r̃− 360L2z − 435K +
1601r̃+ 1512r̃2E − 1185K
cos 2θ̃
174QLz
sin(2θ̃) ˙̃r
2S2Lz
Er̃2 + 16r̃ − 9K
, (3.17)
20r̃ + 18r̃2E − 15K
280K2 − 14008
Kr̃ +
r̃2 +
Er̃3 − 2528
KEr̃2 +
E2r̃4
−45K2 + r̃L2z(83 + 80r̃E)− 115KL2z + 14Kr̃(6 + 5r̃E)
15r̃7
cos(2θ̃)
−2175K2 + 2975Kr̃+ 80r̃2 + 3012KEr̃2 − 112Er̃3 − 168E2r̃4
15r̃4
3075K − 20r̃ − 192Er̃2
sin(2θ̃)
θ ˙̃r
7K − 2L2z
−3K + 16
+K cos(2θ̃)
3K − 16
r̃ − 24
sin(2θ̃)
−4K +
θ ˙̃r. (3.18)
C. Alternative set of constants of the motion
A body in a generic bound orbit in Kerr traces an
open ellipse precessing about the hole’s spin axis. For
stable orbits the motion is confined to a toroidal region
whose shape is determined by E, Lz, K. The motion
can equivalently be characterized by the set of constants
inclination angle ι, eccentricity e, and semi-latus rectum
p defined by Hughes [32]. The constants ι, p and e are
defined by cos ι = Lz/
K, and by r̃± = p/(1± e), where
r̃± are the turning points of the radial motion, and r̃
is the Boyer-Lindquist radial coordinate. This param-
eterization has a simple physical interpretation: in the
Newtonian limit of large p, the orbit of the particle is an
ellipse of eccentricity e and semilatus rectum p on a plane
whose inclination angle to the hole’s equatorial plane is
ι. In the relativistic regime p ∼M , this interpretation of
the constants e, p, and ι is no longer valid because the
orbit is not an ellipse and ι is not the angle at which the
object crosses the equatorial plane (see Ryan [14] for a
discussion).
We adopt here analogous definitions of constants of
motion ι, e and p, namely
cos(ι) = Lz/
K, (3.19)
= r̃±. (3.20)
Here K is the conserved quantity (2.23) or (2.25), and r̃±
are the turning points of the radial motion using the r̃
coordinate defined by Eq. (2.11), given by the vanishing
of the right-hand side of Eq. (2.26).
We now rewrite our results in terms of the new con-
stants of the motion e, p and ι. We can use Eq. (2.26)
together with the equations (3.19) and (3.20) to write E,
Lz and K as functions of p, e and ι. To leading order in
Q and S we obtain
K = p
1− 2S cos ι
3 + e2
1 + e2
) 2Q cos2 ι
3 + e2
, (3.21)
E = − (1− e
2S cos ι
1− e2
1− e2
cos2 ι− 1
, (3.22)
p cos ι
1− S cos ι
(3 + e2)−
1 + e2
) Q cos2 ι
3 + e2
. (3.23)
As discussed in the introduction, the effects quadratic in
S on the conservative dynamics scale as O(a2ǫ6) and thus
are not included in this analysis to O(a2ǫ4).
Inserting these relations into the expressions (3.16)–
(3.18) gives, dropping terms of O(QS), O(Q2) and
O(QS2):
Ė = −
15p2r̃7
75p4 − 100p3r̃ + p2r̃2
11− 51e2
+ 32pr̃3
1− e2
)− 6r̃4
1− e2
4S cos ι
15p7/2r̃9
735p6 − 2751p5r̃ + 10p4r̃2(365− 6e2)− 128pr̃5(1− e2)2 − 48r̃6(e2 − 1)3
64S cos ι
15p3/2r̃6
5p(−23 + 3e2)− 3r̃(−9 + e2 + 8e4)
15p4r̃9
4005p6 − 6499p5r̃ + 2p4r̃2
1577− 1977e2
− 24r̃6
1− e2
)3 − 32p3r̃3
8− 33e2
+ 64pr̃5
1− 2e2 + e4
15p4r̃9
24p2r̃4
5− 27e2 + 22e4
− pr̃3 sin(2θ̃)
6585p2 − 4630pr̃+ 2292r̃2(1 − e2)
θ ˙̃r
15p4r̃9
2p2 cos(2θ̃)
4215p4 − 7495p3r̃ + 4p2r̃2(1151− 951e2)− 1012pr̃3(1− e2) + 300r̃4(1− 2e2 + e4)
15p4r̃9
cos(2ι)
2535p6 − 3307p5r̃ + 12p4r̃2(37− 237e2)− 48r̃6(1− e2)3 + 800p3r̃3(1 + e2) + 128pr̃5(1− 2e2 + e4)
15p2r̃5
cos(2ι)
1 + 2e2 − 3e4
15p2r̃9
84r̃4(1− e2)2(1 + e2)2 + 345p4 − 905p3r̃ − 413pr̃3(1− e2) + 2p2r̃2(446− 201e2)
15p2r̃9
cos(2θ̃)
15p4 − 110p3r̃ + 4p2r̃2(47− 12e2)− 118pr̃3(1− e2) + 24r̃4(1− e2)2(1 + e2)2
15r̃9
cos(2ι)
45p2 − 80pr̃ + 36r̃2(1− e2)
15pr̃6
sin(2θ̃) ˙̃r
15p2 + 10pr̃ − 12r̃2(1− e2)
, (3.24)
L̇z = −
8 cos ι
15p2 − 20pr̃ + 9r̃2(1 − e2)
15p2r̃7
525p4 − 1751p3r̃ + 34p2r̃2(61− 6e2) + 12pr̃3(−69 + 29e2) + 6r̃4(17 + 2e2 − 19e4)
15p2r̃7
375p4 − 93p3r̃ + 468pr̃3(1− e2)− 10p2r̃2(58 + 21e2)− 48r̃4(1 − 2e2 + e4)
cos(2θ̃)
15p2r̃7
450p4 − 922p3r̃ − 60pr̃3(3 + e2)− 9p2r̃2(−83 + 23e2) + 27r̃4(1 + 2e2 − 3e4)
cos(2ι)
13p2 − 8pr̃ + 4r̃2(1− e2)
sin(2θ̃) ˙̃r
− Q cos ι
5p5/2r̃7
615p4 − 753p3r̃ + 15p2r̃2
19− 31e2
+ 20pr̃3
1 + 3e2
+ 9r̃4
1− 6e2 + 5e4
− Q cos ι
5p1/2r̃7
cos(2θ̃)
1185p2 − 1601pr̃+ 756r̃2(1− e2)
− 2Q cos ι
5p5/2r̃7
2 cos(2ι)
45p4 − 18r̃4e2(1− e2)− 45p2r̃2(1 + e2) + 20pr̃3(1 + e2)
− 435p3r̃3 sin(2θ̃) ˙̃θ ˙̃r
2 cos ι
p1/2r̃7
9p2 − 16pr̃ + 36
r̃2(1− e2)
, (3.25)
20pr̃ − 15p2 − 9r̃2(1− e2)
8S cos ι
15p3/2r̃7
525p4 − 1751p3r̃ + 2p2r̃2(1172− 57e2) + 12pr̃3(−99 + 19e2)− 24r̃4(−11 + 4e2 + 7e4)
5p2r̃7
−615p4 + 753p3r̃ + 30p2r̃2(17e2 − 9) + 72r̃4e2(1− e2)− 40pr̃3(1 + 3e2)
5p2r̃7
cos(2ι)
−345p4 + 249p3r̃ − 160pr̃3(1 + e2) + 120p2r̃2(1 + 3e2) + 36r̃4(1 + 2e2 − 3e4)
15p2r̃7
2 cos(2θ̃)
2175p4 − 2975p3r̃ − 56pr̃3(1 − e2) + 2p2r̃2(713− 753e2) + 42r̃4(1 − 2e2 + e4)
15pr̃4
sin(2θ̃)
3075p2 − 20pr̃ + 96r̃2(1 − e2)
−9p2 + 16pr̃ − 36
r̃2(1− e2)
cos(2θ̃) + cos(2ι)
3p2 − 16
pr̃ +
r̃2(1 − e2)
sin(2θ̃) ˙̃r
−2p2 + 7
pr̃ − 4
r̃2(1− e2)
. (3.26)
D. Time averaged fluxes
In this section we will compute the infinite time-
averages 〈Ė〉, 〈L̇z〉 and 〈K̇〉 of the fluxes. These averages
are defined by
〈Ė〉 ≡ lim
∫ T/2
Ė(t)dt. (3.27)
These time-averaged fluxes are sufficient to evolve or-
bits in the adiabatic regime (except for the effect of res-
onances) [12, 25]. In Appendix II, we present two dif-
ferent ways of computing the time averages. The first
approach is based on decoupling the r̃ and θ̃ motion us-
ing the analog of the Mino time parameter for geodesic
motion in Kerr [12]. The second approach uses the ex-
plicit Newtonian parameterization of the orbital motion.
Both averaging methods give the following results:
〈Ė〉 = −32
(1− e2)3/2
e4 − S
cos(ι)
cos(2ι)
cos(2ι)
,(3.28)
〈L̇z〉 = −
(1 − e2)3/2
cos ι
e2 − S
2p3/2 cos ι
+ 7e2 +
cos(2ι)
45 + 148e2 +
cos(2ι)
1 + 3e2 +
, (3.29)
〈K̇〉 = −64
(1 − e2)3/2
e2 − S
2p3/2
+ 37e2 +
cos(ι)
cos(2ι)
cos(2ι)
. (3.30)
Using Eqs. (3.21) and (3.23), we obtain from (3.28) – (3.30) the following time averaged rates of change of the
orbital elements e, p, ι:
〈ṗ〉 = −64
(1− e2)3/2
− S cos(ι)
96p3/2
1064 + 1516e2 + 475e4
149e2
469e2
227e4
cos(2ι)
+ e2 +
[13− cos(2ι)]
, (3.31)
〈ė〉 = −
e(1− e2)3/2
121e2
Se(1− e2)3/2 cos(ι)
5p11/2
1172 + 932e2 +
1313e4
Q(1− e2)3/2
785e2
219e4
+ 13e6 +
2195e2
+ 251e4 +
218e6
cos(2ι)
S2e(1− e2)3/2
2 + 3e2 +
[13− cos(2ι)] , (3.32)
〈ι̇〉 = S sin(ι)(1 − e
2)3/2
p11/2
1− e2
S2 sin(2ι)
240p6
8 + 3e2
8 + e2
Q cot(ι)(1 − e2)3/2
312 + 736e2 − 83e4 −
408 + 1268e2 + 599e4
cos(2ι)
. (3.33)
IV. APPLICATION TO BLACK HOLES
A. Qualitative discussion of results
The above results for the fluxes, Eqs. (3.31), (3.32)
and (3.33) show that the correction terms at O(a2ǫ4)
due to the quadrupole have the same type of effect on the
evolution as the linear spin correction computed by Ryan:
they tend to circularize eccentric orbits and change the
angle ι such as to become antialigned with the symmetry
axis of the quadrupole.
The effects of the terms quadratic in spin are quali-
tatively different. In the expression (3.28) for 〈Ė〉, the
coefficient of cos(2ι) due to the spin self-interaction has
the same sign as the quadrupole term, while the terms
not involving ι have the opposite sign. The terms in-
volving cos(2ι) in Eq. (3.30) for 〈K̇〉 of O(Q) and O(S2)
terms have the same sign, while the terms not involving
ι have the opposite sign. The fractional spin-spin cor-
rection to 〈L̇z〉, Eq. (3.29), has no ι-dependence, and in
expression (3.33) for 〈ι̇〉, the dependence on ι of the two
effects O(Q) and O(S2) is different, too. This is not sur-
prising as the O(Q) effects included here are corrections
to the conservative orbital dynamics, while the effects of
O(S2) that we included are due to radiation reaction.
B. Comparison with previous results
The terms linear in the spin in our results for the time
averaged fluxes, Eqs. (3.28) – (3.33), agree with those
computed by Ryan, Eqs. (14a) – (15c) of [15], and with
those given in Eqs. (2.5) – (2.7) of Ref. [33], when we use
the transformations to the variables used by Ryan given
in Eqs. (2.3) – (2.4) in [33].
Equation (3.28) for the time averaged energy flux
agrees with Eq. (3.10) of Gergely [23] and Eq. (4.15)
of [18] when we use the following transformations:
K = L̄2
Ā2 sin2 κ cos δ − (1− Ā2) cos2 κ
= L̄2
E cos2 κ
(1 + 2L̄2) sin2 κ cos δ
, (4.1)
cos ι = cosκ
E cos2 κ
(1 + 2L̄2) sin2 κ cos δ
, (4.2)
(δ + κ), (4.3)
ξ0 = (ψ0 − ψi) +
, (4.4)
where Ā, L̄, κ, δ, ψ0 and ψi are the quantities used by
Gergely. The first relation here is obtained from the turn-
ing points of the radial motion as follows. We compute
r̃± in terms of E and K and map these expressions back
to r using Eqs. (2.11). The result can then be com-
pared with the turning points in Gergely’s variables, Eq.
(2.19) of [23], using the fact that E is the same in both
cases. Instead of the evolution of the constants of motion
K and Lz, Gergely computes the rates of change of the
magnitude L of the orbital angular momentum and of the
angle κ defined by cosκ = (L · S)/L. Using the trans-
formations (4.1) – (4.4) and the definition of κ we verify
that our Eq. (3.29) agrees with the 〈L̇z〉 computed using
Gergely’s Eqs. (3.23) and (3.35) in [23] and Eq. (4.30)
of [18].
In the limit of the circular equatorial orbits analyzed
by Poisson [22], our Eq. (3.28) agrees with Poisson’s Eq.
(22) when we use the transformations and specializations:
, (4.5)
ι = 0, (4.6)
e2 = 0, (4.7)
cosαA = 1, (4.8)
where v and αA are the variables used by Poisson and the
relation (4.5) is obtained by comparing the expressions
for the constants of motion in the two sets of variables.
The main improvement of our analysis over Gergely’s
is that we express the results in terms of the Carter-type
constant K, which facilitates comparing our results with
other analyses of black hole inspirals. Our computations
also include the spin curvature scattering effects for all
three constants of motion; Gergely [18] only considers
these effects for two of them: the energy and magnitude
of angular momentum, not for the third conserved quan-
tity.
When we expand Eq. (3.28) for small inclination an-
gles and specialize to circular orbits, then after converting
p to the parameter v using Eq. (4.5), we obtain
〈Ė〉 = − 32
11Q− S
= − 32
33− 527
. (4.9)
This result agrees with the terms atO(a2v4) of Eq. (3.13)
of Shibata et al. [24], whose calculations were based on
the fully relativistic expressions. This agreement is a
check that we have taken into account all the contribu-
tions at O(a2ǫ4). The analysis in Ref. [24] could not dis-
tinguish between effects due to the quadrupole and those
due curvature scattering, but we can see from Eq. (4.9)
that those two interactions have the opposite dependence
on ι. Comparing (4.9) with Eq. (3.7) of [24] (which gives
the fluxes into the different modes (l = 2,m, n), where m
and n are the multiples of the ϕ and θ frequencies), we see
that the terms in the (2,±2, 0) and the (2,±1,±1) modes
are entirely due to the quadrupole, while the spin-spin in-
teraction effects are fully contained in the (2,±1, 0) and
(2, 0,±1) modes.
V. NON-EXISTENCE OF A CARTER-TYPE
CONSTANT FOR HIGHER MULTIPOLES
In this section, we show that for a single axisymmetric
multipole interaction, it is not possible to find an ana-
log of the Carter constant (a conserved quantity which
does not correspond to a symmetry of the Lagrangian),
except for the cases of spin (treated by Ryan [15]) and
mass quadrupole moment (treated in this paper). Our
proof is valid only in the approximations in which we
work – expanding to linear order in the mass ratio, to
the leading post-Newtonian order, and to linear order in
the multipole. However we will show below that with
very mild additional smoothness assumptions, our non-
existence result extends to exact geodesic motion in exact
vacuum spacetimes.
We start in Sec. VA by showing that there is no co-
ordinate system in which the Hamilton-Jacobi equation
is separable. Now separability of the Hamilton-Jacobi
equation is a sufficient but not a necessary condition for
the existence of a additional conserved quantity. Hence,
this result does not yield information about the existence
or non-existence of an additional constant. Nevertheless
we find it to be a suggestive result. Our actual derivation
of the non-existence is based on Poisson bracket compu-
tations, and is given in Sec. VB.
A. Separability analysis
Consider a binary of two point masses m1 and m2,
where the mass m1 is endowed with a single axisymmet-
ric current multipole moment Sl or axisymmetric mass
multipole moment Il. In this section, we show that the
Hamilton-Jacobi equation for this motion, to linear order
in the multipoles, to linear order in the mass ratio and to
the leading post-Newtonian order, is separable only for
the cases S1 and I2.
We choose the symmetry axis to be the z-axis and write
the action for a general multipole as
ṙ2 + r2θ̇2 + r2 sin2 θϕ̇2
+ f(r, θ) + g(r, θ)ϕ̇+ E] . (5.1)
For mass moments, g(r, θ) = 0, while for current mo-
ments f(r, θ) = 0. For an axisymmetric multipole of
order l, the functions f and g will be of the form
f(r, θ) =
clIlPl(cos θ)
, g(r, θ) =
dlSl sin θ∂θPl(cos θ)
(5.2)
where Pl(cos θ) are the Legendre polynomials and cl and
dl are constants. We will work to linear order in f and g.
In Eq. (5.1), we have added the energy term needed when
doing a change of time variables, cf. the discussion before
Eq. (2.14) in section III. Since ϕ is a cyclic coordinate,
pϕ = Lz is a constant of motion and the system has
effectively only two degrees of freedom. Note that in the
case of a current moment, there will be correction term
in Lz:
Lz = r
2 sin2 θϕ̇+ g(r, θ). (5.3)
Next, we switch to a different coordinate system
(r̃, θ̃, ϕ) defined by
r = r̃ + α(r̃, θ̃, Lz), (5.4)
θ = θ̃ + β(r̃, θ̃, Lz), (5.5)
where the functions α and β are yet undetermined. We
also define a new time variable t̃ by
1 + γ(r̃, θ̃, Lz)
dt̃. (5.6)
Since we work to linear order in f and g, we can work
to linear order in α, β, and γ. We then compute the
action in the new coordinates and drop the tildes. The
Hamiltonian is given by
p2r(1 + γ − 2α,r) +
(1− 2α
− 2β,θ + γ)
(−α,θ − r2β,r)− E(1 + γ)
2r2 sin2 θ
(1 + γ − 2α
− 2β cot θ)
(1− α
+ γ)− f − gLz
r2 sin2 θ
(5.7)
and the corresponding Hamilton-Jacobi equation is
Ĉ1 +
+ 2V̂ , (5.8)
where we have denoted
Ĉ1 = J(r, θ) [1 + γ − 2α,r] = 1 + γ − 2α,r + j, (5.9)
Ĉ2 = J(r, θ)
1− 2α
− 2β,θ + γ
= 1− 2α
− 2β,θ + γ + j, (5.10)
Ĉ3 = J(r, θ)
−α,θ − r2β,r
= −α,θ − r2β,r, (5.11)
V̂ = J(r, θ)
2r2 sin2 θ
(1 + γ − 2α
− 2β cot θ)
+ γ)− E(1 + γ)
− f − gLz
r2 sin2 θ
2r2 sin2 θ
(1 + γ − 2α
− 2β cot θ + j)
−E(1 + γ + j)− 1
(1− α
+ γ + j)
−f − gLz
r2 sin2 θ
. (5.12)
The unperturbed problem is separable, so make the
perturbed problem separable, we have multiplied the
Hamilton-Jacobi equation by an arbitrary function
J(r, θ), which can be expanded as J(r, θ) = 1 + j(r, θ),
where j(r, θ) is a small perturbation.
To find a solution of the form W =Wr(r)+Wθ(θ), we
first specialize to the case where Ĉ3 = 0:
− Ĉ3 = β,rr2 + α,θ = 0. (5.13)
We differentiate Eq. (5.8) with respect to θ, using Eq.
(5.8) to write (dWr/dr)
2 in terms of (dWθ/dθ)
2 and then
differentiate the result with respect to r to obtain
∂θĈ2
∂θĈ1
2V̂ ∂θĈ1
Ĉ1Ĉ2
. (5.14)
Expanding Eq. (5.14) to linear order in the small quan-
tities then yields the two conditions for the kinetic and
the potential part of the Hamiltonian to be separable:
0 = ∂r∂θ
2α,r −
− 2β,θ
, (5.15)
sin2 θ
2β,r cot
2 θ − 3β,rθ cot θ + β,r csc2 θ
sin2 θ
+ α,rθ
−∂r∂θ
Pl(cos θ) +
dlSlLz
rl sin θ
∂θPl(cos θ)
2α,rθ −
+ 2Er2α,rθ
, (5.16)
where we have used Eq. (5.2) for f and g. Therefore, the
following conditions must be satisfied:
M4(θ)−N(r) =
+ β,θ − 2α,r, (5.17)
M1(θ) = 2β cot
2 θ + β csc2 θ + β,θθ
−3β,θ cot θ, (5.18)
M2(θ) = r
2∂r(r
2β,r), (5.19)
M3(θ) = 2rα,rθ − α,θ +
∂θPl(cos θ)
−SlLz
∂θ(csc θ ∂θPl(cos θ)). (5.20)
Here, the functions M and N are arbitrary integration
constants.
Solving the condition for the kinetic term to be sep-
arable, Eq. (5.17), together with Eq. (5.13) gives the
general solution that goes to zero at large r as
cos(nθ + ν), (5.21)
β = − A
sin(nθ + ν), (5.22)
where A and ν are arbitrary and n is an integer. These
functions must satisfy the conditions (5.18) – (5.20) in
order for the potential term to be separable as well. To
see when this will be the case, we start by considering Eq.
(5.20). Substituting the general ansatz α = a1(r)a2(θ)
shows that a′2 = P
l or a
2 = (cscθ P
′ depending on
whether a mass or a current multipole is present. The
function a1(r) is then determined from
0 = 2ra′1 − a1 +
clIl/r
(l−1)
dlSlLz/r
(5.23)
Hence,
[clIl/(2l)] r
(1−l)
[dlSlLz/(2l+ 1)] r
(5.24)
so that we obtain for mass moments
Pl(cos θ)
, β =
P ′l (cos θ)
(5.25)
and for current moments
dlSlLz
2l+ 1
csc θP ′l (cos θ)
, (5.26)
dlSlLz
(2l+ 1)(l + 1)
(csc θ P ′l (cos θ))
, (5.27)
where we have used the condition (5.13) to solve for β.
Substituting this in Eq. (5.19) determines that l = 2
for mass moments and l+1 = 2 for current moments. For
an l = 2 mass moment, conditions (5.17) and (5.18) are
satisfied as well, with n = 2 and ν = 0. For the case of an
l = 1 current moment, the extra term inH is independent
of θ anyway. But for any other multipole interaction,
the Hamilton-Jacobi equation will not be separable. For
example, for the current octupole Sijk, the last term in
Eq. (5.7) is proportional to S3Lz(5 cos
2 θ − 1)/r5 and
is therefore not separable. From Eq. (5.2) one can see
that, for a general multipole, the functions f or g contain
different powers of cos θ appearing with the same power
of r since the Legendre polynomials can be expanded as
[34]:
Pl(cos θ) =
(−1)n(2l− 2n)!
2ln!(l − n)!(l − 2n)!
(cos θ)l−2n, (5.28)
where N = l/2 for even l and N = (l + 1)/2 for odd l.
It will not be possible to cancel all of these terms with
(5.21) – (5.22) for l > 2.
The case when Ĉ3 is non-vanishing will only be sepa-
rable if all the coefficients are functions of r or of θ only,
and if in addition, the potential also depends only on r or
on θ. Achieving this for our problem will not be possible
because the potential cannot be transformed to the form
required for separability.
B. Derivation of non-existence of additional
constants of the motion
In this subsection, we show using Poisson brackets that
for a single axisymmetric multipole interaction, to linear
order in the multipole and the mass ratio, a first integral
analogous to the Carter constant does not exist, except
for the cases of mass quadrupole and spin.
Suppose that such a constant does exist. We write the
Hamiltonian corresponding to the action (5.1) as H =
H0 + δH and the Carter-type constant as K = K0 +
δK(pr, pθ, Lz, r, θ), where
2r2 sin2 θ
, (5.29)
δH = − clIl
Pl(cos θ)−
dlSlLz
rl+2 sin θ
∂θPl(cos θ),(5.30)
K0 = p
sin2 θ
. (5.31)
Computing the Poisson bracket gives, to linear order in
the perturbations
0 = {H0, δK}+ {δH,K0} (5.32a)
δK + {δH,K0}, (5.32b)
where we have used that {H0,K0} = 0 and the fact that
{H0, δK} = d(δK)/dt. Here, d/dt denotes the total time
derivative along an orbit (r(t), θ(t), pr(t), pθ(t)) of H0 in
phase space. The partial differential equation (5.32a) for
δK thus reduces to a set of ordinary differential equa-
tions that can be integrated along the individual orbits
in phase space.
The unperturbed motion for a bound orbit is in a
plane, so we can switch from spherical to plane polar co-
ordinates (r, ψ). In terms of these coordinates, we have
H0 = p
r/2+p
ψ/2, K0 = p
ψ, and cos θ = sin ι sin(ψ+ψ0),
with cos ι = Lz/
K and the constant ψ0 denoting the
angle between the direction of the periastron and the
intersection between the orbital and equatorial plane.
Then Eq. (5.32) becomes
δK = η(t), (5.33)
η(t) = − 2pψ dlSlLz
sin ι rl+2(t)
∂ψPl(sin ι sin(ψ(t) + ψ0))
cos(ψ(t) + ψ0)
2pψ clIl
rl+1(t)
∂ψPl(sin ι sin(ψ(t) + ψ0)). (5.34)
For unbound orbits, one can always integrate Eq.
(5.33) to determine δK. However, for bound periodic
orbits there is a possible obstruction: the solution for
the conserved quantity K0 + δK will be single valued if
and only if the integral of the source over the closed orbit
vanishes,
∮ Torb
η(t)dt = 0. (5.35)
Here, Torb is the orbital period. In other words, the par-
tial differential equation (5.32) has a solution δK if and
only if the condition (5.35) is satisfied. This is the same
condition as obtained by the Poincare-Mel’nikov-Arnold
method, a technique for showing the non-integrability
and existence of chaos in certain classes of perturbed dy-
namical systems [35].
Thus, it suffices to show that the condition (5.35) is
violated for all multipoles other than the spin and mass
quadrupole. To perform the integral in Eq. (5.35), we use
the parameterization for the unperturbed motion, r =
K/(1+ e cosψ) and dt/dψ = K3/2/(1+ e cosψ)2, so that
the condition for the existence of a conserved quantity
K0 + δK becomes
clIl(1 + e cosψ)
l−1∂ψPl(sin ι sin(ψ + ψ0))−
dlSlLz
K sin ι
(1 + e cosψ)l∂ψ
∂ψPl(sin ι sin(ψ + ψ0))
cos(ψ + ψ0)
(5.36)
In terms of the variable χ = ψ + ψ0 − π/2, Eq. (5.36) can be written as
dχclIl [1 + e(sinψ0 cosχ− cosψ0 sinχ)]l−1
Pl(sin ι cosχ)
dlSlLz
sin ι
[1 + e(sinψ0 cosχ− cosψ0 sinχ)]l
Pl(sin ι cosχ)
. (5.37)
Inserting the expansion (5.28) for Pl(cosχ), taking the derivatives, and using the binomial expansion for the first term
in Eq. (5.37), we get
0 = clIl
Alnjk e
j(sin ι)l−2n(sinψ0)
k(cosψ0)
dχ (sinχ)j−k+1(cosχ)k+l−2n−1
dlSlLz
Blnjk e
j(sin ι)l−2n−1(sinψ0)
k(cosψ0)
dχ (sinχ)j−k+1(cosχ)k+l−2n−2. (5.38)
The coefficients Alnkj and Blnkj are
Alnkj =
(−1)n+k+1(l − 1)!(2l − 2n)!
2ln!(l − 1− j)!k!(j − k)!(l − n)!(l − 2n− 1)!
, Blnkj =
(−1)n+kl!(2l− 2n)!
2ln!(l − j)!k!(j − k)!(l − n)!(l − 2n− 2)!
. (5.39)
The only non-vanishing contribution to the integrals in Eq. (5.38) will come from terms with even powers of both
cosχ and sinχ. These can be evaluated as multiples of the beta function:
0 = clIl
Clnjk e
j(sin ι)l−2n(sinψ0)
k(cosψ0)
j−k δ(j−k+1),even δ(l+k−1),even
dlSlLz
Dlnjk e
j(sin ι)l−2n−1(sinψ0)
k(cosψ0)
j−k δ(j−k+1),even δ(l+k),even. (5.40)
Here, the coefficients are
Clnjk =
2Γ( j
+ 1)Γ(k
− n+ 1)
Alnkj , Dlnjk =
2Γ( j
+ 1)Γ(k
− n− 1
− n+ 3
Blnkj (5.41)
Eq. (5.40) shows that for even l, terms with j =even
(odd) and k =odd (even) give a non-vanishing contribu-
tion for the case of a mass (current) multipole, and hence
K0+δK is not a conserved quantity for the perturbed mo-
tion. Note that terms with j =even and k =odd for even
l occur only for l > 3, so for l = 2 the mass quadrupole
term in Eq. (5.40) vanishes and therefore there exists an
analog of the Carter constant, which is consistent with
our results of Sec. II and our separability analysis. For
odd l, terms with j =odd (even) and k =even (odd) are
finite for Il (Sl). Note that for the case l = 1 of the spin,
the derivatives with respect to χ in Eq. (5.37) evaluate to
zero, so in this case there also exists a Carter-type con-
stant. These results show that for a general multipole
other than I2 and S1, there will not be a Carter-type
constant for such a system.
1. Exact vacuum spacetimes
Our result on the non-existence of a Carter-type con-
stant can be extended, with mild smoothness assump-
tions, to falsify the conjecture that all exact, axisymmet-
ric vacuum spacetimes posess a third constant of the mo-
tion for geodesic motion. Specifically, we fix a multipole
order l, and we assume:
• There exists a one parameter family
(M, gab(λ))
of spacetimes, which is smooth in the parameter λ,
such that λ = 0 is Schwarzschild, and each space-
time gab(λ) is stationary and axisymmetric with
commuting Killing fields ∂/∂t and ∂/∂φ, and such
that all the mass and current multipole moments of
the spacetime vanish except for the one of order l.
On physical grounds, one expects a one parameter
family of metrics with these properties to exist.
• We denote by H(λ) the Hamiltonian on the tan-
gent bundle overM for geodesic motion in the met-
ric gab(λ). By hypothesis, there exists for each λ
a conserved quantity M(λ) which is functionally
independent of the conserved energy and angular
momentum. Our second assumption is that M(λ)
is differentiable in λ at λ = 0. One would expect
this to be true on physical grounds.
• We assume that the conserved quantity M(λ) is
invariant under the symmetries of the system:
L~ξM(λ) = L~ηM(λ) = 0,
where ~ξ and ~η are the natural extensions to the 8
dimensional phase space of the Killing vectors ∂/∂t
and ∂/∂φ. This is a very natural assumption.
These assumptions, when combined with our result of
the previous section, lead to a contradiction, showing
that the conjecture is false under our assumptions.
To prove this, we start by noting that M(0) is a con-
served quantity for geodesic motion in Schwarzschild, so
it must be possible to express it as some function f of
the three independent conserved quantities:
M(0) = f(E,Lz,K0). (5.42)
Here E is the energy, Lz is the angular momentum, and
K0 is the Carter constant. Differentiating the exact re-
lation {H(λ),M(λ)} = 0 and evaluating at λ = 0 gives
{H0,M1} =
{E,H1}+
{Lz, H1}+
{K0, H1},
(5.43)
where H0 = H(0), H1 = H
′(0), and M1 = M
′(0). As
before, we can regard this is a partial differential equa-
tion that determines M1, and a necessary condition for
solutions to exist and be single valued is that the integral
of the right hand side over any closed orbit must vanish:
{E,H1}+
{Lz, H1}+
{K0, H1}
(5.44)
Now strictly speaking, there are no closed orbits in
the eight dimensional phase space. However, the ar-
gument of the previous section applies to orbits which
are closed in the four dimensional space with coordinates
(r, θ, pr, pθ), since by the third assumption above every-
thing is independent of t and φ, and pt and pφ are con-
served. Here (t, r, θ, φ) are Schwarzschild coordinates and
(pt, pr, pθ, pφ) are the corresponding conjugate momenta.
Next, we can pull the partial derivatives ∂f/∂E etc.
outside of the integral. It is then easy to see that the first
two terms vanish, since there do exist a conserved energy
and a conserved z-component of angular momentum for
the perturbed system. Thus, Eq. (5.44) reduces to
{K0, H1} = 0. (5.45)
Since M(0) is functionally independent of E and Lz, the
prefactor ∂f/∂K0 must be nonzero, so we obtain
{K0, H1} = 0. (5.46)
The result (5.46) applies to fully relativistic orbits in
Schwarzschild. We need to take the Newtonian limit of
this result in order to use the result we derived in the
previous section. However, the Newtonian limit is a lit-
tle subtle since Newtonian orbits are closed and generic
relativistic orbits are not closed. We now discuss how the
limit is taken.
The integral (5.46) is taken over any closed orbit in
the four dimensional phase space (r, θ, pr, pθ) which cor-
responds to a geodesic in Schwarzschild. Such orbits are
non generic; they are the orbits for which the ratio be-
tween the radial and angular frequencies ωr and ωθ is a
rational number. We denote by qr and qθ the angle vari-
ables corresponding to the r and θ motions [36]. These
variables evolve with proper time τ according to
qr = qr,0 + ωrτ, (5.47a)
qθ = qθ,0 + ωθτ, (5.47b)
where qr,0 and qθ,0 are the initial values. We denote the
integrand in Eq. (5.46) by
I(qr, qθ, a, ε, ι),
where I is some function, and a, ε and ι are the parame-
ters of the geodesic defined by Hughes [32] (functions of
E, Lz and K0). The result (5.46) can be written as
∫ T/2
dτ I[qr(τ), qθ(τ), a, ε, ι] = 0, (5.48)
where T = T (a, ε, ι) is the period of the r, θ motion.
Since the variables qr and qθ are periodic with period
2π, we can express the function I as a Fourier series
I(qr, qθ, a, ε, ι) =
n,m=−∞
Inm(a, ε, ι)einqr+imqθ . (5.49)
Now combining Eqs. (5.47), (5.48) and (5.49) gives
n,m=−∞
Inm(a, ε, ι)einqr,0+imqθ,0
×Si [(nωr +mωθ)T/2] , (5.50)
where Si(x) = sin(x)/x. Since the initial conditions qr,0
and qθ,0 are arbitrary, it follows that
Inm(a, ε, ι)Si [(nωr +mωθ)T/2] = 0 (5.51)
for all n, m.
Next, for closed orbits the ratio of the frequencies must
be a rational number, so
, (5.52)
where p and q are integers with no factor in common.
These integers depend on a, ε and ι. The period T is
given by 2π/T = qωr = pωθ. The second factor in Eq.
(5.51) now simplifies to
(np+mq)π
, (5.53)
which vanishes if and only if
n = n̄q, m = m̄p, n̄+ m̄ 6= 0, (5.54)
for integers n̄, m̄. It follows that
Inm(a, ε, ι) = 0 (5.55)
for all n, m except for values of n, m which satisfy the
condition (5.54)
Consider now the Newtonian limit, which is the limit
a → ∞ while keeping fixed ε and ι and the mass of the
black hole. We denote by IN(qr , qθ, a, ε, ι) the Newtonian
limit of the function I(qr , qθ, a, ε, ι). The integral (5.48)
in the Newtonian limit is given by the above computation
with p = q = 1, since ωr = ωθ in this limit. This gives
dτIN =
INn,−n(a, ε, ι) ein(qr,0−qθ,0), (5.56)
where INnm are the Fourier components of IN. In the
previous subsection, we showed that this function is non-
zero, which implies that there exists a value k of n for
which IN k,−k 6= 0.
Now as a → ∞, we have ωr/ωθ → 1, and hence from
Eq. (5.52) there exists a critical value ac of a such that
the values of p and q exceed k for all closed orbits with
a > ac. (We are keeping fixed the values of ε and ι). It
follows from Eqs. (5.54) and (5.55) that
Ik,−k(a, ε, ι)
IN k,−k(a, ε, ι)
= 0 (5.57)
for all such values of a. However this contradicts the fact
Ik,−k(a, ε, ι)
IN k,−k(a, ε, ι)
→ 1 (5.58)
as a→ ∞. This completes the proof.
Hence, if the three assumptions listed at the start of
this subsection are satisfied, then the conjecture that all
vacuum, axisymmetric spacetimes possess a third con-
stant of the motion is false.
Finally, it is sometimes claimed in the classical dynam-
ics literature that perturbation theory is not a sufficiently
powerful tool to assess whether the integrability of a sys-
tem is preserved under deformations. An example that
is often quoted is the Toda lattice Hamiltonian [38, 39].
This system is integrable and admits a full set of con-
stants of motion in involution. However, if one approx-
imates the Hamiltonian by Taylor expanding the poten-
tial about the origin to third order, one obtains a sys-
tem which is not integrable. This would seem to indicate
that perturbation theory can indicate a non-integrability,
while the exact system is still integrable.
In fact, the Toda lattice example does not invalidate
the method of proof we use here. If we write the Toda
lattice Hamiltonian as H(q,p), then the situation is that
H(λq,p) is integrable for λ = 1, but it is not integrable
for 0 < λ < 1. Expanding H(λq,p) to third order in λ
gives a non-integrable Hamiltonian. Thus, the perturba-
tive result is not in disagreement with the exact result
for 0 < λ < 1, it only disagrees with the exact result for
λ = 1. In other words, the example shows that pertur-
bation theory can fail to yield the correct result for finite
values of λ, but there is no indication that it fails in ar-
bitrarily small neighborhoods of λ = 0. Our application
is qualitatively different from the Toda lattice example
since we have a one parameter family of Hamiltonians
H(λ) which by assumption are integrable for all values
of λ.
VI. CONCLUSION
We have examined the effect of an axisymmetric
quadrupole moment Q of a central body on test parti-
cle inspirals, to linear order in Q, to the leading post-
Newtonian order, and to linear order in the mass ratio.
Our analysis shows that a natural generalization of the
Carter constant can be defined for the quadrupole inter-
action. We have also analyzed the leading order spin self-
interaction effect due to the scattering of the radiation off
the spacetime curvature due to the spin. Combining the
effects of the quadrupole and the leading order effects
linear and quadratic in the spin, we have obtained ex-
pressions for the instantaneous as well as time-averaged
evolution of the constants of motion for generic orbits un-
der gravitational radiation reaction, complete at O(a2ǫ4).
We have also shown that for a single multipole interaction
other than Q or spin, in our approximations, a Carter-
type constant does not exist. With mild additional as-
sumptions, this result can be extended to exact space-
times and falsifies the conjecture that all axisymmetric
vacuum spacetimes possess a third constant of motion for
geodesic motion.
VII. ACKNOWLEDGMENTS
This research was partially supported by NSF grant
PHY-0457200. We thank Jeandrew Brink for useful cor-
respondence.
Appendix A: Time variation of quadrupole: order of
magnitude estimates
In this appendix, we give an estimate of the timescale
Tevol for the quadrupole to change. The analysis in the
body of this paper is valid only when Tevol ≫ Trr, where
Trr is the radiation reaction time, since we have neglected
the time evolution of the quadrupole. We distinguish be-
tween two cases: (i) when the central body is exactly non-
spinning but has a quadrupole, and (ii) when the central
body has finite spin in addition to the quadrupole.
1. Estimate of the scaling for the nonspinning case
For the purpose of a crude estimate, the relevant in-
teraction is the tidal interaction with energy
QijEij ∼ −
Q̄I cos2 θ, (A1)
where Eij is the tidal field, θ is the angle between the
symmetry axis and the normal to the orbital plane of
m2, and we have written the quadrupole as Q ∼ Q̄I,
where Q̄ is dimensionless and I is the moment of inertia.
For small deviations from equilibrium, the relevant piece
of the Lagrangian is schematically
L ∼ Iψ̇2 + Q̄I m2
ψ2. (A2)
We define the evolution timescale Tevol to be the time
it takes for the angle to change by an amount of order
unity, and since the amplitude of the oscillation scales
roughly as ∼ m2/m1, the evolution time scales as
T−2evol ∼
ω2orbit, (A3)
where ω2orbit = M/r
3. Thus, the ratio of the evolution
timescale compared to the radiation reaction timescale
scales as
Tevol/Trr ∼
. (A4)
2. Estimate of the scaling for the spinning case
When the body is spinning the effect of the tidal cou-
pling is to cause a precession. For the purpose of this
estimate, we calculate the torque on m1 due to the com-
panion’s Newtonian field. The torque N scales as
Ni ∼ ǫimjQmkEjk. (A5)
We assume that the precession is slow, i.e.
ωprec ≪ S̄/m1
, (A6)
where ωprec is the precession frequency and S̄ = S/m
is the dimensionless spin. This gives the approximate
scaling of the precession timescale as (cf. [37])
Tprec/Trr ∼
. (A7)
and the evolution timescale is thus
Tevol/Trr ∼
. (A8)
Because of our assumption (A6) that the precession is
slow, equation (A8) is valid only when
) S̄2
. (A9)
When S̄ is sufficiently small that the condition (A9) is
violated, the relevant timescale is instead given by Eq.
(A3).
3. Application to Kerr inspirals
For Kerr inspirals,
S̄ ∼ a, Q̄ ∼ a2, µ/M ≪ 1 and r ∼M. (A10)
Therefore, the condition (A9) is satisfied, and the pre-
cession time is longer than the radiation reaction time
Tprec/Trr ∼
. (A11)
Note that for Kerr inspirals, since r ∼ M both formulas
(A3) and (A7) give the same scaling.
Moreover, for Kerr inspirals, the amplitude of the pre-
cession will be small, of order the mass ratio µ/M . This is
because of angular momentum conservation: in the rela-
tivistic regime, the orbital angular momentum is a factor
of µ/M smaller than the angular momentum of the black
hole and can therefore not cause a large precession ampli-
tude. Even if the orbital angular momentum at infinity
is large, most of it will be radiated away as outgoing
gravitational waves during the earlier phase of the inspi-
ral. This factor of µ/M is taken into account when we
consider the evolution timescale, which for Kerr inspirals
reduces to
Tevol/Trr ∼
. (A12)
Since 1/a ≥ 1, M/r ∼ 1 and M/µ ≪ 1, the evolution
time is long compared to the radiation reaction time and
we can neglect the time variation of the quadrupole at
leading order.
Appendix B: Computation of time averaged fluxes
1. Averaging method that parallels fully
relativistic averaging
We start by noting that the differential equations
(2.26) and (2.27) governing the r̃ and θ̃ motions decouple
if we define a new time parameter t̂ by
dt̂ =
dt̃. (B1)
This is the analog of the Mino time parameter for
geodesic motion in Kerr [12]. The equations of motion
(2.26)–(2.24) then become
= V̂r̃(r̃), (B2)
V̂r̃(r̃) = 2Er̃
4 + 2r̃3 −Kr̃2 − 4SLzr̃
r̃ − 2L2z
, (B3)
= V̂θ̃(θ̃), (B4)
V̂θ̃(θ̃) = K −
sin2 θ̃
−QE cos 2θ̃, (B5)
= V̂ϕr̃(r̃) + V̂ϕθ̃(θ̃), (B6)
V̂ϕr̃(r̃) =
, V̂ϕθ̃(θ̃) =
sin2 θ̃
. (B7)
The parameters t and t̂ are related by:
= V̂tr̃(r̃) + V̂tθ̃(θ̃) (B8)
V̂tr̃(r̃) = r̃
2, V̂tθ̃(θ̃) =
cos 2θ̃. (B9)
It follows from Eqs. (B2) and (B4) that the functions
r̃(t̂) and θ̃(t̂) are periodic; and we denote their periods
by Λr̃ and Λθ̃. We define the fiducial motion associated
with the constants of motion E, Lz and K to be the
motion with the initial conditions r̃(0) = r̃min and θ̃(0) =
θ̃min, where r̃min and θ̃min are given by the vanishing of
the right-hand sides of Eqs. (B2) and (B4) respectively.
The functions r̂(t̂) and θ̂(t̂) associated with this fiducial
motion are given by
∫ r̂(t̂)
r̃min
V̂r̃(r̃)
= t̂, (B10)
∫ θ̂(t̂)
θ̃min
V̂θ̃(θ̃)
= t̂. (B11)
From Eq. (B8) it follows that
t(t̂) = t0 +
V̂tr̃[r̃(t
′)] + V̂tθ̃[θ̃(t
, (B12)
where t0 = t(0). Next, we define the constant Γ to be
the following average value:
∫ Λr̃
dt′V̂tr̃[r̂(t
′)] +
dt′V̂tθ̃[θ̂(t
′)]. (B13)
Then we can write t(t̂) as a sum of a linear term and
terms that are periodic:
t(t) = t0 + Γt̂+ δt(t̂), (B14)
where δt(t̂) denotes the oscillatory terms in Eq. (B12).
To average a function over the time parameter t̂, it is
convenient to parameterize r̃ and θ̃ in terms of angular
variables as follows. For the average over θ̃ we introduce
the parameter χ by
cos2 θ̂(t̂) = z− cos
2 χ, (B15)
where z− = cos
2 θ̃− with z− being the smaller root of Eq.
(B4):
K + 3QE ±
(K −QE)2 + 4QEL2z
(B16)
and where β = 2QE. Then from the definition (B11)
of θ̂ together with Eq. (B4) and the requirement that χ
increases monotonically with t̂ we obtain
β (z+ − z− cos2 χ). (B17)
Then we can write the average over t̂ of a function Fθ̃(t̂)
which is periodic with period Λθ̃ in terms of χ as
〈Fθ̃〉t̂ =
dt̂Fθ̃(t̂)
Fθ̃[t̂(χ)]
β (z+ − z− cos2 χ)
, (B18)
where
Λθ̃ =
β (z+ − z− cos2 χ)
. (B19)
Similarly, to average a function Fr̃(t̂) that is periodic with
period Λr̃, we introduce a parameter ξ via
1 + e cos ξ
, (B20)
where the parameter ξ varies from 0 to 2π as r̃ goes
through a complete cycle. Then,
= P (ξ), (B21)
P (ξ) ≡
V̂r̃[r̃(ξ)]
pe | sin ξ |
(1 + e cos ξ)
(B22)
The average over t̂ of Fr̃(t̂) can then be computed from
〈Fr̃〉t̂ =
dξ Fr̃/P (ξ)
dξ/P (ξ)
. (B23)
Now, a generic function Fr̃,θ̃[r̃(t̂), θ̃(t̂)] will be biperiodic
in t̂: Fr̃,θ̃[r̃(t̂+Λr̃), θ̃(t̂+Λθ̃)] = Fr̃,θ̃[r̃(t̂), θ̃(t̂)]. Combin-
ing the results (B18) and (B23) we can write its average
as a double integral over χ and ξ as
〈Fr̃,θ̃〉t̂ =
Λθ̃Λr̃
Fr̃,θ̃[r̃(ξ), θ̃(χ)]
β (z+ − z− cos2 χ)P (ξ)
(B24)
To compute the time average of Ė, L̇z, and K̇, we need
to convert the average of a function over t̂ calculated from
(B24) to the average over t. As explained in detail in
[9], in the adiabatic limit we can choose a time interval
∆t which is long compared to the orbital timescale but
short compared to the radiation reaction time. From
Eq. (B12) we have ∆t = Γt̂+ osc.terms. The oscillatory
terms will be bounded and will therefore be negligible in
the adiabatic limit, so we have to a good approximation
〈Ė〉t =
〈Ė V̂t〉t̂, (B25)
where V̂t ≡ V̂tr̃ + V̂tθ̃, cf. Eq. (B8), and similarly for L̇z
and K̇.
The explicit results we obtain using this method are
given in section III, Eqs. (3.28), (3.29), and (3.30).
2. Averaging method using the explicit
parameterization of Newtonian orbits
To perform the time-averaging using this method, we
define a parameter ξ via
1 + e cos ξ
, (B26)
where the parameter ξ varies from 0 to 2π as r̃ goes
through a complete cycle. Note that θ appears in Eqs.
(3.16) – (3.18) only in terms that are linear in Q, so we
can write θ in terms of ξ using the Newtonian relation
x3 = r cos θ = r sin ι sin(ξ + ξ0). (B27)
Here, ξ0 is the angle between the direction of the peri-
helion and the intersection of the orbital and equatorial
plane. Similarly, for the ṙθ̇ terms in Eqs. (3.17) and
(3.26) we can use the Newtonian relations ṙ = e/
p sin ξ
and ξ̇ =
p/r2. From Eqs. (2.27) and (B20) it follows
(1 + e cos ξ)2
−3 + e2 − 2e cos ξ + 2 cos2 ι(8 − e2 + 8e cos ξ + e2 cos 2ξ)
, (B28)
and from Eq. (2.12)
(1 + e cos ξ)
2 sin2 ι sin2(ξ + ξ0)− 1
. (B29)
Using these expressions, we compute the time-averaged fluxes from
〈Ė〉 =
dξ Ė (dt/dt̃) (dt̃/dξ)
dξ (dt/dt̃) (dt̃/dξ)
(B30)
and obtain:
〈Ė〉 = −32
(1− e2)3/2
e4 − S
cos(ι)
cos(2ι)
cos(2ι)
cos(2ξ0) sin
cos(2ξ0) sin
, (B31)
〈L̇z〉 = −
(1− e2)3/2
cos ι
e2 − S
2p3/2 cos ι
+ 7e2 +
cos(2ι)
−3− 45
45 + 148e2 +
cos(2ι)
1 + 3e2 +
e2 cos(2ξ0) sin
, (B32)
〈K̇〉 = −64
(1 − e2)3/2
e2 − S
2p3/2
+ 37e2 +
cos(ι)
cos(2ι)
cos(2ι)
e2 cos(2ξ0) sin
. (B33)
In the adiabatic limit, the terms involving cos(2ξ0) can
be omitted because they average to zero. As explained
by Ryan [15], the radiation reaction timescale for terms
involving ξ0 is much longer than the precession timescale
for most orbits, so the terms involving ξ0 will average
away. This is consistent with our results for the adia-
batic infinite time-averaged fluxes using the Mino time
parameter. The Mino-time averaging method was based
on the assumption that the fundamental frequencies are
incommensurate and the motion fills up the whole torus,
which is equivalent to averaging over ξ0.
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|
0704.0392 | Simulation of Robustness against Lesions of Cortical Networks | Simulation of Robustness against Lesions of Cortical Networks
Abbreviated title: Simulation of Robustness of Cortical Networks
Marcus Kaiser1,2,3,a, Robert Martin2,4,a, Peter Andras1,2 and Malcolm P. Young2
1 School of Computer Science, University of Newcastle, Claremont Tower, Newcastle upon Tyne,
NE1 7RU, UK
2 Henry Wellcome Building for Neuroecology, Institute of Neuroscience, University of Newcastle,
Framlington Place, Newcastle upon Tyne, NE2 4HH, UK
3 Jacobs University Bremen, School of Engineering and Science, Campus Ring 6, 28759 Bremen,
Germany
4 FR 2-1, NI, Informatik, Technische Universität Berlin, Franklinstr. 28/29, 10587 Berlin, Germany
aAuthors contributed equally to this paper
Correspondence: Marcus Kaiser; School of Computer Science, University of Newcastle, Claremont
Tower, Newcastle upon Tyne, NE1 7RU, UK; E-mail: [email protected]
Keywords: cat, macaque monkey, resilience, scale-free networks, small-world networks
ABSTRACT
Structure entails function and thus a structural description of the brain will help to understand its
function and may provide insights into many properties of brain systems, from their robustness and
recovery from damage, to their dynamics and even their evolution. Advances in the analysis of
complex networks provide useful new approaches to understanding structural and functional
properties of brain networks. Structural properties of networks recently described allow their
characterization as small-world, random (exponential) and scale-free. They complement the set of
other properties that have been explored in the context of brain connectivity, such as topology,
hodology, clustering, and hierarchical organization. Here we apply new network analysis methods
to cortical inter-areal connectivity networks for the cat and macaque brains. We compare these
corticocortical fibre networks to benchmark rewired, small-world, scale-free and random networks,
using two analysis strategies, in which we measure the effects of the removal of nodes and
connections on the structural properties of the cortical networks. The brain networks’ structural
decay is in most respects similar to that of scale-free networks. The results implicate highly
connected hub-nodes and bottleneck connections as structural basis for some of the conditional
robustness of brain systems. This informs the understanding of the development of brain networks’
connectivity.
INTRODUCTION
The brain can be remarkably robust to physical damage. Significant loss of neural tissue can be
compensated for in a relatively short time by large-scale adaptation of remaining brain parts (e.g.,
Spear et al., 1988; Stromswold, 2000; Young, 2000). On the other hand, the removal of small
amounts of tissue (e.g. in Broca’s area) can lead to a severe functional deficit. These findings
provide a somewhat contradictory picture of the robustness of the brain and suggest a number of
questions. Can we evaluate effective robustness given this variability in the effects of brain lesions?
Are severity and nature of effects of localized damage predictable? We assess here how connectivity
data of brain area connectivity can be brought to bear on these questions.
The functionality of any system is grounded in its structural properties. For neurosciences, this has
led to exploration of the structural properties of brain networks, such as topology, hodology,
clustering, and hierarchical organization (e.g., Nicolelis et al., 1990; Felleman and van Essen, 1991;
Young, 1992; Young et al., 1994; Hilgetag et al., 1996; Hilgetag et al., 2000; Sporns et al., 2000;
Young, 2000; Petroni et al., 2001; Sporns et al., 2004; Kaiser and Hilgetag, 2006). Recent advances
in the study of networks have extended those traditional structural descriptions (Strogatz, 2001),
allowing to characterize networks as small-world (Watts and Strogatz, 1998), random and scale-free
(e.g., Barabási and Albert, 1999; Albert et al., 2000).
Small-world networks comprise well-connected local neighbourhoods with fewer long-range
connections between neighbourhoods. The length of a path between two nodes, that is the number of
connections that have to be crossed to go from one node to another, is comparable as low as for a
randomly organized network. Scale-free networks are characterized by their specific distribution of
connectivities or degrees—the number of connections that each node has. The degree distribution
follows a power law. Whereas these networks can have highly-connected nodes or hubs also
networks where nodes have maximally 20 connections have been described as scale-free based on
the power-law degree distribution (Jeong et al., 2001). Small-world and scale-free properties are
compatible, but not equivalent (see, e.g., Amaral et al., 2000).
Scale-free networks have higher robustness than random ones against randomly located damage,
whilst being sensitive to damage targeted at their most widely connected nodes (Barabási and Albert,
1999; Young, 2000). This is reminiscent of the properties of the brain described above.
Previous studies have shown that functional networks of the human brain are scale-free (Eguiluz et
al., 2005). However, at the level of resting state networks between cortical areas, it was argued that
these networks are not scale-free (Achard et al., 2006). Here we analyze what pattern occurs at the
level of structural connectivity. In order to establish, whether the brain has properties of scale-free
networks, the integrity and robustness to damage of brain networks’ structure is compared to that of
benchmark random and scale-free networks (an earlier version of this work had been presented as a
conference abstract, see Martin et al., 2001).
MATERIALS AND METHODS
Brain structure connectivity data. We used macaque and cat cortical inter-areal connectivity data
(Young, 1993; Young et al., 1994; Scannell et al., 1995) applying the CoCoMac database for the
primate data (Stephan et al., 2001; Kötter, 2004). In both species, the data comprised connections
among cortical regions of the neocortex.
For the macaque brain, we considered 66 brain structures with 608 connections between them. In the
case of the cat brain, we considered 56 structures and 814 connections. We excluded cross-
hemispheric connections. The data was represented as the binary connectivity matrix of a graph.
Nodes corresponded to the considered brain structures and edges to the reported connections
between them. Note that due to the directed nature of brain connections, the connectivity matrix is
not necessarily symmetric and the resulting graph has hence directed edges too. The edge density of
the macaque brain graph, that is, the number of reported connections divided by the number of all
possible connections, is 26.4% (Tab. 1). For the cat brain, the edge density is 14.2%. There are on
average 9.2 connections for each structure in the macaque brain and 14.5 connections for the cat
brain structures (see supplementary material for the connectivity matrices).
[Table 1 near here]
Benchmark networks for comparison. We constructed rewired, scale-free, small-world, and random
networks to match the number of nodes and connections of the corresponding two brain networks
(Tab. 1). Figure 1 shows examples of small random and scale-free networks to demonstrate
differences in their topology. For random networks, the number of connections of a node is close to
the average value over all nodes. For scale-free networks, however, nodes with a much higher
number of connections can occur; see hub in Fig. 1b.
[Figure 1 near here]
Rewired networks. For rewired networks, each node has the same number of connections as in the
original network, however, targets or sources of connections might have changed.Rewired networks
were derived from the original cortical networks of the cat and macaque by exchanging connections
so that the total number of connections of each node remained the same (the method for
randomization is described in Milo et al., 2002). Whereas the degree distribution remains unchained,
the cluster architecture is lost during rewiring. Thus, rewired networks allow looking at effects of the
degree distribution alone.
Scale-free networks. The algorithm to generate scale-free benchmark networks is based on Barabási
and Albert (1999). However, in a modification of their approach we began with an initial graph of
six and eight fully connected nodes respectively for the macaque and cat benchmark networks. This
was necessary in order to ensure that the clustering coefficient (average percentage of connections
between neighbours of a node; see definition below) of the initial graph matched the highest
clustering coefficient found in the corresponding brain network. As proposed by Barabási and Albert
(1999), further nodes were added one by one to the graph by preferential attachment. At the
beginning of this process, the probability that a new node is connected to an existing node i is
iP )( ,
where kj is the number of connections of the node j (Barabási and Albert, 1999). After establishing a
connection to node i*, the probabilities are recalculated to reflect the nature of the scale-free
networks: if i is connected to j, then it is more likely that i is connected to nodes which are already
connected to j and it is less likely that i is connected to nodes which are not connected to j. The
rescaling was undertaken according to
connectednot areandif,)(
connectedareandif),(
iiiPk
The probability for the connections in both directions is the same. We confirmed that this modified
routine for generating scale-free networks was able to yield a power-law degree distribution (cf.
supplementary material).
Small-world networks. Small-world networks were generated by rewiring regular networks as
described in the literature (Watts and Strogatz, 1998). The rewiring probability was adjusted so that
the resulting networks had similar clustering coefficient than the respective cortical networks (Tab.
Random networks. Whereas all benchmark networks are generated by a random process, we denote
Erdös-Renyi random networks (Erdös and Rényi, 1960) as random networks in the remaining
manuscript. Random networks were generated by establishing each potential connection between
nodes with probability p. This probability was the desired connection density, that means, the
connection density of the corresponding brain networks, 14.2% of the number of all possible
connections for the macaque and 26.4% for the cat. The degree distribution in these random
networks followed a binomial probability distribution. For large numbers of nodes this can be
approximated by a Poisson distribution and hence the term ‘exponential degree distribution’ is also
used (Bollobas, 1985).
Graph similarity. To assess the discrepancy in connectivity between two networks, first their nodes
are permuted according to their number of connections. Second, permutated cortical and benchmark
matrices are compared by looking what ratio of directed edges in the adjacency matrix that occurred
at the same position in both matrices and the total number of directed edges. This percentage is then
the graph similarity S between graph A and B given the number of (directed) connections |E|:
where is element-by-element multiplication with an element in the resulting matrix non-zero if
both elements are non-zero; Σ is the sum of all elements in the matrix and thus yields the number of
directed edges existing in both matrices, as these are denoted by a value of one in the matrix. Note,
that benchmark networks could be more similar than they appear for this measure as not all possible
arrangements of nodes were tested. Testing all possibilities (1092 for the macaque and 1074 for the
cat) would have been computationally unfeasible.
Network characterisation. The clustering coefficient shows the fragmentation of the network. The
coefficient is the ratio of the number of existing edges between neighbours of a node i and the
number of possible edges between all these neighbours. We considered neighbouring nodes of node i
to be all those nodes that have incoming or outgoing connections between them and node i. If a node
i has ki neighbours, then the number of all possible in- and outgoing edges between the neighbouring
nodes is ki
* (ki – 1). The coefficient itself is a local property of each node and we define the average
coefficient of all nodes to be the clustering coefficient of the graph. This is a measure of how well
connected the nodes of the network are.
Following Albert et al. (2000), we considered the average shortest path (ASP) or characteristic path
length to characterize the network connectivity and integrity. The ASP between any two nodes in the
network is the number of sequential connections that are necessary, on average, to link one node to
another by the shortest possible route (Diestel, 1997). In case a network becomes disconnected in the
process of removing edges/nodes and there is no path between two nodes, the pair of nodes is
ignored. If no two connected points are left, the average shortest path is set to zero. We used Floyd’s
algorithm to determine the matrix of the shortest paths between each pair of nodes (Cormen et al.,
2001). Note that due to directed edges, the shortest path from node i to node j may not be the same
as that from node j to node i.
Target determination. In order to determine the importance of a node to the overall network
structure, a simple metric has been used, namely the number of connections formed by this node. In
experiments requiring the targeted removal of nodes from the networks, the most highly connected
node was eliminated.
To provide the corresponding metric for the targeted elimination of connections (edges) from the
network, we chose edge betweenness (Girvan and Newman, 2002), that is, the number of shortest
paths between all pairs of nodes that pass though the edge. Edges with high edge betweenness are
chosen for targeted attack. Indeed, edge betweenness has been shown to highly correlate with
structural network damage for cortical as well as other biological networks (Kaiser and Hilgetag,
2004).
Analysis methods. We used the iterative random and targeted removal of nodes and connections to
analyze the robustness of the networks against damage. Random removal means that we selected a
node or connection and deleted it from the graph irrespective of the degree of the node. In the case of
targeted removal, we selected the most important node or connection left in the network (see above).
After each deletion, we calculated the ASP of the resulting graph. We continued the removal of
nodes or connections until all nodes were removed from the network. To derive estimates of the
variability in these connectivity measures, we considered 50 benchmark networks for each condition.
In the cases of random removal, we repeated the analysis for the brain networks 50 times as well.
RESULTS
Degree distribution of cortical networks
Fig. 2 shows the degree distributions of macaque and cat compared with a distribution of random
networks. In comparison to random networks, the macaque cortical network has highly connected
nodes but also more sparsely connected nodes, reminiscent of scale-free networks. This is also true
for the cat network that shows a remarkable number of areas with few connections compared to
random networks. Table 2 shows the five most-highly connected nodes for the cat and macaque
networks.
[Table 2 near here]
The standard way of observing whether the cortical network resembles a scale-free network would
be to search for a power-law in the degree distribution. However, this approach would be
inappropriate for cortical networks for three reasons. First, the maximum number of connections of a
node equals the number of regions in the network minus one, that means, 65 (macaque) or 55 (cat).
Therefore, the degree distribution only consists of two scales. Second, where degree distributions
with a low maximal degree had been studied before (Jeong et al., 2001), the number of nodes was
considerably higher (>1,800). As less than 100 degrees form the degree distribution, results are
unlikely to be robust. Third, there exists a sampling problem in that the amount of unknown or not
included connections might change the shape of the degree distribution (Stumpf et al., 2005).
Therefore, we will use indirect measures to determine whether cortical networks are similar to scale-
free networks.
[Figure 2 near here]
Graph similarity
Whereas the degree distribution is an abstraction of the underlying network, we looked at a direct
comparison between the cortical and benchmark networks. Whereas a direct measure of network
similarity was computationally unfeasible (see Methods), we compared the adjacency matrices after
ordering nodes by their degrees (see methods). We then looked at the similarity of cortical networks
with different benchmark networks (Fig. 3). For rewired cortical networks, the percentage of
identical edges was 23% for rewired macaque and 38% for the rewired cat network. Interestingly,
benchmark scale-free networks are as similar to the cortical networks as the rewired cortical
networks. In contrast, the similarity of random and small-world networks is significantly lower. This
can be attributed to the degree distribution of scale-free and cortical networks being comparable as
the rewired network only has the degree distribution in common with the original cortical network.
After these structural properties, we tested the effect of topological changes on general network
properties.
[Figure 3 near here]
Sequential elimination of nodes
We tested the influence of sequential node elimination on the network structure. Nodes were
removed one by one from the network, either randomly or targeted. Plotting the ASP as a function of
the fraction of deleted nodes illustrates the characteristic structural disintegration of each network
type (See Fig. 4 for the example of targeted elimination of nodes from the Macaque benchmark
networks. The complete set of curves for the different analysis types is available as supplementaryt
material).
[Figure 4 near here]
Fig. 4A illustrates the effect of random and targeted removal of nodes from the Macaque brain
network. Clearly, the specific decline in ASP is different for the two analysis strategies. Whilst the
random removal causes only a slow rise in the ASP, targeted removal of highly connected nodes has
a much stronger effect on the network structure of the brain network. After a steep rise in ASP the
network fragments into smaller components. The remaining shortest paths, that is the paths between
nodes within components, are smaller than in the original network. This process leads to a network
with pairs of nodes that are connected to each other but not to other nodes of the network. In these
cases, the shortest path decreases to a value of one. Finally, also nodes within pairs are removed
leading to an ASP of zero.
Fig. 4B–C contrast this specific curve to those observed when removing nodes from the different
benchmark networks in a targeted fashion. Whilst the ASP in the random and small-world networks
is hardly affected by the targeted elimination of a large proportion of nodes, in the scale-free, like in
the brain networks, the effect of targeted node elimination manifests itself in a sharp rise in this
measure. Moreover, both, the scale-free and the brain networks show a decline in the ASP around
the fraction of deletions, and the characteristic behaviour of the brain network is within the 95%
confidence interval encountered for the set of scale-free benchmark networks. This is not the case for
the other benchmark networks considered (see Fig. 4).
For the cat brain network (Fig. 5), the random and small-world networks show a different behaviour
for targeted node removal than the original cortical network. Though the cat response to targeted
node removal is largely within the 95% confidence interval for the scale-free benchmark networks,
the peak ASP value and the fraction of deleted nodes where the peak occurs is lower for the cat
cortical network.
[Figure 5 near here]
The decline in ASP at a later stage during the elimination process, as observed for the brain and
scale-free networks may appear unusual and deserves some additional attention. It can have two
reasons. First, it could be that the network gets fragmented into different disconnected components.
Each of these is smaller, and likely to have a shorter ASP. Second, the overall decrease in network
size with successive eliminations can lead to a decrease in shortest path. This is, however, likely to
be a slow process, as it will usually be offset by an increase in ASP due to the targeted nature of the
elimination.
In order to quantitatively compare the different graphs, we consider two measures. The first is the
maximal ASP measured during the removal of nodes; the second is the fraction of deleted nodes, for
which the peak ASP occurs (Fig. 6). For the fraction of peak ASP, only the scale-free benchmark
networks are close to the cortical fraction whereas all other benchmark networks show significantly
higher fractions. This means that both in the cortical as well as the scale-free networks the removal
of highly-connected nodes leads to a rapid increase of ASP so that the fraction of deleted nodes at
which the maximum ASP occurs is earlier than for other networks. However, the peak value for
scale-free networks is greater than that for cortical networks.
[Figure 6 near here]
Sequential elimination of connections
We also tested the similarity of sequential connection elimination. Connections were eliminated one
after another either randomly or targeted. Full details of the networks disintegration are shown in the
supplementary material. Again we compare the maximal ASP measured during the removal of
connections and the fraction of deleted connections, for which the peak ASP occurs (Fig. 7).
[Figure 7 near here]
Only the scale-free benchmark networks yield similar values for both the cat and macaque network
whereas other networks yield similar values for just one of the cortical networks.
DISCUSSION
We have compared brain inter-area connectivity networks with different types of benchmark
networks, including random, scale-free, and small-world networks, and found strong indications that
the brain connectivity networks share some of their structural properties with scale-free networks.
Besides a formal assessment of the network connectivity (degree distribution and graph similarity,
see Figures 3 and 4), the analysis is based on a novel approach, which measures the effect of
removal of components of the different networks on their structural integrity. In particular, we
compared the effect that the removal of nodes and connections had on the ASP found in the brain
connectivity networks and their benchmark counterparts. Note, however, that this analysis is based
on cortical connectivity within one hemisphere. Connections between hemispheres and between the
cortex and subcortical structures such as thalamic regions were not included. The reason for the lack
of interhemispheric connections was that few tracing studies tested for and thus reported fibre tracts
towards the contralateral hemisphere. Whereas information about thalamocortical connections would
have been available, regions with available information about fibre tracts differed between the cat
and macaque. To be consistent between species, the data was not included. For each species, an
inclusion of these regions yielded similar results concerning the removal of nodes or edges
(supplementary material).
Simulated robustness and its relation to lesion studies
How do our simulations relate to experimental lesion studies? Node elimination corresponds
indirectly to inactivation or lesion of the corresponding brain areas, and from this perspective, we
can interpret this analysis in terms of the brain’s robustness to regional damage. The elimination of
connections corresponds indirectly to localized brain lesions that damage the white matter and
interrupt communication between normally connected brain structures. The ASP yields a measure
how well the brain is connected and how well different streams of information can be integrated.
Analysing the spatial organisation of cortical networks shows that the brain is optimized towards a
low ASP (Kaiser and Hilgetag, 2006). A recent clinical study of the EEG correlation network in
Alzheimer patients suggests that increases in ASP lead to a reduced performance in memory tasks
(Stam et al., 2007). In this study, the ASP of the EEG synchronization network has been higher in
Alzheimer patients compared to the control group. Furthermore, there was a negative correlation
between the patients’ ASP and their performance in a standard clinical memory test. Whereas the
study was based on functional rather structural/anatomical networks, recent studies using diffusion
tensor imaging have shown that changes in brain connectivity can be linked to diseases such as
Schizophrenia and Alzheimer.
All observations have been made equally during the analysis of the brain networks of cat and
macaque, despite different edge densities in the two networks. It is therefore prudent to conclude that
it may be extended to other mammalian brain networks. Hence, conditional robustness of brain
function may be based to a large extent on two fairly simple structural properties of brain networks:
firstly, the number of connections of individual nodes (Young, 2000), i.e., their scale-free nature,
and secondly, the heavily connected local clusters with fewer important ‘bottlenecks’ between them
(Kaiser and Hilgetag, 2004). Consequently, it appears feasible to determine the brain structures that
are the most important to the maintenance of network function. Typically, brain networks should be
able to function robustly in the face of damage to structures that have few connections and damage
to connections that do not form part of many shortest connections between pairs of areas. On the
other hand, functional effects should be dramatic when structures with very many connections (hubs)
are damaged and when connections between structures with very different connectivity patterns
(large edge betweenness, cf. Girvan and Newman, 2002) are damaged.
Is the brain a scale-free network?
One important feature of our approach is that the rigorous checking of a series of benchmark
networks allows assessing the significance of any similarities to other network types found. In the
study of a much simpler brain network, it has previously been established that the brain of C.
elegans is small-world, but not scale-free (Amaral et al., 2000). However, we found that effects of
damage on the modelled cat and macaque brain connectivity networks are largely similar to those
observed in scale-free networks. Furthermore, the similarity of scale-free and original cortical
networks, as measured by graph similarity, was higher than for other benchmark networks. This
agrees with other findings: a scale-free network architecture has been found for functional brain
networks in humans (Eguiluz et al., 2005). In addition, the human resting state network of 90 cortical
and subcortical regions showed similar behaviour after the removal of nodes than our structural
network (Achard et al., 2006). This could now be explained by the underlying structural
connectivity.
We note that this issue remains controversial. A study of the human resting state network between
cortical areas (Achard et al., 2006), concluded that the resting state network is not a scale-free
network as (a) it is more resilient towards targeted attack compared to a scale-free benchmark
network, (b) the degree distribution is not a power-law, and (c) late developing areas such as the
dorsolateral prefrontal cortex are among the hubs of the network. The structural network that we
analyzed, however, differed from the resting state functional network. First, the resilience towards
targeted attack was comparable with that of a scale-free network. Second, though the degree does
not follow a power-law distribution this might be due to the small size of the network and
incomplete sampling of connections between regions.
A design for robustness or by-product of functional constraints?
Is the brain optimized for high robustness or is robustness a by-product of other constraints? In our
view, the emergence of highly-connected areas is more likely to be a side effect of brain evolution
and development generating structures for efficient processing. For example, highly-connected areas
(hubs) in the brain could play a functional role as integrators or spreaders of information (Sporns and
Zwi, 2004).
What could be developmental reasons for some regions having a higher connectivity than others?
There are several potential developmental mechanisms for yielding brain networks with the highly-
connected nodes. Work in brain evolution suggests that when new functional structures are formed
by specialization of parts of phylogenetically older structures, the new structures largely inherit the
connectivity pattern of the parent structure (e.g., Preuss, 2000). This means that the patterns are
repeated and small modifications are added during the evolutionary steps that can arise by
duplication of existing areas (Krubitzer and Kahn, 2003). Such inheritance of connectivity by
copying of modules is proposed to lead to scale-free metabolic systems (Ravasz et al., 2002). A
developmental mechanism for varying the edge degree of regions could be the width of the
developmental time window for synaptogenesis at different regions (Kaiser and Hilgetag, 2007).
In conclusion, we have introduced a quantitative method to characterize the robustness of brain
networks and compare it to that of standard network types. We have shown that cortical networks
are affected in ways similar to scale-free networks concerning the elimination of nodes or
connections. However, a direct comparison of degree distributions has been impossible. Our analysis
can be extended to employ more elimination strategies or use different properties to characterize the
damaged networks. In the future, it would be interesting to compare the effect of experimental
lesions with the simulated lesions of our approach. We therefore hope that this theoretical approach
will prove useful in modelling robustness towards lesions.
Acknowledgements
Supported by the Wellcome Trust, EU Framework Five (R.M) as well as German National Merit
Foundation and Fritz-ter-Meer-Foundation (M.K.).
Abbreviations
ASP, Average Shortest Path;
ORI, Original (brain) network;
RND, (Erdös-Renyi) Random network;
SF, Scale-free network;
SW, Small-world network;
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Tables
Table 1. Comparison of brain networks and benchmark networks.
The table shows the average shortest path and the clustering coefficient statistics for the macaque
and cat brain structure networks, and for the respective benchmark random, rewired, small-world,
and scale-free networks. For the benchmark networks, the data shows the mean value and the
standard deviation of 50 generated networks.
Average shortest path Clustering coefficient
Macaque 2.414 0.453
Random mean 2.093 ± 0.009 0.142 ± 0.004
Rewired mean 2.118 ± 0.010 0.239 ± 0.009
Small-world mean 2.439 ± 0.054 0.416 ± 0.022
Scale-free mean 2.078 ± 0.042 0.564 ± 0.042
Cat 1.961 0.542
Random mean 1.749 ± 0.002 0.265 ± 0.003
Rewired mean 1.803 ± 0.006 0.381 ± 0.006
Small-world mean 1.868 ± 0.017 0.461 ± 0.016
Scale-free mean 1.768 ± 0.014 0.535 ± 0.029
Table 2. Overview of the most highly-connected regions in the cat and macaque network.
The table shows the total number of connections of the region (degree) as well as the number of
incoming / afferent (in-degree) and outgoing / efferent (out-degree) connections. The maximal
possible number of connections would have been 110 connections for the cat and 130 connections
for the macaque.
Rank Area Total Incoming Outgoing
1 AES 59 30 29
2 Ia 55 29 26
3 7 54 28 26
4 Ig 52 22 30
5 5al 49 30 19
Macaque
Rank Area Total Incoming Outgoing
1 A7B 43 23 20
2 LIP 42 19 23
3 A46 42 23 19
4 FEF 38 19 19
5 TPT 37 18 19
Figures
Figure 1. Examples of random and scale-free networks. Schematic view of network connectivity
features. (A) Simple scale-free network having highly-connected nodes (hubs) here shown at the
centre. (B) Simple random network; both networks have the same number of nodes and edges.
0 5 10 15 20 25 30 35 40 45
degree
s Macaque
Random
0 5 10 15 20 25 30 35 40 45 50 55
degree
s Cat
Random
Figure 2. Direct comparison of degree distribution. (A) Histogram of the degree distribution of
the macaque (gray) compared to the distribution of random networks (binomial distribution given
the probability p=0.1417 that an edge occurs, black). (B) Histogram of the degree distribution of the
cat (gray) compared to the distribution of random networks (binomial distribution given the
probability p=0.2643 that an edge occurs, black).
rewired scale-free random small-world
Macaque
Figure 3. Similarity of network connectivity. For each type of benchmark network, 1,000
networks were generated. As the cat network has a larger number of edges, the percentages of
similar edges are also higher. The similarity with the cortical networks is as good for the scale-free
networks as for the rewired cortical networks. In contrast, the similarity of random and small-world
networks is significantly lower.
Figure 4. Sequential node eliminations in Macaque cortical networks. The fraction of deleted
nodes (zero for the intact network) is plotted against the average shortest path (ASP) after node
removals. Nodes were removed in order of connectivity, starting with the most highly connected
nodes (targeted elimination) or the node order was determined randomly (random elimination). (A)
Cortical network during targeted (dashed) and random (solid line) elimination. In the subsequent
plots B, C and D, the dashed line shows the average effect of targeted elimination and the thin
dashed lines the 95% confidence interval for the generated networks. The solid line represents the
average effect of random elimination. The dashed grey line represents targeted removal in the
cortical network of A for comparison. (B) Small-world benchmark network. (C) Scale-free
benchmark network. (D) Random benchmark network. (The complete set of figures for cat and
macaque with node and edge elimination and the effect on ASP is available in the supplementary
material).
Figure 5. Sequential node eliminations in cat cortical networks. The fraction of deleted nodes
(zero for the intact network) is plotted against the average shortest path (ASP) after node removals.
Nodes were removed in order of connectivity, starting with the most highly connected nodes
(targeted elimination) or the node order was determined randomly (random elimination). (A)
Cortical network during targeted (dashed) and random (solid line) elimination. Lines in B-C have the
same meaning as in Fig. 4. (B) Small-world benchmark network. (C) Scale-free benchmark network.
(D) Random benchmark network.
A
cortical scale-free rewired random small-world
Macaque
cortical scale-free rewired random small-world
Macaque
Figure 6. Fraction and value of peak ASP for targeted node elimination. The average values and
standard deviations are shown for the 50 generated benchmark networks. (A) Fraction of eliminated
nodes, at which the largest ASP was attained. For the cat cortical network, only the fraction of peak
ASP for the scale-free network is close to the cat network whereas the fractions of other benchmark
networks are higher. The same is the case for the macaque cortical network. (B) Peak value of the
ASP. It is higher for scale-free networks than for cortical networks, in contrast to more similar
values for the other benchmark networks.
A
cortical scale-free rewired random small-world
Macaque
cortical scale-free rewired random small-world
Macaque
Figure 7. Fraction and value of peak ASP for targeted connection elimination. The average
values and standard deviations are shown for the 50 generated benchmark networks. (A) Fraction of
eliminated connections, at which the largest ASP was attained. For the cat network, scale-free and
small-world fractions are similar to the cortical value whereas fractions of rewired and random
networks are significantly higher. For the macaque network, however, all benchmark networks
except for the small-world network show a similar fraction of peak ASP. (B) Peak values of the
ASP. The peak value of the cat cortical network can be matched by the random and rewired
networks, nearly by the scale-free but significantly not by the small-world network. For the
macaque, all networks except for the scale-free network show significantly different values.
|
0704.0393 | A thermodynamic model for the melting of supported metal nanoparticles | A thermodynamic model for the melting of supported metal
nanoparticles
S. C. Hendy∗
Industrial Research Ltd, Lower Hutt, New Zealand and
MacDiarmid Institute for Advanced Materials and Nanotechnology,
School of Chemical and Physical Sciences,
Victoria University of Wellington, New Zealand
(Dated: November 4, 2018)
Abstract
We construct a simple thermodynamic model to describe the melting of a supported metal
nanoparticle with a spherically curved free surface both with and without surface melting. We
use the model to investigate the results of recent molecular dynamics simulations, which suggest
the melting temperature of a supported metal particle is the same as that of a free spherical
particle with the same surface curvature. Our model shows that this is only the case when the
contact angles of the supported solid and liquid particles are similar. This is also the case for the
temperature at which surface melting begins.
http://arxiv.org/abs/0704.0393v1
INTRODUCTION
Despite decades of study, the melting of nanoparticles continues to generate interest
[1, 2, 3, 4]. In general, the melting temperature of spherical nanoparticles has been found
to decrease in proportion to the surface area to volume ratio of the particle [5], as the
surface free energy of a molten droplet is less than that of the corresponding solid particle.
Although free nanoparticle calorimetry has advanced considerably in recent years [6, 7],
most experimental determinations of the melting points of nanoparticles are conducted with
supported particles (gold [8], lead [9] and tin [10], for example). The melting of supported
nanoparticles is also important in carbon nanotube growth and other catalytic processes
[11, 12], and for the stability of devices assembled from nanoparticles [13, 14]. Thus it is of
interest to study the effect of the substrate on the melting point of a supported nanoparticle.
Recent molecular dynamics simulations [15] of supported iron nanoparticles with a
strongly interacting substrate found that the melting point of the particles was reduced
in inverse proportion to the equilibrium surface curvature that results as they relax to wet
the substrate. This statement also holds in the free particle limit since the curvature of a
free spherical particle of radius a is proportional to its surface to volume ratio, 3/a. Inter-
estingly, the simulations in Ref. [15] found that the constant of proportionality between the
decrease in melting point and the surface curvature did not depend on whether the particle
was supported or free. In other words, the melting temperature of a supported particle that
has a free surface with radius of curvature a, was found to be the same as that of a free
spherical particle with the same surface curvature. The simulated nanoparticles in Ref. [15]
also exhibited surface melting prior to complete melting. Surface melting is phenomena
thought to occur both on certain bulk metal surfaces [16] and in certain metal nanoparticles
[17, 18].
In this paper we use a simple thermodynamic model to investigate the role of the sub-
strate in both melting and surface melting of metal nanoparticles. Our model suggests that
the result in Ref. [15], that the relative decrease in melting point is proportional to the solid
particle surface curvature, only holds when the contact angles of the supported solid and
liquid particles with the substrate are similar. We also show that supported clusters will
exhibit surface melting under certain circumstances, and that the surface melting tempera-
ture in free and supported particles in clusters with same surface curvature is the same only
when the contact angles of the supported solid and liquid phases coincide.
GEOMETRY OF SUPPORTED PARTICLES
We start by considering a solid nanoparticle, initially spherical with radius a, that is
placed on a flat substrate. We neglect the effect of faceting, curvature dependent surface
energies and internal strains due to epitaxial mismatch with the substrate. Furthermore, we
will assume that the particle has relaxed to its equilibrium geometry i.e. that the nanopar-
ticle has relaxed to ”wet” the substrate. Provided the nanoparticle is heated sufficiently
slowly, the particle should relax to this geometry prior to melting. With these simplifying
assumptions, the geometry of the relaxed particle can be approximated by a spherical cap,
as shown in Fig. 1, with dimensions parameterised by either the cap height H , or radius of
curvature of the free surface R, which minimizes the surface energy of the nanoparticle and
substrate.
The surface energy Γ of the system can be written as:
Γ = 2πRHγs + πH(2R−H)(γsb − γb) + Γb (1)
where γs is the surface energy density of the free particle surface, γb is the surface energy
density of the substrate, γsb is the particle-substrate interfacial energy density and Γb is the
total energy of the bare substrate. We will assume that the density of the particle ρs does
not depend on the contact angle so that the volume of the supported particle remains the
same as that of the free particle.
Writing the volume of the particle as a function of H and R, it is then straightforward
to show that Γ is minimized if H = − (∆γsb/γs)R where ∆γsb = γb − γs − γsb. We note
that ∆γsb is often called the spreading parameter in the context of wetting phenomena: if
∆γsb > 0 then the particle will relax to fully wet the substrate. Here we are interested in the
case where the particle does not fully wet the substrate (contact angles greater than zero)
i.e. when ∆γsb < 0 and H/R = −∆γsb/γs > 0 at equilibrium. In fact this minimum value
of Γ can be written as
Γ∗ = 2πγs
a2 + Γb, (2)
where R∗s is the corresponding radius of curvature of the supported solid nanoparticle, given
R∗s =
(3 + ∆γsb/γs)
. (3)
Thus Γ∗ and R∗s are the equilibrium surface energy and radius of curvature of the particle
respectively. Note that the contact angle of the particle can range from 0 to 180 degrees
depending on the value of the spreading parameter ∆γsb.
MELTING AND SURFACE MELTING
In what follows we will assume that the density of the solid and liquid phases are identical
i.e. ρs = ρl = ρ. We first consider the situation in which there is no surface melting. In this
case, melting will occur at a temperature when the free energy of the solid particle wetting
the substrate is equal to that of the corresponding liquid droplet wetting the substrate.
If γl is the surface energy density of the free liquid droplet and R
l is the corresponding
equilibrium radius of curvature, then the difference in free energy between the solid and
liquid is
fs − fl + 3
where fs (fl) is the bulk free energy density of the solid (liquid). Now, using fl − fs =
ρL (1− T/Tc), where L is the latent heat of melting and Tc is the bulk melting temperature,
we find the melting temperature Tm of the supported particle is given by:
Tm = Tc
= T freem (R
ρR∗sL
Tc (5)
Thus if R∗s = R
l = R
∗ then we recover the result of Ref. [15], namely that Tm = T
In other words, if the contact angles of the solid and liquid droplets are equal, the melting
temperature of the supported particle is the same as that of a free particle with an identical
surface curvature, a = R∗. However, if the curvature of the supported liquid particle is dif-
ferent from that of the supported solid particle, it can be seen that the melting temperature
will deviate from that found in Ref. [15].
Now we consider surface melting as illustrated in figure 1 which is thought to occur in
many metals prior to melting [16]. We are interested in determining the onset of melting,
when the solid particle is wet by a thin layer of melt (thickness δ) at the solid-vapor interface.
We will assume that this melt forms a layer of uniform thickness with a geometry like that
represented in figure 1 with δ = R− r = H − h. The total free energy of the surface melted
particle is then a function of δ: F (δ) = Vs(δ)(fs − fl) + V fl +Γ(δ) where Vs(l) is the volume
of the solid (liquid) and Γ is the thickness dependent surface energy. In particular
Γ = π (2RHγl + r(2r − h)γsb + δ(2R− δ)γlb + 2rhγsl(δ))
where γsl(δ) = γsl+∆γsl exp (−δ/ξ) and ξ is a correlation length that describes the thickness
dependence of the interfacial energy in thin metallic liquid films [16] (in Pb, for example, ξ
has been measured to be ∼ 0.6 nm [19]). As the surface melting proceeds, the curvature of
the particle will relax to minimize the free energy i.e R∗ = R∗(δ) where R∗ minimizes the
free energy F for a given δ.
In an isolated spherical nanoparticle of radius a, by minimizing the free energy F (δ) with
respect to δ and setting δ = 0, one can show that surface melting begins at a temperature,
Ts(a) given by
T frees (a) = Tc
(γs − γl)
. (6)
provided ∆γsl > 0 and a > ξ(γs − γl)/∆γsl (if a is less than this, full melting will precede
surface melting i.e. T frees > T
m [17], and equation (5) will hold).
For surface melting to occur in a supported solid nanoparticle with equilibrium curvature
Rs, a minimum in the free energy F (δ) must appear at δ = 0. It is straightforward to show
that a minimum in F at δ = 0 occurs at the temperature Ts:
Ts (Rs) = T
s (Rs) +
γs∆γlb − γl∆γsb
= T frees (Rs) +
cos θs − cos θl
1− cos θs
where θs and θl are the contact angles for solid particle and liquid particle respectively
(defined via Young’s relation γs(l) cos θs(l) = γb − γs(l)b). Once again, if the contact angles of
the solid and liquid droplets are equal, then the temperature at which surface melting occurs
is identical to that of a free particle with the same surface curvature, Rs i.e. Ts = T
s (Rs).
Further, if cos θs > cos θl, so that the substrate favors contact with the solid over that with
the liquid, the corresponding Ts increases and vice versa.
Complete melting will occur once the free energy of the surface melted particle, F (δ),
equals that of the corresponding liquid droplet, Fl i.e. at the temperature Tm and liquid film
thickness δm which satisfy F (δm) = Fl. It is not possible to obtain an analytic expression
for δm or Tm, but numerical solutions to the resulting equations are shown in figure 2 as
a function of Rs for Pb particles. The figures clearly show the strong dependence of the
melting temperature on the liquid droplet contact angle: a difference of ∼ 10o in the molten
particle contact angle can shift the melting point by ∼ 100 K. Note that the melting point
of a free particle with radius Rs no longer coincides with that of a supported particle with
radius of curvature Rs when cos θs = cos θl, as in general the radius of curvature of the
critical surface melted droplet will not be that of the solid particle (although the curves lie
close to each other).
CONCLUSION
We conclude that the melting temperature (and surface melting temperature, if the par-
ticle exhibits surface melting) of supported nanoparticles depends on the radius of curvature
(or the contact angle) of both the supported solid and liquid particles. In general, we do
not expect these curvatures to be the same: on a non-ideal solid substrate for example,
epitaxial effects may favor one phase over the other. It is likely that the ideal substrate used
in Ref [15] resulted in very similar solid and liquid particle contact angles. We have shown
that is only in this ”ideal” case that the melting temperature of free and supported particles
with the same curvature is coincident, whether they exhibit surface melting or otherwise.
Thus, results from free particle melting, where the curvature of the solid and liquid particles
remain substantially the same, have only limited applicability to supported particle melting.
∗ Electronic address: [email protected]
[1] H. Haberland, in Atomic Clusters and Nanoparticles: Les Houches Session LXXIII (Springer,
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1135-1143 (1995).
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(1996).
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(2000).
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C. Seigert and R. Blaikie, Nanotechnology 15, 1382 (2004).
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D. M. A. Mackensie, A. Ayesh, K. C. Tee, A. Awasthi and S. C. Hendy, Appl. Phys. Lett. 89,
213105 (2006).
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[16] J. F. van der Veen, B. Pluis and A. W. Denier van der Gon, in Kinetics of Ordering and
Growth at Surfaces 343-354, (Plenum Press, New York, 1990).
[17] T. Bachels, H.-J. Güntherodt and R. Schäfer, Phys. Rev. Lett. 85, 1252 (2000).
[18] U. Tartaglino, T. Zykova-Timan, F. Ercolessi and E. Tosatti, Phys. Rep. 411, 291321 (2005).
[19] B. Pluis, T. N. Taylor, D. Frenkel and J. F. van der Veen, Phys. Rev. B 40, 1353 (1989).
FIG. 1: The model for the geometry of a supported nanoparticle in equilibrium. We assume that
the particle is a spherical cap of height H and radius of curvature R (left - the dashed lines simply
illustrate the radius of curvature). At the onset of surface melting, we assume that the geometry is
close to that of the solid particle in its equilibrium geometry and that that the solid particle (radius
of curvature r and height h) is initially wet by a molten layer of uniform thickness δ = R−r = H−h
(right).
2 4 6 8 10 12 14
γlb=0.05
γlb=0.15
γlb=0.10
Rs (nm)
2 4 6 8 10 12 14
cos θs=cos θl
γlb=0.05
γlb=0.15
γlb=0.10
Rs (nm)
2 4 6 8 10 12 14
γlb=0.05
γlb=0.15
γlb=0.10
Rs (nm)
2 4 6 8 10 12 14
cosθs=cosθl
γlb=0.05
γlb=0.15
γlb=0.10
FIG. 2: The melting temperature Tm and critical liquid film thickness δm for supported Pb clusters
as a function of the radius of curvature Rs of the relaxed solid particle for γlb = 0.05, 0.10, 0.15 J
m−2 and in the case where cos θs = cos θl (γlb ≃ 0.13 J m
−2). Also shown is the melting temperature
of a free particle with radius Rs. Other surface energies used are γsv = 0.61, γlv = 0.48, γsl = 0.05,
γb = 0.25 and γsb = 0.1 J m
−2 giving a contact angle of 75.8o for the solid supported cluster, and
contact angles for the liquid droplets ranging from 78.0o to 65.4o respectively. Other parameters
used were ξ = 0.63 nm, ρ = 10950 kg m−3, L = 22930 J kg−1 and Tc = 600.65 K [9].
Introduction
Geometry of supported particles
Melting and surface melting
Conclusion
References
|
0704.0394 | Average optimality for risk-sensitive control with general state space | Average optimality for risk-sensitive control with general state space
The Annals of Applied Probability
2007, Vol. 17, No. 2, 654–675
DOI: 10.1214/105051606000000790
c© Institute of Mathematical Statistics, 2007
AVERAGE OPTIMALITY FOR RISK-SENSITIVE CONTROL WITH
GENERAL STATE SPACE1
By Anna Jaśkiewicz
Wroc law University of Technology
This paper deals with discrete-time Markov control processes on
a general state space. A long-run risk-sensitive average cost crite-
rion is used as a performance measure. The one-step cost function is
nonnegative and possibly unbounded. Using the vanishing discount
factor approach, the optimality inequality and an optimal stationary
strategy for the decision maker are established.
1. Introduction and the model. This paper deals with discrete-time
Markov control processes on a general state space. The one-step cost function
is nonnegative and possibly unbounded. The decision maker is supposed to
be risk-averse with a constant risk coefficient γ > 0. The risk-sensitive aver-
age cost criterion is used as a performance measure. The aim of the work is to
establish the optimality inequality for risk-sensitive dynamic programming
and derive an optimal stationary policy. The result is proved under two
different sets of compactness-continuity assumptions, namely, for Markov
control processes with weakly continuous transition probabilities [Condition
(W)], as well as transition probabilities that are continuous with respect
to setwise convergence [Condition (S)]. A similar problem for risk-neutral
stochastic control models has been examined in [27] using the vanishing dis-
count factor approach. However, it is well known that, for risk-sensitive con-
trol models, an analogous approximation of the average cost via a sequence
of the corresponding discounted models does not work. Instead of this, fol-
lowing [9, 15, 16], we introduce an auxiliary discounted minimax problem.
A variational formula that expresses the mutual relationship between the
relative entropy function and the logarithmic moment-generating function
enables us to connect the discounted minimax model with the original one.
Received March 2006; revised September 2006.
1Supported by MEiN Grant 1 P03A 01030.
AMS 2000 subject classifications. Primary 60J05, 90C39; secondary 60A10.
Key words and phrases. Risk-sensitive control, Borel state space, average cost optimal-
ity inequality.
This is an electronic reprint of the original article published by the
Institute of Mathematical Statistics in The Annals of Applied Probability,
2007, Vol. 17, No. 2, 654–675. This reprint differs from the original in pagination
and typographic detail.
http://arxiv.org/abs/0704.0394v1
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http://dx.doi.org/10.1214/105051606000000790
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http://dx.doi.org/10.1214/105051606000000790
2 A. JAŚKIEWICZ
Next, assuming that a certain family of functions is bounded [Condition (B)]
and using Fatou’s lemma (for weakly or setwise convergent measures), we
obtain the optimality inequality.
The predecessor of our result is Theorem 4.1 in [16], where the optimality
inequality for the risk-sensitive dynamic programming with a countable state
space was established. Instead of boundedness assumption (B), Hernández-
Hernández and Marcus [16] assume that there exists a stationary policy
which induces a finite average cost that is equal some constant in each
state. On the other hand, it is well known that an optimal risk-sensitive
average cost may depend on the initial state (see Example 1). This behavior
happens if the risk factor is too large. Instead of this restriction on the
risk coefficient, we use Condition (B), which makes the process reach “good
states” sufficiently fast.
There is a rich literature in risk-sensitive control, going back at least to
the seminal works of Howard and Matheson [18] and Jacobson [19], which
covered the finite horizon case. The average cost criterion on the infinite
horizon was studied in [5, 8, 14, 15, 16, 31] for a denumerable state space
and in [10, 11, 20] for a general state space. It is also worth mentioning
that risk-sensitive control finds natural applications in portfolio managment,
where the objective is to maximize the growth rate of the expected utility
of wealth; see [3, 4, 30] and the references cited therein.
The paper is organized as follows. Below a Markov control model with
the long-run average cost criterion as a performance measure is described, as
well as some basic notation is set up. In Section 2 we introduce preliminaries
and present the auxiliary discounted minimax problem, which is, in turn,
solved in Section 3. The main result is established in Section 4. Section 5
contains a discussion of Condition (B), and in the Appendix a variational
formula for the logarithmic moment-generating function is stated.
A discrete-time Markov control process is specified by the following ob-
jects:
(i) The state space X is a standard Borel space (i.e., a nonempty Borel
subset of some Polish space).
(ii) A is a Borel action space.
(iii) K is a nonempty Borel subset of X×A. We assume that, for each
x ∈X , the nonempty x-section
A(x) = {a ∈A : (x,a) ∈K}
of K is compact and represents the set of actions available in state x.
(iv) q is a regular conditional distribution from K to X.
(v) The one-step cost function c is a Borel measurable mapping from K
to [0,+∞].
RISK-SENSITIVE CONTROL 3
Then the history spaces are defined as H0 = X, Hk = (X ×A)
k ×X and
H∞ = (X ×A)
∞. As usual, a policy π = {πk, k = 0,1, . . .} ∈ Π is a sequence
of transition probabilities from Hk to A such that πk(A(xk)|hk) = 1, where
hk = (x0, a0, . . . , xk) ∈Hk. The class of stationary policies is identified with
the class F of measurable functions f from X to A such that f(x) ∈A(x). It
is well known that F is nonempty [6]. By the Ionescu–Tulcea theorem [24],
for each policy π and each initial state x0 = x, a probability measure P
and a stochastic process {(xk, ak)} are defined on H∞ in a canonical way,
where xk and ak describe the state and the decision at stage k, respectively.
By Eπx we denote the expectation operator with respect to the probability
measure Pπx .
Let γ > 0 be a given risk factor. For any initial state x ∈X and policy
π ∈ Π, we define the following risk-sensitive average cost criterion:
J(x,π) = lim sup
logEπx exp
c(xk, ak)
Our aim is to minimize J(x,π) within the class of all policies and find a
policy π∗, for which
J∗(x) := inf
J(x,π) = J(x,π∗).
Throughout the paper the following assumption will be supposed to hold
true even without explicit reference:
∃π̃ ∈ Π J(x, π̃) < +∞.(G)
Remark 1. Throughout the remainder, we assume that the risk factor
γ > 0 is arbitrary and fixed. Therefore, here and subsequently, we shall not
indicate that some quantities depend on γ [e.g., we write J(x,π) instead of
Jγ(x,π), dropping the index γ].
2. Preliminaries. Let Pr(X) be the set of all probability measures on
X. Fix ν ∈ Pr(X). The relative entropy function R(·‖ν) is a mapping from
Pr(X) into R defined as follows:
R(µ‖ν) :=
dµ, µ≪ ν,
+∞, otherwise.
It is well known that R(µ‖ν) is nonnegative for any µ ∈ Pr(X) and R(µ‖ν) =
0 if and only if µ = ν (consult Lemma 1.4.1 in [12]).
We shall consider the following auxiliary minimax problem, associated
with our original Markov control process. The set X is the state space,
4 A. JAŚKIEWICZ
while A and Pr(X) are the action sets for the decision maker and op-
ponent, respectively. The process then operates as follows. In a state xn,
n = 0,1, . . . , the controller chooses an action an ∈ A(xn), while the oppo-
nent selects µn(·)[xn, an] ∈ Pr(X). As a consequence, the controller pays
γc(xn, an)−R(µn‖q(·|xn, an)) to his opponent, and the system moves to the
next state according to the probability distribution µn(·)[xn, an].
We shall deal with the following classes of strategies. It will cause no
confusion if we continue to use the same letters to denote strategies for
the controller. Namely, π stands for a randomized control strategy (policy),
whereas f denotes a stationary strategy. We write Π and F to denote the sets
of corresponding strategies. For the opponent’s class of strategies, we confine
to the stationary one, which is identified with the class P of stochastic kernels
p on X given K.
Let (Ω,F) be the measurable space consisting of the sample space Ω =
(X × A)∞ and its product σ-algebra F . Then for an initial state x ∈ X,
and strategies π and p, there exists a unique probability measure Pπpx and,
again, a stochastic process {(xk, ak)} is defined (Ω,F) in a canonical way,
where xk denotes the state at time k and ak is the action for the controller.
With some abuse of notation, we let hk stand for the history of the process
up to the kth state, that is,
hk = (x0, a0, x1, . . . , ak−1, xk).
The corresponding expectation operator is denoted by Eπpx .
For fixed x ∈ X, π ∈ Π and p ∈ P , we define the following functional
costs:
Vβ(x,π, p) =
βkEπpx [γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))],(1)
where β ∈ (0,1) is the discount factor, and
j(x,π, p) = lim sup
Eπpx [γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))].
Note that, since the function R(·‖·) is lower semicontinuous on Pr(X) ×
Pr(X) and p and q are stochastic kernels [i.e., measurable functions of (x,a)],
it follows that the mapping
(x,a) 7→R(p(·|x,a)‖q(·|x,a))
is measurable (Lemma 1.4.3(f) in [12]). Observe that Vβ(x,π, p) and j(x,π, p)
might be undetermined, because c can be unbounded. We thus restrict the
set of admissible strategies for the opponent in the following way.
RISK-SENSITIVE CONTROL 5
Definition 1. Given π = {πk} ∈ Π, we say that p ∈ P is a π-admissible
strategy iff
A(xk)
R(p(·|xk, a)‖q(·|xk, a))πk(da|hk) < +∞,(2)
and moreover, there exists a constant C ≥ 0, possibly depending on π and
p, such that
A(xk)
[γc(xk, a) −R(p(·|xk, a)‖q(·|xk, a))]πk(da|hk) + C ≥ 0,
for all histories of the process hk, k ≥ 0, induced by p and π. We denote
this set by Q(π). [Note that this set is nonempty, since p = q ∈Q(π) for any
π ∈ Π.]
Let us introduce the following notation. For any π ∈ Π, p ∈ Q(π) and
n≥ 1, define
Jn(x,π) = logE
x exp
c(xk, ak)
jn(x,π, p) =
Eπpx [γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))].
Now we are ready to present the result that was originally proved in
[16] for Markov strategies. However, it still remains valid when arbitrary
strategies for the decision maker are considered. Therefore, for the sake of
clarity, we state the result with its proof.
Proposition 1. Let x∈X and p ∈Q(π). Then:
(a) supp∈Q(π) jn(x,π, p) ≤ Jn(x,π) for each n≥ 1,
(b) lim supn→∞ supp∈Q(π)
jn(x,π, p) ≤ γJ(x,π).
Proof. (a) Let p ∈ Q(π) be any stochastic kernel. For n = 1, we con-
clude
j1(x,π, p) ≤ E
x (γc(x,a0)) ≤ logE
γc(x,a0) = J1(x,π),
where the first inequality holds since the relative entropy is nonnegative, and
the second one is due to Jensen’s inequality. Now assume that the hypothesis
is true for some n≥ 1. Clearly,
jn+1(x,π, p) =
Eπpx [γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))]
= Eπpx
[γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))], n≥ 1.
6 A. JAŚKIEWICZ
Denote by π(1) the “1-shifted” strategy, that is,
(·|hk) = πk+1(·|x0, a0, hk), k ≥ 0.
Then, we have
jn+1(x,π, p)
= Eπpx [γc(x,a0) + jn(x1, π
(1), p) −R(p(·|x,a0)‖q(·|x,a0))]
≤ Eπpx (γc(x,a0))
+ Eπpx (E
x {[Jn(x1, π
(1)) −R(p(·|x,a0)‖q(·|x,a0))]|a0})
= Eπx log e
γc(x,a0)
+ Eπpx
Jn(x1, π
(1))p(dx1|x,a0) −R(p(·|x,a0)‖q(·|x,a0))
log eγc(x,a0)π0(da0|x)
eJn(x1,π
(1))q(dx1|x,a0)π0(da0|x)
eγc(x,a0)+
γc(xk,ak)q(dx1|x,a0)π0(da0|x)
≤ log
eγc(x,a0)+
γc(xk,ak)q(dx1|x,a0)π0(da0|x)
= Jn+1(x,π).
Clearly, the first inequality follows from the induction hypothesis. The third
inequality is due to Jensen’s inequality, whilst the second one follows from
Lemma A in the Appendix. Since p ∈Q(π) is arbitrary, we get the desired
conclusion.
Part (b) follows directly from part (a). �
Remark 2. Note that in the proof of Proposition 1 we did not really
have to use the fact that p ∈ Q(π). The only assumption which plays an
essential role is condition (2). Namely, it guarantees that jn(x,π, p) is well
defined for all n≥ 1, x ∈X and π ∈ Π. However, in Definition 1 we restrict
the opponent’s class of strategies to the set Q(π) in order to be able to apply
the Hardy–Littlewood theorem. In actual fact, later on it will be clear that
the set Q(π), where π ∈ Π, is sufficiently large. Namely, the supremum of
certain discounted functional costs over the set Q(π) will not change if we
add new elements to Q(π); see the proofs of Lemmas 1 and 2.
RISK-SENSITIVE CONTROL 7
Let π̃ be as in assumption (G) and let p ∈Q(π̃). Then from the Hardy–
Littlewood theorem (Theorem H.2 in [13]), we get
lim sup
(1 − β)Vβ(x, π̃, p) ≤ lim sup
jn(x, π̃, p)
and from Proposition 1(b),
lim sup
p∈Q(π̃)
jn(x, π̃, p) ≤ γJ(x, π̃).
Combining these two inequalities, we conclude that
lim sup
(1− β)Vβ(x, π̃, p) ≤ γJ(x, π̃) for every p ∈Q(π̃).
This in turn yields
lim sup
(1− β)Vβ(x) ≤ γJ(x, π̃),(4)
where Vβ(x) is the upper value of functional cost (1), that is,
Vβ(x) = inf
p∈Q(π)
Vβ(x,π, p).
Consequently, inequality (4) and assumption (G) together lead to the fol-
lowing:
Vβ(x) < +∞(5)
for each x ∈X and β ∈ (0,1). In addition, Vβ(x) ≥ 0. Now defining
ρ := inf
J(x,π), mβ := inf
Vβ(x)
and observing that
lim sup
(1 − β)mβ ≤ γρ,(6)
one can deduce that there exists a sequence of discount factors {βn} con-
verging to 1 for which
(1− βn)mβn = l,(7)
where l is a certain nonnegative constant.
8 A. JAŚKIEWICZ
3. A solution to the auxiliary discounted minimax problem. The main
thrust of this section is to solve the auxiliary discounted minimax problem
introduced in the previous section. In other words, we look for a discounted
functional equation whose solution is the function Vβ . This is done by an ap-
proximation of the above-mentioned minimax models by ones with bounded
cost functions. These models in turn are solved by a fixed point argument in
Proposition 1. Next, we show in Lemma 1 that the corresponding solutions
equal the upper values of some discounted costs on the infinite horizon. Fi-
nally, the limit passage in Lemma 2 gives the desired discounted functional
equation with the function Vβ as a solution.
We shall need the following two sets of compactness-semicontinuity as-
sumptions, which will be used alternatively.
Condition (S).
(i) The set A(x) is compact.
(ii) For each x ∈X and every Borel set D ⊂X, the function q(D|x, ·) is
continuous on A(x).
(iii) The cost function c(x, ·) is lower semicontinuous for each x ∈X.
Condition (W).
(i) The set A(x) is compact and the set-valued mapping x 7→ A(x) is
upper semicontinuous, that is, {x ∈X : A(x) ∩ B 6= ∅} is closed for every
closed set B in A.
(ii) The transition law q is weakly continuous on K, that is, the function
(x,a) 7→
u(y)q(dy|x,a), (x,a) ∈K,
is continuous function for each bounded continuous function u.
(iii) The cost function c is lower semicontinuous on K.
By Lb(X) and Bb(X), we denote the set of all bounded lower semicontin-
uous and bounded Borel measurable functions on X, respectively. Further,
let N stand for the set of positive integers. Choose N ∈ N and define the
truncated cost function
cN (x,a) = min{N,c(x,a)}.
The following result was proved under Condition (W) for bounded cost
functions by a fixed point argument; see page 72 in [10]. However, a simple
and obvious modification of the proof gives the conclusion under Condition
(S) as well.
RISK-SENSITIVE CONTROL 9
Proposition 2. Under (W) [(S)], for any discount factor β ∈ (0,1) and
a number N ∈N, there exists a unique function wNβ ∈ Lb(X) [w
β ∈Bb(X)]
such that
= min
a∈A(x)
N (x,a)
q(dy|x,a)
for each x ∈X, and
0 ≤ (1 − β)wNβ (x) ≤Nγ.(9)
Moreover, there exists a stationary strategy f0 ∈ F (possibly depending on β
and N) that attains the minimum in (8).
Let β and N be fixed just in the next lemma.
Lemma 1. Assume (W) or (S). Then, it holds
wNβ (x)
= inf
p∈Q(π)
Eπpx β
k[γcN (xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))]
for any initial state x ∈X.
Proof. Note that (8) can be rewritten in the following equivalent form:
wNβ (x) = min
a∈A(x)
γcN (x,a) + log
q(dy|x,a)
.(11)
Applying Lemma A in the Appendix to (11), we get
wNβ (x)
= min
a∈A(x)
µ∈∆(x,a)
γcN (x,a) −R(µ‖q(·|x,a)) + β
wNβ (y)µ(dy)
∆(x,a) := {µ ∈ Pr(X) :R(µ‖q(·|x,a)) < +∞}, (x,a) ∈K.
Moreover, the measure
µ0(dy)[x,a] =
q(dy|x,a)
q(dy|x,a)
achieves the supremum in (12). Put
p0(dy|x,a) = µ0(dy)[x,a] for each (x,a) ∈K.(13)
10 A. JAŚKIEWICZ
Note that p0 ∈Q(π) for any strategy π ∈ Π. This directly follows from the
definition of R(p0(·|x,a)‖q(·|x,a)) and (9). Simple calculations give the up-
per bound
R(p0(·|x,a)‖q(·|x,a)) ≤ 2
1 − β
for every (x,a) ∈K.
Let p0 be defined as in (13). By (12), we then have
wNβ (x) ≤ γc
N (x,a) −R(p0(·|x,a)‖q(·|x,a)) + β
wNβ (y)p
0(dy|x,a).
By iteration of this inequality n times, it follows
wNβ (x) ≤
βkEπp
x [γc
N (xk, ak) −R(p
0(·|xk, ak)‖q(·|xk, ak))]
+ βn+1Eπp
β (xn+1),
where π is any strategy for the controller. Now, letting n→∞ and making
use of (9), we conclude
wNβ (x) ≤
βkEπp
x [γc
N (xk, ak) −R(p
0(·|xk, ak)‖q(·|xk, ak))].
Since π is arbitrary, we get
wNβ (x) ≤ inf
βkEπp
x [γc
N (xk, ak) −R(p
0(·|xk, ak)‖q(·|xk, ak))]
≤ inf
p∈Q(π)
βkEπpx [γc
N (xk, ak)(14)
−R(p(·|xk, ak)‖q(·|xk, ak))].
Note that inequality (14) is valid because p0 ∈Q(π).
On the other hand, by (12), we can write
wNβ (x) ≥ γc
N (x, f0(x)) −R(p(·|x, f0(x))‖q(·|x, f0(x)))
wNβ (y)p(dy|x, f
0(x)),
with f0 as in Proposition 2 and any p ∈Q(f0). Proceeding along the same
line, we infer
wNβ (x) ≥
x [γc
N (xk, f
0(xk)) −R(p(·|xk, f
0(xk))‖q(·|xk, f
0(xk)))].
RISK-SENSITIVE CONTROL 11
Since p ∈Q(f0) is arbitrary, we easily deduce
wNβ (x) ≥ sup
p∈Q(f0)
x [γc
N (xk, f
0(xk))
−R(p(·|xk, f
0(xk))‖q(·|xk, f
0(xk)))]
≥ inf
p∈Q(π)
βkEπpx [γc
N (xk, ak)
−R(p(·|xk, ak)‖q(·|xk, ak))].
Finally, combining (14) with (15) completes the proof. �
In the remainder of the paper, we shall use the following notation. Let
L(X) denote the set of all lower semicontinuous functions on X, whereas
B(X) stands for the set of all Borel measurable functions on X.
Lemma 2. Let (W) [(S)] hold and β ∈ (0,1). Then, we have the follow-
(a) The function
wβ(x) := lim
wNβ (x)
is finite and nonnegative for each x ∈X. Moreover, wβ ∈L(X) [wβ ∈B(X)].
(b) The functional equation holds
ewβ(x) = min
a∈A(x)
eγc(x,a)
eβwβ(y)q(dy|x,a)
for all x ∈X. Furthermore, there exists a Borel measurable selector fβ ∈ F
of the minima in (16).
(c) For any x ∈X, wβ(x) = Vβ(x).
Proof. Let x ∈X and β ∈ (0,1) be fixed. From (10), it is easily seen
that the sequence {wNβ (x)} is nondecreasing in N. Therefore, wβ(x) =
limN→∞w
β (x) exists and by (9), it is nonnegative. Clearly, under (S),
wβ ∈B(X), whereas, under (W), wβ ∈L(X); see Proposition 10.1 in [26].
In order to prove that wβ(x) is finite for each x ∈X, observe first that,
for any π ∈ Π, p ∈Q(π) and N ∈N,
Vβ(x,π, p) =
βkEπpx [γc(xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))]
βkEπpx [γc
N (xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))].
12 A. JAŚKIEWICZ
Moreover, from Lemma 1, we have
Vβ(x) = inf
p∈Q(π)
Vβ(x,π, p)
≥ inf
p∈Q(π)
βkEπpx [γc
N (xk, ak) −R(p(·|xk, ak)‖q(·|xk, ak))]
= wNβ (x).
Hence, letting N →∞, it follows
Vβ(x) ≥ lim
wNβ (x) = wβ(x).(17)
By (5), Vβ(x) is finite for each x ∈X, so is wβ(x). This finishes the proof of
part (a).
In order to prove part (b), note that by (11) and part (a) the limit
a∈A(x)
γcN (x,a) + log
q(dy|x,a)
exists. Since the first and the second term in (18) are nondecreasing and
(W) or (S) holds, then we may interchange the limit with the minimum
(see Proposition 10.1 in [26]). Furthermore, making use of the Lebesgue
monotone convergence theorem, we conclude (16). The existence of a Borel
measurable selector fβ ∈ F follows from the compactness–semicontinuity
assumptions and Proposition D.5 in [17].
We now turn to proving part (c). Again, taking a logarithm on both sides
of (16), it follows
wβ(x) = min
a∈A(x)
γc(x,a) + log
eβwβ(y)q(dy|x,a)
.(19)
Applying Lemma A in the Appendix to (19), we easily obtain
wβ(x)
= min
a∈A(x)
µ∈∆(x,a)
γc(x,a) −R(µ‖q(·|x,a)) + β
wβ(y)µ(dy)
∆(x,a) = {µ ∈ Pr(X) :R(µ‖q(·|x,a)) < +∞}, (x,a) ∈K.
Observe that by (20), for any p ∈Q(fβ), the following holds:
wβ(x) ≥ γc(x, fβ(x)) −R(p(·|x, fβ(x))‖q(·|x, fβ(x)))
wβ(y)p(dy|x, fβ(x)).
RISK-SENSITIVE CONTROL 13
Iterating this inequality n times, we immediately obtain
wβ(x) ≥
x [γc(xk, fβ(xk))
−R(p(·|xk, fβ(xk))‖q(·|xk, fβ(xk)))]
+ βn+1E
x wβ(xn+1)(21)
x [γc(xk, fβ(xk))
−R(p(·|xk, fβ(xk))‖q(·|xk, fβ(xk)))].
Now note that, by Definition 1,
x [γc(xk, fβ(xk)) −R(p(·|xk, fβ(xk))‖q(·|xk, fβ(xk)))] ≥−C,
for some C ≥ 0 and k ≥ 1. Thus, letting n→∞ in (21), it follows
wβ(x) ≥
x [γc(xk, fβ(xk)) −R(p(·|xk, fβ(xk))‖q(·|xk, fβ(xk)))]
= Vβ(x, fβ, p).
Since p ∈Q(fβ) is arbitrary, we see that
wβ(x) ≥ sup
p∈Q(fβ)
Vβ(x, fβ, p) ≥ Vβ(x).(22)
Inequalities (17) and (22) combined conclude the proof of part (c). �
4. A solution to the risk-sensitive control problem. For any x ∈X and
any discount factor β ∈ (0,1), define
hβ(x) := Vβ(x) −mβ
with mβ = infx∈X Vβ(x). Obviously, hβ is nonnegative.
The following boundedness assumption is supposed to hold true. As men-
tioned in the Introduction, we put off discussing it until Section 5:
Condition (B). For any x ∈X , supβ∈(0,1) hβ(x) < +∞.
Remark 3. A similar assumption and its equivalent variants were used
to study the expected average cost criterion for Markov decision processes
in the risk-neutral setting [17, 27, 28]. Roughly speaking, Hernández-Lerma
and Lasserre [17], Schäl [27], and Sennott [28] assume that the family of the
so-called normalized β-discounted cost functions is bounded. This assump-
tion, however, simply holds for ergodic Markov decision processes. More
14 A. JAŚKIEWICZ
precisely, if the n-step transition probabilities converge to the unique in-
variant probability measure geometrically fast, and the cost functions are
bounded (or more generally satisfy a certain growth hypothesis), then the
aforementioned family of functions is pointwise relatively compact [21, 22].
It is worth pointing out that this requirement is crucial to obtain the opti-
mality inequality in the risk-neutral case; see [27, 28]. In Section 5 we provide
an example that illustrates that also in the risk-sensitive case Condition (B)
cannot be weakened.
We shall need the following two versions of Fatou’s lemma for converging
measures.
Lemma 3. Let {µn} be a sequence of probability measures converging to
µ ∈ Pr(X) and let {hn} be a sequence of measurable nonnegative functions
on X. Then,
h(y)µ(dy) ≤ lim inf
hn(y)µn(dy)
in the following cases:
(a) {µn} converges setwise to µ [i.e.,
f(y)dµn(y) →
f(y)dµ(y)∀f ∈
Bb(X)], and h(x) = lim infn→∞ hn(x);
(b) {µn} converges weakly to µ, and h(x) = inf{lim infn→∞ hn(xn) :xn →
x}; moreover, h ∈ L(X).
Proof. Part (a) is due to Royden [25], page 231, whereas part (b) was
proved by Serfozo [29]. For the proof of lower semicontinuity of h, the reader
is referred to Lemma 3.1 in [22]. �
Now we are in a position to state the main result of the paper. This theo-
rem concerns a study of the risk-sensitive average cost optimality inequality,
which is sufficient to establish the existence of an optimal stationary policy.
Theorem 1. Assume (B) and (W) [or (S)]. Then, for each risk factor
γ > 0, there exist a constant l̂ and a nonnegative function h ∈ L(X) [h ∈
B(X)] and f̂ ∈ F such that
h(x) + l̂ ≥ min
a∈A(x)
γc(x,a) + log
eh(y)q(dy|x,a)
= γc(x, f̂(x)) + log
eh(y)q(dy|x, f̂(x))
RISK-SENSITIVE CONTROL 15
for all x ∈X. Moreover,
= inf
J(x,π) = J(x, f̂).
In other words, l̂/γ is the optimal risk-sensitive average cost and f̂ is a
risk-sensitive average cost optimal stationary policy.
Remark 4. (a) There are two papers [16, 27] that can be treated as
predecessors of our work. They both deal with the optimality inequality but
within two different frameworks. The first work [16] establishes the optimal-
ity equation for the risk-sensitive dynamic programming on a denumarable
state space. In the other one, the result is obtained for Markov control pro-
cesses on an uncountable state space for the risk factor γ = 0. From this
point of view, our result is an extention of Theorem 4.1 in [16] to a general
state space and Theorem 3.8 in [27] to the risk-sensitive case. Moreover, the
common feature of the discussed results is that their proofs are based on the
vanishing discount factor approach. Our proof also relies on this method, and
similarly, as in [27] or [21, 22], makes use of the Fatou lemmas for setwise
and weakly convergent measures.
(b) Finally, it is also worth mentioning that there are papers studying the
optimality equation in the risk-sensitive dynamic programming, which is of
the following form:
h(x) + l̂ = min
a∈A(x)
γc(x,a) + log
eh(y)q(dy|x,a)
.(24)
The constant l̂
is (under suitable assumptions) an optimal cost with respect
to the risk-sensitive average cost criterion. Let us mention and discuss a few
representative papers that deal with equation (24). In [8, 15] Markov control
models satisfying a simultaneous Doeblin condition, on a finite and countable
state space, respectively, are considered. The cost functions are supposed to
be bounded and the risk factor must be sufficiently small. Otherwise, as
argued in [8], the optimality equation need not have a solution.
In [10] Di Masi and Stettner extend the result to a general state space
by retaining bounded cost functions and replacing a simultaneous Doeblin
condition with a very strong assumption on transition probabilities. In [11],
however, they replace this assumption by one imposed on the risk coeffi-
cient. Finally, the class of Markov control models that requires neither any
ergodicity conditions nor the smallness of the risk factor was pointed out by
Jaśkiewicz in [20].
Fairly recently Borkar and Meyn [5] considered Markov decision processes
with unbounded cost functions on a denumarable state space. Their result
16 A. JAŚKIEWICZ
assumes the following: the state space is irreducible under all Markov poli-
cies, the costs are norm-like, and there exists a policy that induces a finite
average risk-sensitive cost. Moreover, their proof is based on a multiplicative
ergodic theorem that was studied in more detail in [1].
Proof of Theorem 1. Let {βn} be a sequence of discount factors
converging to 1 for which (7) holds. Defining
l̂ := l = lim
(1 − βn)mβn
and applying (6), we note that
≤ inf
J(x,π)(25)
for any x ∈X. Assume for a while that inequality (23) is satisfied and there
exists f̂ ∈ F as in the statement of Theorem 1. We prove that f̂ is an optimal
policy. From (23), we have
h(x) ≥ γc(x, f̂(x)) − l̂ + log
eh(y)q(dy|x, f̂(x)).
By iteration of this inequality n times, we obtain
h(x) ≥ logEπx exp
γc(xk, f̂(xk)) + h(xn+1)
− (n+ 1)l̂.
Since h is nonnegative, we infer
+ l̂≥
Jn+1(x, f̂)
with Jn+1(x, f̂) defined in (3). Letting n→∞, it follows
≥ J(x, f̂), x ∈X.(26)
Hence, (25) and (26) together imply
= J(x, f̂) = inf
J(x,Π)
for each x ∈X.
We next focus on showing inequality (23). Let n≥ 1 and put hn := hβn ,
fn := fβn. Note that (19) can be rewritten in the following form:
(1 − βn)mβn + hn(x) = min
a∈A(x)
γc(x,a) + log
eβnhn(y)q(dy|x,a)
= γc(x, fn(x)) + log
eβnhn(y)q(dy|x, fn(x)).
RISK-SENSITIVE CONTROL 17
(i) Assume first (S) and define
h(x) = lim inf
hn(x).
Taking the lim inf on both sides of (27), we get
lim inf
((1 − βn)mβn + hn(x))
= l̂ + h(x) = lim inf
a∈A(x)
γc(x,a) + log
eβnhn(y)q(dy|x,a)
Making use of Lemma 3(a) and the measurable selection theorem (see Propo-
sition D.5(a) in [17]), one can prove that there exists f̂ ∈ F such that (23)
holds.
(ii) Now assume (W). Fix x0 ∈X and choose any xn → x0, n→∞. Take
a subsequence {nk} of positive integers such that
lim inf
hn(xn) = lim
hnk(xnk).
Then by (27),
lim inf
((1− βn)mβn + hn(xn))
= l̂ + lim inf
hn(xn) = l̂ + lim
hnk(xnk)
= lim
a∈A(xnk )
γc(xnk , a) + log
eβnkhnk (y)q(dy|xnk , a)
= lim
γc(xnk , fnk(xnk)) + log
eβnkhnk (y)q(dy|xnk , fnk(xnk))
Note that G = {x0} ∪ {xn} is compact in X. From the upper semicontinu-
ity of x 7→A(x), compactness of every A(z) and Berge’s theorem (see [2] or
Theorem 7.4.2 in [23]), it follows that
z∈GA(z) is compact in A. There-
fore, {fnk(xnk)} has a subsequence converging to some a0 ∈A. By (W)(i),
a0 ∈ A(x0), that is, (x0, a0) ∈ K. Without loss of generality, assume that
fnk(xnk) → a0, k →∞. By the lower semicontinuity of the cost function c
and (28), we have
l̂ + lim inf
hn(xn) ≥ γc(x0, a0) + lim
eβnkhnk (y)q(dy|xnk , fnk(xnk)).
This and Lemma 3(b) imply that
l̂ + lim inf
hn(xn) ≥ γc(x0, a0) + log
eh̃(y)q(dy|x0, a0),
where eh̃ is the generalized lim inf of the sequence eh̃k = ehnk . Clearly, h≤ h̃.
By Lemma 3(b), h ∈L(X). Thus,
l̂ + lim inf
hn(xn) ≥ γc(x0, a0) + log
eh(y)q(dy|x0, a0).(29)
18 A. JAŚKIEWICZ
Since xn → x0 was chosen arbitrarily, we infer from (29) that
l̂ + h(x0) ≥ γc(x0, a0) + log
eh(y)q(dy|x0, a0).
The last inequality shows that, for any x ∈X, there exists an ax ∈A(x) such
l̂ + h(x) ≥ γc(x,ax) + log
eh(y)q(dy|x,ax)
≥ min
a∈A(x)
γc(x,a) +
eh(y)(y)q(dy|x,a)
By our compactness–semicontinuity assumptions and Proposition D.5(b) in
[17], there exists some f̂ ∈ F such that (23) holds. �
5. A discussion. This section is devoted to a discussion of Condition (B).
We start with revisiting Example 3.1 in [8].
Example 1. Put X = {0,1}, A = {a}, c(x) := c(x,a) = x and the tran-
sition matrix is as follows:
ρ 1 − ρ
where ρ ∈ (0,1). Recall that the following was proved.
Let us consider three cases for the risk factor γ:
(I) γ <− log(1− ρ),
(II) γ = − log(1− ρ),
(III) γ >− log(1− ρ).
Then if (I) or (II) hold, the optimal risk-sensitive average cost equals 0
and is independent of the initial state. In case (III) we have J∗(0) = 0 and
J∗(1) = 1 +
log(1−ρ)
> 0. In addition, it is interesting to observe that, for
(II) and (III) cases, there does not exist a function h :X 7→ R such that
optimality inequality (23) is satisfied. Indeed, to see this take x = 1 and
consider (III). The optimality inequality is then as follows:
γJ∗(1) + h(1) = γ + log(1 − ρ) + h(1) ≥ γ + log(eh(1)(1− ρ) + eh(0)ρ).
Note that the right-hand side is strictly greater than γ + log(eh(1)(1 − ρ)),
which equals to the left-hand side. Similar calculations for case (II) also
lead to a contradiction. Hence, although an optimal cost is constant, the
optimality inequality need not have a solution.
Now we turn to checking Condition (B). Let Vβ be as in Lemma 2. Clearly,
Vβ = w
β for N ≥ 1 and Vβ(0) = 0. Then, by (8) under (I), we get
Vβ(1) = γ + log[e
βVβ(1)(1 − ρ) + ρ] < γ + log[eVβ(1)(1 − ρ) + ρ].
RISK-SENSITIVE CONTROL 19
Hence,
Vβ(1) < log
eγ(1− ρ)
1− eγ(1− ρ)
∀β ∈ (0,1),
and consequently, supβ∈(0,1) hβ(x) < +∞.
Now let the risk factor γ be as in (III). Then by (8),
Vβ(1) > γ + log(1 − ρ) + βVβ(1),
which in turn implies that
Vβ(1) >
γ + log(1− ρ)
Thus, hβ(1) = Vβ(1) goes to the infinity when β ր 1.
For case (II), we obtain
Vβ(1) = − log(1− ρ) + log[e
βVβ(1)(1− ρ) + ρ]
= βVβ(1) + log
1 + e−βVβ(1)
1 − ρ
If Vβ(1) ր +∞ when β ր 1, then the right-hand side of (31) also goes to the
infinity. On the contrary, assume that supβ∈(0,1) Vβ(1) ≤C for some constant
C > 0. Then,
Vβ(1) ≥
log[1 + e−Cρ/(1 − ρ)]
which leads to a contradiction when β ր 1. In consequence, in case (II) the
family {hβ(1)} does not satisfy Condition (B) either.
Therefore, the following conclusion can be drawn. Condition (B) is nec-
essary to obtain a solution to the optimality inequality.
For a verification of Condition (B), one can use Lemma 4 below. For a
similar result in the risk-neutral, case we refer to [27, 28]. For some η ≥ 0,
define the stopping time
τ = τ(β) := inf{n≥ 0 :Vβ(xn) ≤mβ + η}.
Lemma 4. For η ≥ 0, β ∈ (0,1) and x ∈X,
hβ(x) ≤ η + inf
logEπx exp
γc(xk, ak)
Proof. By Lemma 2(b), (c) and the fact that Vβ(y) ≥ 0, y ∈ X , we
Vβ(x) = min
a∈A(x)
γc(x,a) + log
eβVβ(y)q(dy|x,a)
< γc(x,a) + log
eVβ(y)q(dy|x,a)
20 A. JAŚKIEWICZ
for each x ∈X. Subtracting mβ from both sides in (32), we obtain
Vβ(x) −mβ < γc(x,a) + log
e(Vβ(y)−mβ )q(dy|x,a).
Iteration of this inequality up to the stopping time τ yields
Vβ(x) −mβ < logE
c(xk,ak)+η
= η + logEπx exp
c(xk, ak)
Since π ∈ Π is an arbitrary policy, we easily get the conclusion. �
Note that the fact
Eπx exp
γc(xk, ak)
< +∞(33)
has the following interpretation: before the process will reach “good states,”
the incurred costs at “early stages” should not be too large. Indeed, let us
define a set D as follows. We say that
x ∈D iff Vβ(x) ≤mβ + η
for a certain η ≥ 0. Clearly, D 6= ∅. Denote by τD the first return time of
the process, governed by fβ, to set D. Certainly, if (33) holds with τ := τD,
then Condition (B) is satisfied.
In Example 1 we can take D = {0} and η = 0, since Vβ(0) ≤ 0 + 0. If γ is
as in (I), then (33) holds:
E1 exp
τ0−1∑
γc(xk)
enγ(1− ρ)n−1ρ =
eγ(1− ρ)
1− eγ(1− ρ)
In other cases (33) fails to hold and, in addition, the earlier calculations
show that hβ(1) = +∞.
Summing up, the presented example shows that, without Condition (B)
imposed on the family of functions {hβ(x)}, β ∈ (0,1), a solution to the
optimality inequality need not exist, and moreover, the optimal risk-sensitive
average cost may depend on the initial state. In view of the above discussion,
Condition (B) is designed to prevent the accrual of infinite expected costs.
Namely, the costs incurred at transient states, that may be occupied only
at “early stages,” have an important and definite influence on a long-run
performance measure. Therefore, Condition (B) requires the model to be
sort of communicating insofar as certain sets of “good states” to be reached
sufficiently fast. Then, the optimal risk-sensitive average cost is constant and
the optimality inequality takes place. In addition, it is worth mentioning that
RISK-SENSITIVE CONTROL 21
the ergodicity itself of a Markov process/chain does not help so much as in
the risk-neutral case. In other words, for an ergodic Markov chain, it may
happen that the optimal risk-sensitive average cost depends on the initial
state as in Example 1. Moreover, in this example one can even prove in
a straightforward way that under case (I) [either under Condition (B) or
for sufficiently small risk factors], the optimality equation (24) is satisfied.
Therefore, it would be interesting to know whether Condition (B) (together
with some compactness–continuity assumptions) is sufficient to obtain a
solution to the optimality equation. There is a conjecture that, since in the
risk-neutral case a counterpart of Condition (B) is not sufficient [7], neither
is it in the risk-sensitive setting. But this question is beyond the scope of
the paper and remains open.
APPENDIX
The lemma below establishes a variational formula for the logarithmic
moment-generating function. The reader is referred to Theorem 4.5.1 and
Proposition 1.4.2 in [12] for its proof.
Lemma A. Let X be a Polish space, h a measurable function mapping
on X into R, which is either bounded from below or bounded from above,
and ν a probability measure on X .
(a) Then, we have the variational formula
ehdν = sup
−R(µ‖ν) +
where
∆ = {µ ∈ Pr(X ) :R(µ‖ν) < +∞}.
(b) Let µ0 denote the probability measure on X , which is µ0 ≪ ν and
satisfies
(x) =
eh(x)∫
eh dν
Then, the supremum in the variational formula is attained uniquely at µ0.
Acknowledgments. A part of this research was done while the author
was a Humboldt research fellow and visiting the University of Ulm. The
author gratefully acknowledges support from the Alexander von Humboldt
Foundation.
The second part of this paper was written at the Institute of Mathematics
and Computer Science, Wroc law University of Technology.
The author is greatly indebted to Professor Ulrich Rieder for drawing
her attention to paper [16], suggesting the problem and for several helpful
conversations.
22 A. JAŚKIEWICZ
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Institute of Mathematics and Computer Science
Wroc law University of Technology
Wybrzeże Wyspiańskiego 27
PL-50-370 Wroc law
Poland
E-mail: [email protected]
http://www.ams.org/mathscinet-getitem?mr=2206685
http://www.ams.org/mathscinet-getitem?mr=2247715
http://www.ams.org/mathscinet-getitem?mr=0752692
http://www.ams.org/mathscinet-getitem?mr=0198505
http://www.ams.org/mathscinet-getitem?mr=0151555
http://www.ams.org/mathscinet-getitem?mr=0378841
http://www.ams.org/mathscinet-getitem?mr=1250112
http://www.ams.org/mathscinet-getitem?mr=1645435
http://www.ams.org/mathscinet-getitem?mr=0705462
http://www.ams.org/mathscinet-getitem?mr=1731299
http://www.ams.org/mathscinet-getitem?mr=1093001
mailto:[email protected]
Introduction and the model
Preliminaries
A solution to the auxiliary discounted minimax problem
A solution to the risk-sensitive control problem
A discussion
Appendix
Acknowledgments
References
Author's addresses
|
0704.0395 | A Study of $B_{d}^0 \to J/\Psi \eta^{(\prime)}$ Decays in the pQCD
Approach | ZJOU-PHY-TH-07-02
NJNU-TH-07-11
A Study of B0d → J/Ψη(′) Decays in the pQCD Approach
Xin Liua∗, Zhen-Jun Xiaob†, Hui-Sheng Wangc
a. Department of Physics, Zhejiang Ocean University,
Zhoushan, Zhejiang 316000, P.R. China
b. Department of Physics and Institute of Theoretical Physics,
Nanjing Normal University, Nanjing, Jiangsu 210097, P.R. China and
c. Department of Applied Mathematics and Physics,
Anhui University of Technology and Science,
Wuhu, Anhui 241000, P.R. China
(Dated: November 4, 2018)
Abstract
Motivated by the very recent measurement of the branching ratio of B0d → J/ψη decay, we
calculate the branching ratios of Bd
0 → J/ψη and Bd0 → J/Ψη′ decays in the perturbative
QCD (pQCD) approach. The pQCD predictions for the branching ratios of considered decays
are: BR(B0d → J/Ψη) = (1.96
+9.68
−0.65) × 10−6, which is consistent with the first experimental
measurement within errors; while BR(B0d → J/Ψη′) = (1.09
+3.76
−0.25) × 10−6, very similar with
B0d → J/Ψη decay and can be tested by the forthcoming LHC experiments. The measurements
of these decay channels may help us to understand the QCD dynamics in the corresponding
energy scale, especially the reliability of pQCD approach to these kinds of B meson decays.
PACS numbers: 13.25.Hw, 12.38.Bx, 14.40.Nd
∗ [email protected]
† [email protected]
http://arxiv.org/abs/0704.0395v1
Very recently, the first observation of B0d → J/Ψη decay was reported by Belle Collab-
oration [1], and the branching ratio measured is
BR(B0d → J/Ψη) = (9.5± 1.7(stat)± 0.8(syst))× 10−6, (1)
which is consistent with the currently available theoretical predictions [1, 2, 3].
Up to now, the theoretical calculations for the branching ratios of Bd → J/Ψη(′) decays
were obtained by using the heavy quark factorization approximation in Ref. [2], or from
the measured J/Ψπ0 and J/ΨK0 branching ratios[3, 4, 5] based on the assumption of
the SU(3) flavor symmetry of strong interaction. In this paper, we will calculate the
branching ratios of B0d → J/Ψη and B0d → J/Ψη(′) decays directly by employing the
low energy effective Hamiltonian [6] and the perturbative QCD (pQCD) factorization
approach [7, 8, 9].
The paper is organized as follows: we present the formalism used in the calculation of
B0d → J/ψη(′) decays in Sec. I. In Sec. II, we show the numerical results and compare
them with the measured values. A short summery and some conclusions are also included
in this section.
I. FORMALISM AND PERTURBATIVE CALCULATIONS
The pQCD approach has been developed earlier from the QCD hard-scattering ap-
proach [7], and has been used frequently to calculate various B meson decay channels
[7, 8, 9, 10]. For two body charmless hadronic Bd,s → Mη(′) (here M stands for the
pseudo-scalar or vector light mesons composed of the light quarks u, d, s) decays, the
pQCD predictions generally agree well with the measured values [9, 10, 11].
In Refs. [12, 13], the authors calculated B → D∗sK,D
s and Bs → D(∗)+D(∗)−
decays and found that the pQCD approach works well for such decays. Here we try to
apply the pQCD approach to calculate the B meson decays involving the heavier J/Ψ
meson as one of the two final state mesons.
A. Formulism
In pQCD approach, the decay amplitude of B → J/ΨP (P = η, η(′) here) decay can
bo written conceptually as the convolution,
A(B →M1M2) ∼
d4k1d
4k3 Tr
C(t)ΦB(k1)ΦJ/Ψ(k2)ΦP (k3)H(k1, k2, k3, t)
, (2)
where the term “Tr” denotes the trace over Dirac and color indices. C(t) is the Wilson
coefficient which results from the radiative corrections at short distance. In the above
convolution, C(t) includes the harder dynamics at larger scale thanMB scale and describes
the evolution of local 4-Fermi operators from mW (the W boson mass) down to t ∼
Λ̄MB) scale, where Λ̄ ≡ MB −mb. The function H(k1, k2, k3, t) is the hard part and
can be calculated perturbatively. The function ΦM is the wave function which describes
hadronization of the quark and anti-quark to the mesonM . While the functionH depends
on the process considered, the wave function ΦM is independent of the specific process.
Using the wave functions determined from other well measured processes, one can make
quantitative predictions here.
Using the light-cone coordinates the B meson and the two final state meson momenta
can be written as
(1, 1, 0T ), P2 =
(1, r2, 0T ), P3 =
(0, 1− r2, 0T ), (3)
respectively, where r = MJ/Ψ/MB, and the light meson masses m
η have been ne-
glected. The longitudinal polarization vector of the J/Ψ meson, ǫL, is given by ǫL =
2MJ/Ψ
(1,−r2, 0T ). Putting the light (anti-) quark momenta in B, J/Ψ and η(
′) mesons
as k1, k2, and k3, respectively, we can choose
k1 = (x1P
1 , 0,k1T ), k2 = (x2P
2 , 0,k2T ), k3 = (0, x3P
3 ,k3T ). (4)
Then, for B → J/Ψη decay for example, the integration over k−1 , k−2 , and k+3 in eq.(2)
will lead to
A(B → J/Ψη′) ∼
dx1dx2dx3b1db1b2db2b3db3
C(t)ΦB(x1, b1)ΦJ/Ψ(x2, b2)Φη(x3, b3)H(xi, bi, t)St(xi) e
−S(t)] ,(5)
where bi is the conjugate space coordinate of kiT , and t is the largest energy scale in
functionH(xi, bi, t). The large logarithms ln(mW/t) are included in the Wilson coefficients
C(t). The large double logarithms (ln2 xi) on the longitudinal direction are summed by
the threshold resummation [14], and they lead to St(xi) which smears the end-point
singularities on xi. The last term, e
−S(t), is the Sudakov form factor which suppresses the
soft dynamics effectively [15]. Thus it makes the perturbative calculation of the hard
part H applicable at intermediate scale, i.e., MB scale. We will calculate analytically the
function H(xi, bi, t) for the considered decays in the first order in αs expansion and give
the convoluted amplitudes in next section.
B. The B0d → J/Ψη(
′) Decays
The low energy effective Hamiltonian for decay modes B0d → J/ψη(
′) can be written as
Heff =
[VcbV
cd (C1(µ)O
1(µ) + C2(µ)O
2(µ))] , (6)
with the four-fermion operators
Oc1 = d̄αγ
µ(1− γ5)cβ · c̄βγµ(1− γ5)bα , Oc2 = d̄αγµ(1− γ5)cα · c̄βγµ(1− γ5)bβ (7)
where the Wilson coefficients Ci(µ) (i = 1, 2), we will use the leading order (LO) expres-
sions, although the next-to-leading order (NLO) results already exist in the literature [6].
This is the consistent way to cancel the explicit µ dependence in the theoretical formulae.
For the renormalization group evolution of the Wilson coefficients from higher scale to
lower scale, we use the formulae as given in Ref.[16] directly.
FIG. 1: Typical Feynman diagrams contributing to the Cabibbo- and color- suppressed B0d →
J/Ψη(
′) decays.
As for B meson wavefunction, we make use of the same parameterizations as used
in the studies of different processes [16]. For vector J/ψ meson, in terms of the nota-
tion in Ref. [17], we decompose the nonlocal matrix elements for the longitudinally and
transversely polarized J/ψ mesons into
ΦJ/Ψ(x) =
mJ/ψǫ/ LΨ
L(x) + ǫ/ LP/Ψ
, (8)
Here, ΨL denote for the twist-2 distribution amplitudes, and Ψt for the twist-3 distri-
bution amplitudes. x represents the momentum fraction of the charm quark inside the
charmonium.
The J/ψ meson asymptotic distribution amplitudes read as [18]
ΨL(x) = 9.58
x(1− x)
x(1 − x)
1− 2.8x(1− x)
Ψt(x) = 10.94
(1− 2x)2
x(1− x)
1− 2.8x(1− x)
. (9)
It is easy to see that both the twist-2 and twist-3 DAs vanish at the end points due to
the factor [x(1− x)]0.7.
From the effective Hamiltonian (6), the Feynman diagrams corresponding to the con-
sidered decay are shown in Fig.1. With the meson wave functions and Sudakov factors,
the hard amplitude is given as
Feη = 8πCFm
dx1dx3
b1db1b3db3 φB(x1, b1)
(1− r2)
(1 + x3(1− r2))φAη (x3, b3) + r0(1− 2x3)
·φPη (x3, b3)
(1− 2x3) + r2(1 + 2x3)
φTη (x3, b3)
·αs(t1e) he(x1, x3, b1, b3) exp[−Sab(t1e)]
1− (1− x1)r2φPη (x3, b3)− x1r2φAη (x3, b3)
·αs(t2e)he(x3, x1, b3, b1) exp[−Sab(t2e)]
. (10)
where r0 = m
0/mB; CF = 4/3 is a color factor. The function he, the scales t
e and the
Sudakov factors Sab are displayed in Appendix A.
For the non-factorizable diagrams 1(c) and 1(d), all three meson wave functions are
involved. The integration of b3 can be performed using δ function δ(b3 − b1), leaving only
integration of b1 and b2. For the concerned operators, the corresponding decay amplitude
Meη =
dx1dx2 dx3
b1db1b2db2 φBs(x1, b1)
2rrcφ
J/Ψ(x2, b2)φ
η (x3, b2)− 4rr0rcφtJ/Ψ(x2, b2)φTη (x3, b2)
2 + x3(1− 2r2)
φLJ/Ψ(x2, b2)φ
η (x3, b2)
x3r0 + (x2 − x3)r0r2
φLJ/Ψ(x2, b2)φ
η (x3, b2)
·αs(tf )hf(x1, x2, x3, b1, b2) exp[−Scd(tf )]} . (11)
where rc = mc/mB,mc is the mass for c quark.
For the B0d → J/Ψη′ decay, the Feynman diagrams are obtained by replacing the η
meson in Fig. 1 with the meson η′. The corresponding expressions of decay amplitudes will
be similar with those as given in Eqs.(10-11), since the η and η′ are all light pseudoscalar
mesons and have the similar wave functions. The expressions of B0d → J/Ψη′ decay can
be obtained simply by the following replacements
φAη −→ φAη′ , φPη −→ φPη′ , φTη −→ φTη′ , r0 −→ r′0. (12)
For the η−η′ system, there exist two popular mixing basis: the octet-singlet basis and
the quark-flavor basis [19, 20]. Here we use the quark-flavor basis [19] and define
ηq = (uū+ dd̄)/
2, ηs = ss̄. (13)
The physical states η and η′ are related to ηq and ηs through a single mixing angle φ,
= U(φ)
cosφ − sinφ
sin φ cos φ
. (14)
The three input parameters fq, fs and φ in the quark-flavor basis have been extracted
from various related experiments [19, 20]
fq = (1.07± 0.02)fπ, fs = (1.34± 0.06)fπ, φ = 39.3◦ ± 1.0◦, (15)
where fπ = 130 MeV. In the numerical calculations, we will use these mixing parameters
as inputs. It worth of mentioning that the effects of possible gluonic component of η′
meson will not considered here since it is small in size [10, 21, 22].
For B0d → J/Ψη decay, by combining the contributions from different diagrams, the
total decay amplitude can be written as
M(B0d → J/Ψη) = VcbV ∗cdF1(φ)
FeηfJ/Ψ
+MeηC2
where the relevant mixing parameter is F1(φ) = cos φ/
It should be mentioned that the Wilson coefficients Ci = Ci(t) in Eq. (16) should
be calculated at the appropriate scale t using equations as given in the Appendices of
Ref. [16]. Here the scale t in the Wilson coefficients should be taken as the same scale
appeared in the expressions of decay amplitudes in Eqs. (10) and (11). This is the way
in pQCD approach to eliminate the scale dependence. In order to estimate the effect of
higher order corrections, however, we introduce a scale factor at = 1.0± 0.2 and vary the
scale tmax as described in Appendix A.
Similarly, the decay amplitudes for B0d → J/Ψη′ decay can be obtained easily from
Eq.(16) by the following replacements of F1(φ) → F ′1(φ) = sinφ/
II. NUMERICAL RESULTS AND DISCUSSIONS
In this section, we will calculate the branching ratios for those considered decay modes.
The input parameters and the wave functions to be used are given in Appendix B. In
numerical calculations, central values of input parameters will be used implicitly unless
otherwise stated.
With the complete decay amplitudes, we can obtain the decay width for the considered
decays,
Γ(B0d → J/ψη(
′)) =
(1− r2)
M(B0d → J/ψη(
. (17)
By employing the quark-flavor scheme of η−η′ system and using the mixing parameters
as given in Eq. (15), one finds the branching ratios for the considered two decays with
error bars as follows:
Br( B0d → J/Ψη) =
1.96+0.71−0.50(ωb)
+9.65
−0.39(at)
+0.32
+0.13(a2)
+0.14
−0.13(fJ/Ψ)
× 10−6, (18)
Br( B0d → J/Ψη′) =
1.09+0.32−0.24(ωb)
+3.73
+0.01(at)
+0.28
+0.01(a2)
+0.08
−0.07(fJ/Ψ)
× 10−6, (19)
where the main errors are induced by the uncertainties of ωb = 0.40 ± 0.05 GeV, at =
1.0 ± 0.2, a2 = 0.115 ± 0.115 and fJ/Ψ = 0.405 ± 0.014 GeV , respectively. One can see
that the pQCD predictions are sensitive to the variations of ωb and at.
For B0d → J/Ψη decay, the central value of the pQCD prediction for Br(B0d → J/Ψη)
is a factor of 4 smaller than the measured value as given in Eq. (1) [1]. But the pQCD
prediction is in fact still consistent with Belle’s first measurement if we take the large
theoretical and experimental errors into account. By varying the scale factor at in the
range of at = [0.8, 1.0], for example, the central value of Br(B → J/Ψη) will change in the
range of [0.2, 1.1]×10−5 accordingly. It is not difficult to understand such at dependence.
Since the J/Ψ meson is much heavier than light mesons, and therefore moving not as fast
as those light meson when B meson is decaying. So a small decrease of the scale ti will
lead to a larger Wilson coefficients C1,2(t) and αs(ti), and consequently results in a larger
decay rate.
For B0d → J/Ψη′ decay, only experimental upper limit (at 90% C.L) is available now:
BR(B0 → J/Ψη′) < 6.3 × 10−5 [4, 5]. The pQCD prediction for the branching ratio of
B0d → J/Ψη′ decay is very similar in magnitude with that of B0d → J/Ψη, consistent with
the upper limit and will be tested in the forthcoming LHC experiments.
At the leading order, only the tree Feynman diagrams as shown in Fig. 1 contribute to
B0d → J/Ψη(′) decays. There exists no CP violation in these decays within the standard
model, since there is only one kind of Cabibbo-Kabayashi-Muskawa (CKM) phase involved
in the corresponding decay amplitudes, as can be seen from eq. (16).
In short, we calculated the branching ratios of B0d → J/Ψη and B0d → J/Ψη′ decays at
the leading order by using the pQCD factorization approach. Besides the usual factoriz-
able diagrams, the non-factorizable spectator diagrams are also calculated analytically in
the pQCD approach. By keeping the transverse momentum kT , the end-point singularity
disappears in our calculation.
From our calculations and phenomenological analysis, we found the following results:
• Using the quark-flavor scheme, the pQCD predictions for the branching ratios are
Br(B0d → J/Ψη) =
1.96+9.68−0.65
× 10−6, (20)
Br(B0d → J/Ψη′) =
1.09+3.76−0.25
× 10−6, (21)
where the various errors as specified previously have been added in quadrature.
• The major theoretical errors of the pQCD predictions are induced by the uncertain-
ties of the hard energy scale ti’s and the parameters ωb.
Acknowledgments
X. Liu would like to acknowledge the financial support of The Scientific Research
Start-up Fund of Zhejiang Ocean University under Grant No.21065010706. This work was
partially supported by the National Natural Science Foundation of China under Grant
No.10575052, and by the Specialized Research Fund for the Doctoral Program of Higher
Education (SRFDP) under Grant No. 20050319008.
APPENDIX A: RELATED FUNCTIONS
We show here the function hi’s, coming from the Fourier transformations of the function
H(0),
he(x1, x3, b1, b3) = K0
x1x3(1− r2)mBb1
θ(b1 − b3)K0
x3(1− r2)mBb1
x3(1− r2)mBb3
+ θ(b3 − b1)K0
x3(1− r2)mBb3
x3(1− r2)mBb1
St(x3), (A1)
hf(x1, x2, x3, b1, b2) =
θ(b2 − b1)I0(MB
x1x3(1− r2)b1)K0(MB
x1x3(1− r2)b2)
+ (b1 ↔ b2)
K0(MBF(1)b2), for F
(1) > 0
0 (MB
| b2), for F 2(1) < 0
, (A2)
where J0 is the Bessel function, K0 and I0 are the modified Bessel functions with
K0(−ix) = −(π/2)Y0(x) + i(π/2)J0(x), and F(j)’s are defined by
F 2(1) = (x1 − x2)x3(1− r2) + r2c , (A3)
F 2(2) = (x1 − x2)x3(1− r2) + r2c . (A4)
The threshold resummation form factor St(xi) is adopted from Ref. [17]
St(x) =
21+2cΓ(3/2 + c)√
πΓ(1 + c)
[x(1 − x)]c, (A5)
where the parameter c = 0.3. This function is normalized to unity.
The Sudakov factors used in the text are defined as
Sab(t) = s
x1mB/
2, b1
x3mB/
2, b3
(1− x3)mB/
2, b3
ln(t/Λ)
− ln(b1Λ)
ln(t/Λ)
− ln(b3Λ)
, (A6)
Scd(t) = s
x1mB/
2, b1
x2mB/
2, b2
(1− x2)mB/
2, b2
x3mB/
2, b1
(1− x3)mB/
2, b1
ln(t/Λ)
− ln(b1Λ)
ln(t/Λ)
− ln(b2Λ)
, (A7)
where the function s(q, b) are defined in the Appendix A of Ref. [16]. The scale ti’s in the
above equations are chosen as
t1e = at ·max(
x3(1− r2)MB, 1/b1, 1/b3),
t2e = at ·max(
x1(1− r2)MB, 1/b1, 1/b3),
tf = at ·max(
x1x3(1− r2)MB,
(x1 − x2)x3(1− r2) + r2cMB, 1/b1, 1/b2), (A8)
where at = 1.0±0.2 and r =MJ/Ψ/MB. The scale ti’s are chosen as the maximum energy
scale appearing in each diagram to kill the large logarithmic radiative corrections.
APPENDIX B: INPUT PARAMETERS AND WAVE FUNCTIONS
The masses, decay constants, QCD scale and B0d meson lifetime are
(f=4)
= 250MeV, fπ = 130MeV, fJ/Ψ = 405MeV,
0 = 1.08GeV, MB0d = 5.28MeV, MJ/Ψ = 3.097GeV,
MW = 80.41GeV, τB0
= 1.54× 10−12s. (B1)
For the CKM matrix elements, here we adopt the Wolfenstein parametrization for the
CKM matrix, and take λ = 0.2272, A = 0.818, ρ = 0.221 and η = 0.340 [4].
For the B meson wave function, we adopt the model
φB(x, b) = NBx
2(1− x)2exp
(ωbb)
, (B2)
where ωb is a free parameter and we take ωb = 0.40± 0.05 GeV in numerical calculations,
and NB = 91.745 is the normalization factor for ωb = 0.40 for the B meson.
The wave function for dd̄ components of η(′) meson is given by
Φηdd̄(p, x, ζ) ≡
P/φAηdd̄(x) +m
(x) + ζm
0 (v/n/− v · n)φTηdd̄(x)
, (B3)
where p and x are the momentum and the momentum fraction of ηdd̄ respectively, while
φAηdd̄, φ
and φTηdd̄ represent the axial vector, pseudoscalar and tensor components of the
wave function respectively. We here assume that the wave function of ηdd̄ is same as the
π wave function based on SU(3) flavor symmetry. The parameter ζ is either +1 or −1
depending on the assignment of the momentum fraction x.
The explicit expression of chiral enhancement scale m
0 = m
0 is given by [21]
[m2η cos
2 φ+m2η′ sin
(m2η′ −m2η) cosφsinφ], (B4)
and numerically m
0 = 1.07MeV for mη = 547.5 MeV, mη′ = 957.8 MeV, fq = 1.07fπ,
fs = 1.34fπ and φ = 39.3
For the distribution amplitude φAηq , φ
and φTηq , we utilize the results for π meson
obtained from the light-cone sum rule [23] including twist-3 contributions:
φAηq(x) =
fqx(1 − x)
1 + a
5(1− 2x)2 − 1
21(1− 2x)4 − 14(1− 2x)2 + 1
, (B5)
φPηq(x) =
30η3 −
3(1− 2x)2 − 1
−3η3ω3 −
ρ2ηq −
ρ2ηq(s)a
35(1− 2x)4 − 30(1− 2x)2 + 3
,(B6)
φTηq(x) =
fq(1− 2x)
+ (5η3 −
η3ω3 −
ρ2ηq −
ρ2ηqa
2 )(10x
2 − 10x+ 1)
, (B7)
with the updated Gegenbauer moments [24]
2 = 0.115, a
4 = −0.015, ρηq = 2mq/mqq, η3 = 0.015, ω3 = −3.0. (B8)
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http://arxiv.org/abs/hep-ex/0603003
http://www.slac.stanford.edu/xorg/hfag
http://arxiv.org/abs/hep-ph/0701146
http://arxiv.org/abs/hep-ph/0703162
http://arxiv.org/abs/hep-ph/0703187
Formalism and Perturbative Calculations
Formulism
The Bd0 J/(') Decays
Numerical results and Discussions
Acknowledgments
Related Functions
Input parameters and wave functions
References
|
0704.0396 | Finite-temperature phase transitions in a two-dimensional boson Hubbard
model | Finite-temperature phase transitions in a two-dimensional boson Hubbard model
Min-Chul Cha1 and Ji-Woo Lee2
Department of Applied Physics, Hanyang University, Ansan 426-791, Korea
Department of Physics, Myongji University, Yongin 449-728, Korea
We study finite-temperature phase transitions in a two-dimensional boson Hubbard model with
zero-point quantum fluctuations via Monte Carlo simulations of quantum rotor model, and construct
the corresponding phase diagram. Compressibility shows a thermally activated gapped behavior in
the insulating regime. Finite-size scaling of the superfluid stiffness clearly shows the nature of the
Kosterlitz-Thouless transition. The transition temperature, Tc, confirms a scaling relation Tc ∝ ρ
with x = 1.0. Some evidences of anomalous quantum behavior at low temperatures are presented.
PACS numbers: 73.43.Nq, 74.25.Dw, 05.30.Jp
Recently quantum phase transitions[1, 2] have drawn a
lot of attention in systems of interacting particles. Typ-
ically strong interactions suppress the itineracy of par-
ticles to induce a strongly correlated insulating phase,
whereas with weak interactions a conducting phase is
stable. The criticality of these zero-temperature phase
transitions can be investigated at low, but finite, tem-
peratures. How quantum fluctuations associated with a
quantum critical point(QCP) have influence on phases at
finite temperatures [3, 4, 5] is a theoretically interesting
and an experimentally relevant question.
At finite temperatures, it is expected that a quantum
phase transition turns into a classical one with the same
order parameter or disappears. Remnant quantum fluc-
tuations near a QCP may bring anomalous properties[3],
which can be captured by scaling relations, and lead
to crossover behaviors as temperature rises. Some pos-
sibilities such as reentrant behaviors due to the inter-
play of quantum and thermal fluctuations have been
proposed[6].
These issues can be clarified by direct investigations of
a generic quantum mechanical model. So far most of the
theoretical investigations heavily rely on the exact solu-
tion of the quantum Ising model, available strictly in one
dimension[4]. Interacting bosonic systems simulated via
Monte Carlo methods, not suffering from negative sign
problems, will be an ideal place to study these problems.
In previous works, a quantum XY model, equivalent to
hard-core bosons at half-filling, showed the Kosterlitz-
Thouless(KT) transition[7] at finite temperature in two
dimensions[8, 9]. In the model with nearest neighbor
repulsion, destruction of the solid order as well as the
superfluidity by thermal fluctuations was observed [10].
However, generic finite-temperature phase diagrams have
not been constructed.
In this work, we investigate the thermally driven phase
transitions of a two-dimensional quantum rotor model,
which is believed to share the same critical properties of
a soft-core generic boson Hubbard model[11], via Monte
Carlo simulations. The results are summarized in the
phase diagram as shown in Fig. 1. Finite-size scaling
properties of the superfluid stiffness confirm that the na-
ture of the classical phase transition associated with the
destruction of superfluidity is consistent with that of the
KT transition, and clearly support the scenario of the
universal jump at the critical point[12]. Finite tempera-
ture, T , sets the size in the temporal direction, leading to
a scaling behavior[4, 11] Tc ∝ ρ0x with x = 1.0, where Tc
is the transition temperature and ρ0 is the superfluid stiff-
ness at zero temperature. The compressibility diverges at
the transition. In the insulating regime at low tempera-
ture, thermally activated behavior of the compressibility
with a finite energy gap is observed. Some anomalous de-
pendence of energy and specific heat on T , possibly due
to quantum fluctuations, are observed for T < 0.25U .
The Hamiltonian of a boson Hubbard model reads
nj(nj − 1)− µ
nj − w
ibj + b
jbi),(1)
where bj(b
j) is the boson annihilation(creation) operator
0 0.01 0.02 0.03 0.04 0.05
µ=0.9
Normal fluid
Mott Insulator
Superfluid
KT transition
Gapped fluid
FIG. 1: (Color online) Phase diagram on the space of hop-
ping strength, t(=
n0(n0 + 1)w), and temperature, T , in
unit of U . The solid line denotes the classical phase transi-
tions, which terminates at a QCP at T = 0. The dotted line
represents crossover between gapped fluid and normal fluid.
http://arxiv.org/abs/0704.0396v1
at the j-th site, and nj is the number operator. U and
w stand for the strengths of the on-site repulsion and of
the nearest neighbor hopping, respectively, and µ is the
chemical potential.
It is convenient to put µ/U + 1/2 = n0 + n̄ with an
integer n0 and −1/2 < n̄ ≤ 1/2 so that n0 represents the
background number of bosons per site and n̄ is a charge
offset. When n̄ = 0, the density of bosons is fixed to
a commensurate filling across the transition. For non-
integer n̄, however, an integer filling in a Mott insulator
shifts to a non-integer filling in a compressible fluid. We
study the phase transition of the latter case in (2+1)-
dimensional L×L×Lτ square lattices, where L denotes
the size in a spatial dimension and Lτ in the temporal
dimension.
Since the phase transition of the model in Eq. (1)
is characterized by the establishment of phase coher-
ence, we may rewrite the Hamiltonian in terms of
the phase angle θj of bosons by replacing bj(b
j) =√
−iθj (
nj + 1e
iθj ) with nj =
. Under the as-
sumption that the nature of the transition is governed
only by the fluctuations of θj , not those of the hop-
ping strength, we replace nj → n0 so that bj(b†j) =√
−iθj (
n0 + 1e
iθj ). Then, the Hamiltonian is re-
duced to a quantum rotor model
nj(nj − 1)− µ
nj − 2t
cos(θi − θj),(2)
where t =
n0(n0 + 1)w. Here we take the number of
bosons nj ≥ 0.
Through a path integral mapping, we can construct
the corresponding classical action[14]
Jτr (J
r − 1)− ǫµJτr − ln IJxr (2ǫt)− ln IJyr (2ǫt)(3)
with the partition function
∇· ~J=0
{ ~Jr}
~J], (4)
where ǫ = β/Lτ is a lattice constant in the imaginary
time axis for an inverse temperature β, ~Jr is an integer
current at site r = (j, τ) with a spatial index j and a tem-
poral index τ , which is conserved at each site as denoted
by ∇ · ~J = 0, and Im(x) is the modified Bessel func-
tion given by the relation eK cos θ =
m=−∞ Im(K)e
In this work, we investigate the properties of the model
in Eq. 3 via Monte Carlo simulations using a recently
proposed worm algorithm [13]. In order to reduce the
systematic errors in discretizing the imaginary time axis,
we need to take ǫ
tU ≪ 1. We take Uǫ = 0.5 - 2 for
t ≪ U and set the energy unit U = 1.
The superfluid stiffness in a finite system is given by[14]
ρL = β
−1L2−d〈W 2x 〉, (5)
0.02 0.03 0.04 0.05
L=128
0 0.5 1
L exp(-b (t
) µ=0.9
=0.0409
b=1.85
0 10 20 30 40 50 60
0.005
0.015
L=128
0 0.1 0.2 0.3 0.4
L exp(-b (β
) µ=0.9
t=0.034
t=0.034
b=3.35
=28.8
FIG. 2: (Color online) Finite-size scaling behaviors of the su-
perfluid stiffness as a function of (a) hopping strength and
(b) temperature. For both cases, data collapsing onto a sin-
gle curve works fine in terms of the scaling parameter L/ξ as
shown in insets, consistent with the nature of the KT transi-
tion and the universal jump at the critical point.
where Wx = L
r and 〈...〉 denotes the averages
over the probabilites determined by the partition func-
tion of Eq. (4), and d is the spatial dimensionality. Sim-
ilarly the compressibility is
κ = βL−d[〈N2〉 − 〈N〉2], (6)
with N = L−1τ
r . The energy expectation is given
〈H〉 = L−1τ 〈
〉 , (7)
and the specific heat is CV = L
−d(∂〈H〉/∂T ).
We consider the case for µ = 0.9 so that n0 = 1 and
n̄ = 0.4. Figure 2 shows the finite-size scaling behav-
ior of the superfluid stiffness as a function of (a) t and
(b) β. Finite-size scaling properties of the transition can
be obtained by plotting the curves in terms of a scal-
ing variable L/ξ, where ξ is the correlation length. Here
we assume an essential singularity[15] ξ ∼ exp(bδ−1/2),
where δ = t − tc (or β − βc) is a tuning parameter and
b is a non-universal scaling factor. In terms of this scal-
ing variable, we obtain high-quality data collapsing onto
a single curve for different sizes, consistent with the na-
ture of the KT transition. The scaling behavior also sup-
ports the scenario of the universal jump of the superfluid
stiffness[12], (π/2)βcρ∞ = 1, at the critical point in the
thermodynamic limit. By extrapolating the single curves
to the critical point, we find that (π/2)βcρ∞ ≈ (a)1.01
and (b)1.06. These numbers are, however, sensitive to
fitting parameters b and tc(βc).
Figure 3a shows the behavior of the compressibility.
The finite-size scaling ansatz of the compressibility is
written in the form
κ = Lz−dX̃κ(L(t− tc)1/ν , β/Lz), (8)
where X̃κ is a dimensionless scaling function and z is the
dynamical critical exponent. For the generic superfluid-
insulator transition(GSIT), z = 2 is expected[11]. The
crossing behavior of the compressibility curves for differ-
ent sizes at the critical point t0c = 0.023±0.001, therefore,
represents the scaling properties near the QCP, where t0c
is the critical hopping strength at zero temperature. For
different values of µ, we have similar results with t0c just
shifted.
We find that the compressibility diverges at the transi-
tion. In the superfluid side, κ ∼ 1/(t− t0c). This strongly
supports that the longe-range density fluctuations drive
the transition. In the insulating side, the compressibility
has an activated form e−∆gap/T with a finite energy gap
∆gap. This dependence is shown in Fig. 3b for different
t, from which we can calculate ∆gap as shown in the in-
set. For small t, we need a large number of Monte Carlo
steps to obtain equilibrium and have bigger error bars in
determination of ∆gap. The gap vanishes around t = t
as expected.
Thus we have a so-called ’V-shaped’ phase diagram
(Fig.1). In the insulating regime, the Mott insulator
exists at T = 0, which turns into an activated gapped
fluid with a finite energy gap at low temperature. It
gradually disappears in a high-temperature normal fluid.
This crossover line can be specified by the condition
∆gap/T ≈ 1. The phase coherence in a superfluid at
T = 0 is destroyed by quantum fluctuations to form a
QCP or by thermal fluctuations at T > 0 to define clas-
sical phase transitions. The phase boundary in Fig. 1 is
obtained by tuning t for given T (black circles) as well
as by tuning β for a given t (red squares). Note that the
phase boundary follows a scaling relation Tc ∝ |t− t0c |zν ,
which implies that β determines the correlation length in
the temporal direction, where ν is the correlation length
0 0.01 0.02 0.03 0.04 0.05
L=12 Lτ=18
L=16 Lτ=32
L=20 Lτ=50
L=24 Lτ=72
L=28 Lτ=98
µ=0.9
0 50 100 150 200
1e-05
0.0001
0.001
0 0.01 0.02
µ=0.9
t=.005
t=.010
t=.015
t=.020
t=.021
t=.022
t=.023(b)
FIG. 3: (Color online) (a) Compressibility of the boson Hub-
bard model shows behavior of the GSIT with z = 2.0, diverg-
ing at the transition. (b) In the insulating regime, we have
thermally activated behaviors, κ ∼ e−∆gap/T , from which
∆gap can be evaluated. Inset: ∆gap as a function of t, van-
ishing at the QCP.
critical exponent. The boundary in Fig. 1 is consistent
with the expectation zν = 1[11] for the GSIT.
It is interesting to check the predicted scaling relation
[4, 11] Tc ∝ ρ0x in this model. Figure 4 shows that
the zero-temperature superfluid stiffness ρ0, denoted by
dotted line, which obtained via extrapolation of values at
T > 0, follows ρ0 ∝ |t− t0c |, implying that x = 1.0. It is
consistent with the hyperscaling argument[11] suggesting
x = z/(d+ z − 2).
We expect that this quantum criticality disappears
as temperature rises, which means quantum fluctuations
possibly leave some tracks in bulk properties at low tem-
peratures. Figure 5 shows the specific heat, CV , and the
energy expectation values, 〈H〉, as a function of T for dif-
ferent t. Sharp rises of CV in the conducting regime or
round up-rises in the insulating regime are followed by
indents, regions indicated by N, which apparently rep-
resent anomalous behavior due to quantum fluctuations
and disappear at high temperatures for T & 0.25. This
feature strongly suggests a crossover in normal fluid from
quantum mechanical to classical regime. Similarly the
curves of 〈H〉 show bumps, indicated by H, only in the
range where quantum critical fluctuations are expected
to have effects.
In summary, we have investigated the phase transi-
tions at finite temperature in a two-dimensional quan-
tum rotor model in which intrinsic zero-point fluctua-
tions are present. Finite-size scaling of the superfluid
stiffness shows an essential singularity of the KT phase
transition and the universal jump at the critical point.
The compressibility diverges at the transition. In the
insulating regime, the compressibility shows a thermally
activated behavior, κ ∼ e−∆gap/T , from which we can
successfully evaluate the gap. This indicates that the in-
sulating behavior at low temperature gradually crosses
over to the behavior of normal fluid as temperature in-
creases. The transition temperature Tc shows a scaling
behavior Tc ∝ |t − t0c |, showing that finite T limits the
length of quantum fluctuations in the temporal direction,
and a hyperscaling relation Tc ∝ ρ0. The behavior of the
specific heat and the energy suggests that, as tempera-
ture rises, quantum critical regime near a QCP crosses
over to classical regime.
MCC would like to thank Gerardo Ortiz for helpful
discussions and the hospitality of Department of Physics,
Indiana University, where parts of this work were carried
out. This work was supported by Korea Research Fund
grant No. R05-2004-000-11004-0.
0.02 0.025 0.03 0.035 0.04 0.045 0.05
β=200
β=400
µ=0.9
FIG. 4: (Color online) Superfluid stiffness for different β. As
β increases, the size dependence becomes smaller. This allows
us to extrapolate the curves to obtain zero-temperature su-
perfluid stiffness, ρ0, in the thermodynamic limit as denoted
by dotted line. It shows that ρ0 ∝ |t− t
c | with t
c ≈ 0.22.
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Fisher, Phys. Rev. B 40, 546 (1989).
[12] D. R. Nelson and J. M. Kosterlitz, Phys. Rev. Lett., 39,
1201 (1977).
[13] F. Alet and E. S. Sørensen, Phys. Rev. E 67, 015701(R)
(2003); Phys. Rev. E 68, 026702 (2003).
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Phys. Rev. B 49, 12115 (1994).
[15] J. M. Kosterlitz, J. Phys. C 7, 1046 (1974).
0 0.1 0.2 0.3 0.4 0.5
t = 0.005
t = 0.020
t = 0.040
t = 0.050
t = 0.070
t = 0.100
-0.85
t = .005
t = .010
t = .020
t = .025
t = .030
t = .040
0 0.1 0.2 0.3 0.4
t = .040
t = .050
t = .060
t = .070
t = .080
t = .090
t = .010
FIG. 5: (Color online) Specific heat, CV , as a function of T
for different t. Sharp rises in the conducting regime, signa-
ture of the superfluid transition, or round up-rise of CV in the
insulating regime are followed by indents which disappear in
high temperature region, T & 0.25. Insets: The curves of the
energy expectation values, 〈H〉, have bumps at low tempera-
tures possibly due to the effects of quantum fluctuations.
|
0704.0397 | Conditional generation of path-entangled optical NOON states | APS/123-QED
Conditional generation of path-entangled optical NOON states
Anne E. B. Nielsen and Klaus Mølmer
Lundbeck Foundation Theoretical Center for Quantum System Research,
Department of Physics and Astronomy, University of Aarhus, DK-8000 Århus C, Denmark
(Dated: November 4, 2018)
We propose a measurement protocol to generate path-entangled NOON states conditionally from
two pulsed type II optical parametric oscillators. We calculate the fidelity of the produced states
and the success probability of the protocol. The trigger detectors are assumed to have finite dead
time, and for short pulse trigger fields they are modeled as on/off detectors with finite efficiency.
Continuous-wave operation of the parametric oscillators is also considered.
PACS numbers: 03.65.Wj, 03.67.-a, 42.50.Dv
I. INTRODUCTION
Nonclassical states of light have many applications,
and a number of different protocols exist for the genera-
tion of various classes of states. The two-mode maximally
entangled N -photon states
|NOON〉 =
|N, 0〉+ eiφ|0, N〉
, (1)
the so-called NOON states, are particularly interesting
because a single-photon phase shift of χ induced in one of
the two components changes the relative phase of the two
terms by Nχ. This special property of NOON states may
be utilized to enhance spatial resolution in (quantum) mi-
croscopy and lithography [1], and in interferometry it has
been shown that a certain measurement strategy, using
NOON states, leads to a phase estimation error scaling
as L−1/4N
T if the phase to be estimated is known to
lie within an interval from −π/L to π/L, where NT is the
total number of photons used in the measurements [2].
This is better than the classical shot noise limit ofN
and NOON states are thus useful to perform accurate
measurements and may be a valuable field resource in
sensors. NOON states are also a source of entanglement
with applications in quantum information protocols and
in fundamental studies such as tests of Bells inequality
It is thus of great interest to be able to produce NOON
states, and various NOON state generation schemes have
been suggested theoretically [4, 5, 6, 7, 8, 9] and studied
in experiments [7, 8, 10, 11]. The N = 1 and N = 2
NOON states may be generated by combining either a
single photon and a vacuum state or two single-photon
states on a 50 : 50 beam splitter, but this simple approach
is not directly extendable to N > 2, and we shall thus
mainly be concerned with generation of N = 3 NOON
states in the present paper, even though the suggested
protocol is, in principle, applicable for all N . Mitchell,
Lundeen, and Steinberg have generated NOON states
with N = 3 from a pair of down converted photons
and a local oscillator photon using certain polarization
transforming components and post-selection [10]. In this
experiment, however, the successful generation of the
NOON state is witnessed by a destructive detection of
the state. In the present paper we propose and analyze
in detail a nondestructive generation protocol, which con-
ditions the successful generation of the N -photon NOON
state on the registration of N photo detection events in
other field modes, and which uses as resource only linear
optics and the output from two optical parametric oscil-
lators (OPOs). The protocol does not rely on efficient
photo detection. The analysis is carried out in terms of
Wigner function formalism, and effects of finite detector
efficiency and finite detector dead time are considered.
Conditional generation of nonclassical states occupying
a single mode has been investigated both experimentally
and theoretically [12, 13, 14, 15, 16, 17, 18, 19, 20]. With
the correlated output from a single nondegenerate OPO
it is, for instance, possible to generate n-photon Fock
states of light in the signal beam conditioned on n photo
detections in the idler (trigger) beam [16, 19, 20], and
in principle the entanglement of a highly squeezed two-
mode field from an OPO makes it possible to prepare any
state in the signal beam that can be either measured as
an eigenstate of a suitable observable of the idler beam or
produced as the final state of a generalized measurement.
The basic idea of the protocol proposed in the present
paper is to mix the output from two OPOs and employ
the entanglement to prepare a two-mode state in two of
the output beams by detection of the desired output state
in the remaining beams.
In Sec. II we explain the NOON state generation pro-
tocol in detail. In Sec. III we analyze the performance
of the protocol quantitatively for pulsed OPO sources.
We provide the fidelity of the generated states and the
success probability. In Sec. IV we consider production of
NOON states from continuous-wave OPO sources, and
Sec. V concludes the paper.
II. EXPERIMENTAL SETUP FOR NOON
STATE GENERATION
The experimental setup is illustrated in Fig. 1. Two
pulses of two-mode squeezed states are generated by two
identical OPOs via type II parametric down conversion.
The field mode operators of the modes generated by the
http://arxiv.org/abs/0704.0397v2
first OPO are denoted â+ and â−, respectively, while the
field mode operators of the modes generated by the sec-
ond OPO are denoted b̂+ and b̂−, respectively. For def-
initeness, we assume that the plus modes are vertically
polarized and that the minus modes are horizontally po-
larized. The modes are separated spatially by the first
two polarizing beam splitters, and the third polarizing
beam splitter combines the â− and b̂+ modes, which are
subsequently subjected to the NOON state measurement
proposed in [21] and illustrated for N = 3 in Fig. 1. The
idea behind this measurement is to apply the highly non-
linear operator ÂN = â
− − (b̂+eiθ)N to the state. The
result is only nonzero if either the â− mode or the b̂+
mode contains at least N photons. On the other hand,
if the squeezing is sufficiently small, it is unlikely to have
more than a total of N photons in the two trigger modes,
and by conditioning on the successful application of ÂN ,
we select the pulses of the system where N photon pairs
are generated in one OPO and zero photon pairs in the
other. It is equally probable that the photons originate
from the first OPO or from the second OPO, and, as we
shall see in detail below, the result is that a NOON state
is generated conditionally in the output modes â+ and
As stated in [21], ÂN can be rewritten as a simple
product of single photon annihilation operators
âN− −
â− −
ei2πn/N
, (2)
and it is thus possible to implement ÂN by means of
beam splitters and photo detectors. We first consider odd
values of N . Beam splitters are used to divide the input
into N distinct spatial modes labeled by n = 1, . . . , N .
The beam splitter reflectivities are chosen to obtain the
same expectation value of the intensity in each of the
modes. The vertically polarized modes are then phase
shifted by the factor ei2πn/N+iπ relative to the horizon-
tally polarized modes, i.e., b̂+ → −b̂+ei2πn/N , and finally
polarizing beam splitters with principal planes oriented
at 45◦ relative to the horizontal polarization transform
â− and −b̂+ei2πn/N into (â−− b̂+ei2πn/N )/
2 (the trans-
mitted mode) and (â− + b̂+e
i2πn/N )/
2 (the reflected
mode) [22]. The annihilation of a photon in each of
the modes transmitted by the beam splitters witnesses
the overall application of the operator ÂN . If one ob-
serves both reflected and transmitted modes simultane-
ously, one conditions on detection events in all the trans-
mitted modes and no detection events in all the reflected
modes. If detection events are instead observed in all the
reflected modes and in none of the transmitted modes,
an operator of the form (2) is also obtained, but θ is ef-
fectively transformed into θ + π due to the phase shift
at the polarizing beam splitter, and the value of φ of the
generated NOON states is changed by Nπ (see below).
The success probability is thus increased by a factor of
two if both outcomes are accepted.
FIG. 1: Experimental setup for NOON state generation.
OPO, optical parametric oscillator; PBS, polarizing beam
splitter; PS, phase shifter; and APD, avalanche photo diode.
The part of the setup enclosed in the dashed box performs the
NOON state measurement, and here it is shown for N = 3.
Note that the polarizing beam splitters inside the box are ori-
ented at 45◦. The numbers denote beam splitter reflectivities
of 1/3 and 1/2, and the three phase shifters transform b̂+
into −b̂+e
2πin/3, where n = 1, 2, 3, respectively. See text for
details.
For even values of N a similar measurement scheme
is applicable, but it is sufficient to divide the field into
N/2 spatial modes initially, and in this case the NOON
state generation is conditioned on detection events in
both transmitted and reflected modes (see [21]).
III. PERFORMANCE OF THE PROTOCOL
After this presentation of the basic idea and the phys-
ical setup we now consider the actual outcome of the
detection process. For short pulse OPO output the dead
time of the photo detectors may typically be longer than
the pulse duration, and we shall thus assume that it is
impossible to obtain more than a single detection event
per detector per pulse, i.e., if the detector efficiency is
unity, the detectors are only able to distinguish between
vacuum and states different from the vacuum state. Such
detectors are denoted on/off detectors, and they are dis-
cussed in detail in Ref. [23]. The finite dead time of the
detectors is not severe to the measurement procedure de-
scribed in [21] because the on/off detector model and the
conventional photo detector model, represented by the
annihilation operator, lead to identical signal states if the
total number of photons in the idler modes is guaranteed
to be less than or equal to the number of conditioning
detection events, i.e., N .
We analyze the performance of the setup using Gaus-
sian Wigner function formalism [15, 19, 20], which is
applicable because the squeezed states generated by the
OPOs and the vacuum states coupled into the system
via the beam splitters are all Gaussian. In general, the
Wigner function of an n-mode Gaussian state with zero
mean field amplitude takes the form
WV (x1, p1, . . . , xn, pn) =
det(V )
TV −1y, (3)
where y ≡ (x1, p1, . . . , xn, pn)T and V is the 2n × 2n
covariance matrix. If ĉi denotes the field mode annihi-
lation operator of mode i, the elements of V are given
in terms of the real and imaginary parts of the expecta-
tion values 〈ĉ†i ĉj〉 and 〈ĉiĉj〉. Note that for a multi-mode
Gaussian state we are free to include only the modes of
interest in (3) because the partial trace operation over un-
observed modes is equivalent to integration over the cor-
responding quadrature variables. A unit efficiency ‘on’
detection in mode i projects mode i on the subspace of
Hilbert space that is orthogonal to the vacuum state, i.e.,
the Wigner function is multiplied by (1− 2πW0(xi, pi)),
whereW0(x, p) = exp(−x2−p2)/π is the Wigner function
of the vacuum state, the variables xi and pi are integrated
out, and the state is renormalized. Since the Gaussian
nature of a state is preserved under linear transforma-
tions, and since a detector with single-photon efficiency
η is equivalent to a beam splitter with transmission η fol-
lowed by a unit efficiency detector [23], effects of non-unit
detector efficiency are easily included in the covariance
matrix.
To calculate 〈ĉ†i ĉj〉 and 〈ĉiĉj〉 explicitly we note that
the state generated by the OPOs is [24]
|ψi〉 = (1− r2)
rn+m|n, n,m,m〉, (4)
where r is the squeezing parameter and the modes are
listed in the order: â+, â−, b̂+, b̂−. We assume that N is
odd and consider the transmitted trigger modes (which
we number from 1 to N), the â+ mode (mode N + 1),
and the b̂− mode (mode N + 2). By expressing the field
operators of the trigger modes (those observed by the
unit efficiency detectors) in terms of â−, b̂+, and field
operators representing vacuum states we find
〈ĉ†j ĉk〉 = 〈ψi|
− − e−2πij/N b̂
(â− − e2πik/N b̂+eiθ)|ψi〉
1 + e2πi(k−j)/N
, (5)
where j ∈ {1, 2, . . . , N}, k ∈ {1, 2, . . . , N}, λ ≡ ηr2/(1−
r2), and we allow of a constant phase shift θ of b̂+ relative
to â−. Furthermore
〈ĉ†N+1ĉN+1〉 = 〈ψi|â
+â+|ψi〉 = r2/(1− r2), (6)
〈ĉ†N+2ĉN+2〉 = 〈ψi|b̂
−b̂−|ψi〉 = r2/(1− r2), (7)
〈ĉk ĉN+1〉 = 〈ψi|
(â− − e2πik/N+iθ b̂+)â+|ψi〉
1− r2
, (8)
〈ĉk ĉN+2〉 = 〈ψi|
(â− − e2πik/N+iθ b̂+)b̂−|ψi〉
1− r2
e2πik/N+iθ , (9)
〈ĉj ĉk〉 = 〈ĉN+1ĉN+1〉 = 〈ĉN+2ĉN+2〉 = 〈ĉN+1ĉN+2〉 =
〈ĉ†N+1ĉN+2〉 = 〈ĉ
k ĉN+1〉 = 〈ĉ
k ĉN+2〉 = 0. (10)
For even values ofN the factors
η/(2N) are replaced by
η/N . Note that loss in the signal beam may be taken
into account by performing the transformations â+ →√
ηsâ+ and b̂− →
ηsb̂− in the above expressions, where
1− ηs is the loss.
The NOON state fidelity FN of the signal state con-
ditioned on N photo detection events in the transmitted
trigger modes is
WNOON(xN+1, pN+1, xN+2, pN+2)
(1− 2πW0(xi, pi))
WV (x1, p1, . . . , xN+2, pN+2)
dxidpi
, (11)
whereWNOON is the Wigner function of the NOON state
(1), and
(1− 2πW0(xi, pi))
WV (x1, p1, . . . , xN+2, pN+2)
dxidpi
, (12)
is the success probability, i.e., the probability to obtain
the conditioning detection events and produce the NOON
state in a given pulse of the OPO system. We expand
the product
(1− 2πW0(xi, pi)) =
(−2πW0(xi, pi))di ,
where the sum is over all d ≡ (d1, d2, . . . , dN ) with
di ∈ {0, 1}, and define the diagonal matrix Jd =
diag(d1, d1, d2, d2, . . . , dN , dN ) and the n×n identity ma-
trix In. Furthermore, we divide the covariance matrix
into four parts
Vtt Vts
V Tts Vss
, (14)
where Vtt is the 2N×2N covariance matrix of the trigger
modes, Vss is the 4 × 4 covariance matrix of the signal
modes, while Vts contains the correlations between the
trigger and the signal modes, and we define the vector
ys = (xN+1, pN+1, xN+2, pN+2)
T and the matrix
Ud = Vss − V Tts Jd(JdVttJd + I2N )−1JdVts. (15)
This allows us to write Eqs. (11) and (12) in the following
compact forms [25]
det(I2N + JdVtt)
WNOON(ys)WUd(ys)dys, (16)
det(I2N + JdVtt)
. (17)
Since WNOON is a product of a polynomial and a Gaus-
sian the integral in Eq. (16) may be evaluated analytically
and for N = 3 and η = 1 we find
(1− r2)2(2− r2)2(3− 2r2)(6 − 5r2)
18(4− 3r2)
, (18)
where the optimal value φ = Nθ + π + 2πn, n ∈ Z,
is assumed. Expressions for PN are given in table I for
N = 1, 2, 3, and 4, and FN and PN are plotted for N = 3
in Figs. 2 and 3, respectively. We observe that high prob-
abilities are only found in the parameter regime, where
the fidelity is low. If, for instance, we want a NOON
state fidelity of at least 0.9, we choose r = 0.14, and if
η = 0.25, P3 is of order 10
−8. With a repetition rate
of order 106 s−1 (see [16]) one state is produced every
second minute on average. The production rate is very
dependent on detector efficiency, and if η is increased to
unity, the rate is increased by approximately a factor of
For odd values of N we may observe both reflected
and transmitted trigger modes and condition on detec-
tion events in all the transmitted trigger modes and no
detection events in all the reflected trigger modes, or, vice
versa. In this case we also include the reflected trigger
modes in the covariance matrix, which we now denote by
V +. By a similar analysis as above we obtain the success
probability
P+N = 2
2N (−2)
det(I4N + J
, (19)
N PN P
(λ+1)2
(λ+1)2
λ3(λ+4)
(λ+2)2(λ+3)(λ+6)
2λ3(3λ+4)
(λ+1)2(λ+2)2(2λ+3)(5λ+6)
λ4(λ2+6λ+6)
(λ+1)2(λ+2)2(λ2+8λ+8)
TABLE I: Success probabilities calculated from Eqs. (17) and
(19). λ ≡ ηr2/(1− r2).
0 0.2 0.4 0.6 0.8 1
FIG. 2: NOON state fidelity F3 (solid lines) and F
3 (dashed
lines) as a function of squeezing parameter r for η = 1 (upper
lines), η = 0.25 (middle lines), and η → 0 (lower lines). Note
that in the latter case F3 = F
where J+d ≡ diag(d1, d1, . . . , dN , dN , 1, 1, . . . , 1, 1), while
the NOON state fidelity F+N is given by Eq. (11) with V
replaced by the matrix
V − (V +R )
T (V +RR + I2N )
−1V +R , (20)
where V +RR is the covariance matrix of the reflected trig-
ger modes and V +R consists of the correlations between
the reflected trigger modes and the signal and transmit-
ted trigger modes. Explicit results for P+N are given in
table I for N = 1 and 3. F+3 and P
3 are compared to
F3 and P3 in Figs. 2 and 3, and it is observed that F
and P+3 are both larger than F3 and P3 if r is not large
(and η > 0). For r → 1, P+3 → 0 because in this limit it
is very unlikely to obtain no detection events in all the
reflected or in all the transmitted trigger modes.
In the limit of very small detector efficiency a simple
expression for the NOON state fidelity for the case of N
trigger detectors is easily derived without using Wigner
function formalism. In general, if the state of interest
is expressed in the photon number basis, the mathemat-
ical operation corresponding to an ‘on’ detection is to
multiply each term by
1− (1 − η)n, where n is the
number of photons in the mode observed by the non-
unit efficiency detector, trace out the detected mode,
and renormalize. If nη ≪ 1 for all contributing terms,
1− (1− η)n ≈ √nη ∝
n, and the on/off detector
0 0.2 0.4 0.6 0.8 1
FIG. 3: Success probability P3 (solid lines) and P
3 (dashed
lines) as a function of squeezing for η = 1 (upper lines) and
η = 0.25 (lower lines). The dotted lines represent the approx-
imate expression (23).
model becomes equivalent to the photo detector model.
In this case the density operator of the output state is
obtained as
ρ = M
〈p|〈q|(âN− − (b̂+eiθ)N )|ψi〉
〈ψi|((â†−)N − (b̂
−iθ)N )|q〉|p〉
(1− r2)N+2
2N !r2N
(r2)n+m
(n−N)! |n,m〉〈n,m|
−e−iNθ
n!(m+N)!
(n−N)!m!
|n,m〉〈n−N,m+N |
(r2)n+m
(m−N)!
|n,m〉〈n,m|
−eiNθ
(n+N)!m!
n!(m−N)!
|n,m〉〈n+N,m−N |
where M is a normalization constant and the traces are
over the â− and b̂+ modes. This leads to the NOON state
fidelity
N = 〈NOON|ρ|NOON〉 = (1− r
2)N+2, (22)
where again φ = Nθ + π + 2πn, n ∈ Z, is assumed.
It is interesting to compare this result with the fidelity
(1 − r2)N+1 obtained for production of N -photon states
from a single two-mode squeezed state by conditioning
on N detection events in the idler beam and using detec-
tors with very small efficiency. If a single-photon state is
produced by this method and transformed into an N = 1
NOON state as explained in the Introduction, the NOON
state fidelity is F1,s = (1−r2)2, and the success probabil-
ity is P1,s = λ/(λ + 1). Choosing squeezing parameters
such that F1,s = F1, we find that P
1 = (4/3)P1,s in the
high fidelity limit. It is thus possible to achieve a higher
success probability using the scheme with two OPOs, but
the price to pay is a more technically involved setup,
and NOON states with two different values of φ are pro-
duced. For N = 2 the present protocol and combination
of two single-photon states on a 50 : 50 beam splitter,
each produced conditionally from a single OPO, lead to
identical fidelities and success probabilities. Finally we
note that the photo detector model underestimates FN
for η > 0 because 1 − (1 − η)n = η
i=0 (1 − η)i < nη
for n = 2, 3, . . . while 1− (1− η)n = nη for n = 0, 1, i.e.,
the ‘wrong’ terms containing more than N photons are
given a too large weight. This is also what we observe in
Fig. 2.
In the limit of small r and for odd values of N the
success probability is given approximately by the simple
expression
〈ψi|((â†−)N − (b̂
−iθ)N )
(âN− − (b̂+eiθ)N )|ψi〉 =
(2N)N
λN (N odd). (23)
Again η/(2N) must be replaced by η/N to obtain PN for
even values of N . The approximation to P3 is shown in
Fig. 3.
IV. NOON STATES FROM
CONTINUOUS-WAVE OPO SOURCES
Our protocol is not limited to pulsed fields, and for
completeness we now briefly consider NOON state gen-
eration from continuously driven OPOs. We assume
N = 3. For continuous-wave fields each of the three
detected trigger beams and the two signal beams are
described by time dependent field operators ĉi(t). The
trigger detections take place in particular modes local-
ized around the three detection times tc1, tc2, and tc3,
and we want to determine the NOON state fidelity of
an output state occupying one temporal mode in each
signal beam. Following the general multimode formal-
ism in Refs. [15, 20], the five relevant modes are spec-
ified by the mode functions fi(t), and the correspond-
ing five single mode operators are ĉi =
fi(t)ĉi(t)dt.
In general, we are free to choose the two output mode
functions at will, and in the present case it is natural to
choose the mode function which gives rise to the largest
three-photon state fidelity when three-photon states are
generated from a single type II continuous-wave OPO.
Since we are mainly interested in the parameter region
where the squeezing is small and the NOON state fidelity
is large, we use the optimal three-photon state mode
function derived for very low beam intensity in [20], i.e.,
f4(t) = f5(t) =
k=1 ck
γ/2 exp(−γ|t− tck|/2), where
0 0.05 0.1 0.15
FIG. 4: NOON state fidelity as a function of ǫ/γ for states
generated from a pair of continuous-wave OPO sources when
conditioning on three simultaneous trigger detection events
tc1 = tc2 = tc3.
0 0.5 1 1.5 2
FIG. 5: Fidelity of NOON states from continuous-wave OPO
sources as a function of separation between trigger detection
events (tc3 − tc1)γ for N = 3, tc3 − tc2 = tc2 − tc1, and ǫ/γ =
0.01.
the coefficients ck are functions of the intervals between
the detection times and γ is the leakage rate of the OPO
output mirror. We furthermore assume that the trigger
mode functions are nonzero only in an infinitesimal time
interval centered at the detection time, which is valid if
the trigger detections take place on a time scale much
shorter than γ−1. Since we consider a low intensity con-
tinuous beam, and since we formally assume that the
trigger detectors only register the light field in infinitesi-
mal time intervals around the detection times, the anni-
hilation operator detector model is perfectly valid in this
case and detector dead time is insignificant.
We may now proceed as above and eliminate all the
irrelevant modes from the analysis by writing down the
Gaussian Wigner function of the five interesting modes.
The only difference is that this time 〈ĉ†i ĉj〉 and 〈ĉiĉj〉 are
expressed in terms of the two time correlation functions
of the OPO output. Also, the operators applied to the
Wigner function to take conditioning into account are
different because the annihilation detector model is used.
The reader is referred to Refs. [15, 20] for details.
The resulting fidelity is shown as a function of ǫ/γ in
Fig. 4, where ǫ is the nonlinear gain in the OPO, and as a
function of the temporal distance between the condition-
ing detection events in Fig. 5. As in the pulsed case the
fidelity decreases when the degree of squeezing increases.
The fidelity also decreases when the temporal distance
between the conditioning detection events increases from
zero, but it is permissible to have a small time interval
between the trigger detection events. We note that the
curves represent a lower limit to the theoretically achiev-
able fidelity since a better fidelity may be obtained for
another choice of output state mode functions.
V. CONCLUSION
In conclusion we have analyzed a method to gener-
ate path entangled NOON states from the output from
two optical parametric oscillators. The method relies
on the joint detection of photons in a number of trig-
ger beams, and we presented a theoretical analysis of the
role of detector efficiency and dead time for the fidelity
of the states obtained and the success probability of the
protocol. Our specific NOON state protocol applies for
general photon numbers of the states, but in practice it
is not realistic to go beyond the case of N = 3, studied
here. This is due to the reduction of the fidelity due to
unwanted contributions from higher number states, when
the OPO output power gets too high, combined with the
severe reduction of the probability to obtain the number
of conditioning detection events needed when the OPO
output power is too low. The N = 3 NOON states, which
can be produced at 90% fidelity at the rate of one state
produced every 10− 100 seconds, seem to be at the limit
of realistic experiments of the proposed kind. Finally, we
also determined the NOON state fidelity for continuous-
wave fields, where the best NOON state occupies a pair
of suitably selected temporal mode functions, and where
we find high fidelities as long as the trigger events occur
within a short time window compared to the lifetime of
the OPO cavity field.
We presented this analysis for the production of op-
tical NOON states, but we note that recent theoreti-
cal proposals and experiments with four wave mixing
of matter waves [26], engineered quadratic interactions
among trapped ions [27], and entanglement between field
and atomic degrees of freedom [28, 29] bring promise for
similar conditional generation of atomic and interspecies
atom-field NOON states.
This work was supported by the European Integrated
project SCALA.
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|
0704.0398 | Renewals for exponentially increasing lifetimes, with an application to
digital search trees | Renewals for exponentially increasing lifetimes, with an application to digital search trees
The Annals of Applied Probability
2007, Vol. 17, No. 2, 676–687
DOI: 10.1214/105051606000000862
c© Institute of Mathematical Statistics, 2007
RENEWALS FOR EXPONENTIALLY INCREASING LIFETIMES,
WITH AN APPLICATION TO DIGITAL SEARCH TREES
By Florian Dennert and Rudolf Grübel
Universität Hannover
We show that the number of renewals up to time t exhibits distri-
butional fluctuations as t→∞ if the underlying lifetimes increase at
an exponential rate in a distributional sense. This provides a proba-
bilistic explanation for the asymptotics of insertion depth in random
trees generated by a bit-comparison strategy from uniform input;
we also obtain a representation for the resulting family of limit laws
along subsequences. Our approach can also be used to obtain rates
of convergence.
1. Introduction. Let (Yk)k∈N be a sequence of independent, nonnegative
random variables and let (Sn)n∈N0 ,
S0 := 0, Sn :=
Yk for all n ∈N,
be the associated sequence of partial sums. Regarding the Yk’s as successive
lifetimes and Sn as the time of the nth renewal, we interpret
Nt := sup{n ∈N0 :Sn ≤ t}
as the number of renewals up to and including time t; (Nt)t≥0 is the renewal
process. Standard renewal theory assumes that the Yk’s all have the same
distribution, in which case Nt, appropriately rescaled, is asymptotically nor-
mal as t→∞. For this result, and for renewal theory in general, we refer
the reader to Section XI in [3].
In this note we consider exponentially increasing lifetimes. We show that
in such a case the distribution of Nt does not converge and that asymp-
totic distributional fluctuations appear (Section 2). Such fluctuations occur
frequently in the analysis of algorithms. The renewal theoretic framework
Received January 2006.
AMS 2000 subject classifications. Primary 60K05; secondary 68Q25.
Key words and phrases. Asymptotic distributional behavior, limiting periodicities, re-
newal processes.
This is an electronic reprint of the original article published by the
Institute of Mathematical Statistics in The Annals of Applied Probability,
2007, Vol. 17, No. 2, 676–687. This reprint differs from the original in pagination
and typographic detail.
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2 F. DENNERT AND R. GRÜBEL
provides a probabilistic view of this phenomenon in connection with digital
search trees (Section 3). We also indicate how our approach can be used to
obtain rates of convergence (Section 4).
2. Renewals for increasing lifetimes. We assume that the lifetimes in-
crease exponentially with rate α, where α> 1 is fixed throughout the sequel,
in the sense that
α−kYk →distr Y∞ and α−kEYk →EY∞(1)
for some random variable Y∞ and as k→∞. Here “→distr” denotes conver-
gence in distribution, so that the first part of (1) means that
Ef(α−kYk) =Ef(Y∞)
for all bounded continuous functions f :R→ R. Below we will use the fact
that in order to prove Xn →distr X it is sufficient to show that Ef(Xn)→
Ef(X) holds for all bounded and uniformly continuous functions. For details
and a general treatment of convergence in distribution we refer the reader
to [1]. Of course, only the distribution µ = L(Y∞) of Y∞ matters, so we
will occasionally write α−kYk →distr µ instead. Finally, throughout this note
a condition involving moments is meant to imply that these moments are
finite.
An important role will be played by
S∞ :=
α−kY∞,k,
where (Y∞,k)k∈N0 is a sequence of independent and identically distributed
random variables with L(Y∞,0) = L(Y∞), Y∞ as in (1). From EY∞ <∞ we
obtain ES∞ = α(α − 1)−1EY∞ <∞ and therefore P (S∞ <∞) = 1; more-
over, we then also have that
−kY∞,k converges almost surely and
hence in distribution to S∞ as n→∞. We will also assume that L(Y∞) has
no atoms, that is,
P (Y∞ = y) = 0 for all y ∈R+.(2)
Finally, it is an elementary analytic fact that, for a sequence (xn)n∈N of real
numbers with limit x ∈R,
α−kxn−k =
The following lemma can be regarded as a random version of (3).
Lemma 1. If (1) and (2) are satisfied, then α−nSn →distr S∞ as n→∞,
and P (S∞ = y) = 0 for all y ∈R.
RENEWALS FOR INCREASING LIFETIMES 3
Proof. Suppose that (Uk)k∈N is a sequence of independent random
variables on some probability space (Ω,A, P ), all uniformly distributed on
the unit interval. Let Fk be the distribution function of Yk, F the distribution
function of Y∞. We use a variant of the quantile construction:
Ỹk := F
k (Uk), Ỹ∞,k := F
−1(Uk) for all k ∈N.
We then have L(Ỹ1, . . . , Ỹn) = L(Y1, . . . , Yn) for all n ∈N, which implies
L(α−nSn) =L(α−nS̃n) with S̃n :=
Using α−nS̃n =
k=0 α
−k(α−(n−k)Ỹn−k) we obtain
α−nS̃n −
α−kỸ∞,n−k
α−kE|α−(n−k)Ỹn−k − Ỹ∞,n−k|.(4)
With Y ′k := F
k (U1) and Y
∞ := F
−1(U1) we have
E|α−kỸk − Ỹ∞,k|=E|α−kY ′k − Y ′∞|.(5)
From (1) it follows that α−kY ′k →distr Y ′∞ and Eα−kY ′k → EY ′∞. Because
of Y ′k, Y
∞ ≥ 0 Theorem 5.4 in [1] applies and gives the L1-convergence of
α−kY ′k to Y
∞, that is, E|α−kY ′k − Y ′∞| → 0 as k →∞. Using this together
with (3), (4) and (5) we obtain
α−nS̃n −
α−kỸ∞,n−k
= 0.(6)
Now let f :R→R be bounded and uniformly continuous. We have
|Ef(α−nSn)−Ef(S∞)|=
Ef(α−nS̃n)−Ef
α−kỸ∞,n−k
α−kỸ∞,k
α−kỸ∞,k
f(α−nS̃n)− f
α−kỸ∞,n−k
α−kỸ∞,k
α−kỸ∞,k
For the first integral on the right-hand side we use (6), for the second an
elementary estimate shows that the difference between the arguments of f
converges to 0 in probability. In both cases we now use uniform continuity
4 F. DENNERT AND R. GRÜBEL
when the arguments of f are close to each other and boundedness otherwise.
This leads to
Ef(α−nSn) =Ef(S∞),
which gives the convergence in distribution. The statement about the atoms
of S∞ follows from (2) and the fact that S∞ is equal in distribution to
Y∞ + α
−1S∞ with Y∞ and S∞ independent. �
The above proof is based on classical weak convergence arguments. An
alternative proof can be obtained via the Wasserstein distance
dW (µ, ν) = inf{E|X − Y | :L(X) = µ,L(Y ) = ν},
its known relation to weak convergence and convergence of the first moments,
and the same variant of the quantile construction, which in this context is
known as the comonotone coupling.
We write ⌊x⌋ for the greatest integer less than or equal to x and {x} for
the fractional part of x ∈R.
Theorem 2. Suppose that (1) and (2) are satisfied and let
Qη := L(⌊− logαS∞ + η⌋), 0≤ η ≤ 1.(7)
If (tn)n∈N is a sequence of real numbers with tn →∞ and such that {logα tn}→
η for some η ∈ [0,1], then
Ntn − ⌊logα tn⌋→distr Qη as n→∞.
Proof. We use the abbreviations kn := ⌊logα tn⌋ and ηn := {logα tn}. In
particular, logα tn = kn + ηn. Further, let Z∞ := − logαS∞. By a standard
renewal theoretic argument,
P (Nt = j) = P (Sj ≤ t)−P (Sj+1 ≤ t) for all t≥ 0, j ∈N0,
hence
P (Ntn − kn = j) = P (Skn+j ≤ tn)− P (Skn+j+1 ≤ tn)
= P (− logα(α−kn−jSkn+j) + ηn ≥ j)
−P (− logα(α−kn−j−1Skn+j+1) + ηn ≥ j + 1)
→ P (⌊Z∞ + η⌋= j) as n→∞,
where in the last step Lemma 1 and three general facts about convergence
in distribution were used: First, the continuous mapping theorem, which
implies that − logα(α−mSm)→distr − logαS∞ as m→∞; secondly, the in-
terplay with convergence in probability, see Theorem 4.1 in [1], which yields
RENEWALS FOR INCREASING LIFETIMES 5
− logα(α−nSn)+ ηn →distr − logαS∞+ η as n→∞; finally, that L(S∞) and
therefore also L(− logαS∞+η) assign probability 0 to single points and that
this implies
P (− logα(α−nSn) + ηn ≥ z) = P (− logαS∞ + η ≥ z) for all z ∈R.
A structural consequence of the representation (7) is the →distr-continuity
of η 7→Qη on the open unit interval; at η = 0 this function is right continuous,
at η = 1 it is left continuous. The extreme members are translates of each
other in the sense that Q0({j}) =Q1({j + 1}) for all j ∈ Z.
The total variation distance dTV of probability measures is defined by
dTV(µ, ν) := sup
|µ(B)− ν(B)|,
for µ, ν concentrated on Z this can be written as
dTV(µ, ν) =
|µ({j})− ν({j})|.(8)
For a sequence of probability measures that are concentrated on a fixed
countable set Scheffé’s lemma implies that weak convergence is equivalent
to convergence in total variation distance, hence (7) can be rewritten as
dTV(L(Ntn − ⌊logα tn⌋),Q{log
tn}) = 0.
Because of the continuity of [0,1] ∋ η 7→Qη this in turn leads to a statement
that avoids the use of subsequences,
dTV(L(Nt − ⌊logα t⌋),Q{log
t}) = 0.(9)
In Section 4 we will investigate the rate of convergence in (9) in a particular
case.
3. An application to digital search trees. The nodes of a (rooted, di-
rected) binary tree can be represented by finite strings of 0’s and 1’s if we
interpret 0 as a move to the left and 1 as a move to the right. The length of
the string is the depth (or level) of the node it represents, the root node corre-
sponds to the empty string and has level 0. The sequence (Tn)n∈N associated
with a sequence (xn)n∈N of numbers from the unit interval by the DST (dig-
ital search tree) algorithm is obtained as follows: For T1, we put x1 into the
root node. If x1, . . . , xn have been stored into Tn then the position of xn+1 is
determined by traveling through the tree with the direction given by the bi-
nary expansion of xn+1 until an empty node has been found. This algorithm
and its properties are discussed in the standard texts of the area, for exam-
ple, [8, 10, 11]. As an example we consider the first ten numbers given in [8],
6 F. DENNERT AND R. GRÜBEL
Fig. 1. Binary tree.
Appendix A, (
2, log 2, log 3, log 10). Let xi be the
fractional part of the ith entry, 1≤ i≤ 10; the relevant first four bits of the
respective binary expansions are given by (0110,1011,0011,0010, 0100,0111,
0011,1011,0001,0100). This leads to the binary tree given in Figure 1.
Consider now the sequence (Tn)n∈N0 of random trees that the DST algo-
rithm associates with a sequence (Un)n∈N of independent random variables,
where we assume that the Un’s are uniformly distributed on the unit inter-
val and that T0 is the empty tree. Let Xn(θ) be the depth of the first free
node of Tn along the path determined by a sequence θ ∈ {0,1}N. Such a θ
defines a family of nested intervals of length 2−k, k = 1,2,3, . . . , and it is
easy to see that (Xn(θ))n∈N0 is a Markov chain with X0(θ) ≡ 0 and tran-
sition probabilities pk,k+1 = 1− pk,k = 2−k for all k ∈ N0. Conditioning on
the value of Un+1 we see that the distribution of Xn(θ) is the same as the
distribution of Zn+1, the insertion depth of Un+1. This quantity is known
as “unsuccessful search” in the literature on the analysis of algorithms. [Of
course, this distributional equality does not hold for the joint distributions:
n 7→Xn(θ) is increasing, n 7→ Zn+1 is not.] For example, the next number
in Knuth’s list is x11 = 1/ log 2, the binary expansion of the fractional part
{x11} begins with 011100 and hence x11 would be inserted at level 4 as the
right child of x6.
The Markov chain (Xn(θ))n∈N0 is of the simple birth type and can there-
fore be described by its respective holding times Y1, Y2, Y3, . . . in the states
k = 0,1,2, . . . . These are independent, and Yk has a geometric distribution
with parameter pk−1,k, that is, for all k ∈N,
P (Yk = j) = (1− 2−k+1)j−12−k+1 for all j ∈N.
Here we interpret the case k = 1 as Y1, the holding time in 0, being constant
and equal to 1. As a result of its simple stochastic dynamics, (Xn(θ))n∈N0
is equal to the renewal process N associated with the sequence (Yk)k∈N,
observed at discrete times, that is, (Xn(θ))n∈N0 = (Nn)n∈N0 . It is easy to
see that for this sequence (Yk)k∈N of lifetimes conditions (1) and (2) are
satisfied and that L(Y∞) = Exp(2), with Exp(λ) the exponential distribution
RENEWALS FOR INCREASING LIFETIMES 7
with parameter λ (and mean 1/λ). Hence Theorem 2 can be applied: If the
sequence (n(m))m∈N ⊂ N is such that n(m) →∞ and {log2 n(m)} → η as
m→∞, then
Xn(m)(θ)− ⌊log2 n(m)⌋→distr Qη.(10)
Here Qη , 0≤ η ≤ 1, is the distribution of ⌊− log2 S∞+η⌋, S∞ :=
k=0 2
−kY∞,k
and Y∞,k, k ∈N0, are independent and identically distributed with L(Y∞,1) =
Exp(2). Alternatively, we can write S∞ :=
k=1 Ỹk with Ỹk, k ∈N, again in-
dependent and L(Ỹk) = Exp(2k) for all k ∈N.
The explicit representation of the family of limit distributions on the basis
of the convolution product of the distributions Exp(2k), k ∈N, can be used
to obtain a series expansion for the distribution functions associated with
Qη , 0 ≤ η ≤ 1. For this, we start with a partial fraction expansion: For all
n ∈N and all z ∈C with |ℜ(z)|< 2,
(1− 2−kz)−1 =
an,k(1− 2−kz)−1,(11)
where an,k :=
j=1(1− 2j)−1
j=1 (1− 2−j)−1. Reading (11) as an equality
relating characteristic functions we obtain
Exp(21) ⋆Exp(22) ⋆ · · · ⋆Exp(2n) =
an,kExp(2
k).(12)
Note, however, that the right-hand side in (12) is not the usual mixture of
probability distributions as the coefficients alternate in sign. With
ak := b
(1− 2j)−1, b :=
(1− 2−j)−1,
letting n→∞ in (12) leads to L(S∞) =
k=1 akExp(2
k), so that
Qη((−∞, x]) = P (⌊− log2(S∞) + η⌋ ≤ x)
= P (S∞ > 2
η−1−x)(13)
ak exp(−2k+η−1−x) for all x ∈ Z.
This representation of the limiting distribution functions as an alternating
series has already been obtained by Louchard [9] in the context of digital
search trees and by Flajolet [4] in the context of approximate counting; see
also Section 6.4 in [10] and Section 6.3 in [8] for related results. These authors
use a completely different approach, more analytic in flavor and relying on
combinatorial identities due to Euler.
8 F. DENNERT AND R. GRÜBEL
Our main point here, however, is not a rederivation of (13) but the rep-
resentation of the family {Qη : 0≤ η < 1} in terms of one particular random
variable, which is first shifted by η and then discretized. This representation
can, for example, be used to obtain information about the tail behavior of
the limit distributions. Janson [7] notes that (13) by itself would only give
an exponential rate of decrease for the tail probabilities, he then provides an
analytic argument that improves this to a superexponential rate by show-
ing that the associated Fourier transform is an entire function. Using the
representation S∞ =
k=1 2
−kZk with Zk independent and L(Zk) = Exp(1)
together with the fact that Exp(1) has a density bounded by 1, we get
P (S∞ ≤ 2−j)≤ P (Z1 ≤ 2−j+1)P (Z2 ≤ 2−j+2) · · ·P (Zj−1 ≤ 2−1)
≤ 2−j+12−j+2 · · ·2−1
= 2−j(j−1)/2
for all j ∈ N. Because of Qη([k,∞)) ≤ P (S∞ ≤ 2−k+1) for all k ∈ N, k ≥ 2,
this leads to
Qη([x,∞)) = o(exp(−ρx2)) as x→∞, for all ρ < (log 2)/2.
The fact that a representation by discretization is possible in many situ-
ations where fluctuations were first found by calculation seems to belong to
the folklore of the subject, at least in simple instances such as the asymp-
totic distributional behavior of the maximum of a sample from a geometric
distribution. The geometric case together with some renewal theoretic tech-
niques (for identically distributed lifetimes) was used in [5] to obtain results
of the above type for von Neumann addition. In [2] a discretization represen-
tation occurs on the level of stochastic processes, leading to a probabilistic
approach to fluctuation phenomena in the context of multiplicities of the
maximum in a random sample from a discrete distribution. In a recent pa-
per, Janson [7] studies the effects of discretizing random variables and the
resulting distributional fluctuations and gives a range of interesting exam-
ples. Of course, the explanation for periodicities can be, and indeed often is,
quite different and mechanisms other than discretization may be responsible;
see, for example, [6] and the references given there.
4. Rates of convergence. The renewal theoretic approach can also be
used to obtain rates of convergence. We sketch one of the possibilities, for
a particular choice of distances, and give details for the DST situation from
Section 3. Let, for t > 0, k(t) := ⌊logα t⌋ and η(t) := {logα t}.
The Kolmogorov–Smirnov distance of two probability measures µ and ν
on the real line is defined by
dKS(µ, ν) := sup
|µ((−∞, x])− ν((−∞, x])|.
RENEWALS FOR INCREASING LIFETIMES 9
If X and Y are real random variables, then we abbreviate dKS(L(X),L(Y ))
to dKS(X,Y ); if F and G are the associated distribution functions, then
dKS(X,Y ) = ‖F − G‖∞, where the supremum norm for general bounded
functions f :R → R is given by ‖f‖∞ := supx∈R |f(x)|. The Kolmogorov–
Smirnov distance is obviously invariant under strictly monotone transfor-
mations. For example,
dKS(αX + β,αY + β) = dKS(X,Y ) for all α,β ∈R, α 6= 0,
and for X,Y > 0,
dKS(X,Y ) = dKS(logX, logY ).
With the notation as in the proof of Theorem 2,
|P (Nt − k(t) = j)−P (⌊− logα(S∞) + η(t)⌋= j)|
≤ |P (− logα(α−k(t)−jSk(t)+j) + η(t)≥ j)−P (− logα(S∞) + η(t)≥ j)|
+ |P (− logα(α−k(t)−j−1Sk(t)+j+1) + η(t)≥ j +1)
− P (− logα(S∞) + η(t)≥ j + 1)|.
With the auxiliary quantities
Zt := ⌊− logα(S∞) + η(t)⌋, φ(m) := dKS(α−mSm, S∞)
and the above properties of the Kolmogorov–Smirnov distance this leads to
|P (Nt − k(t) = j)−P (Zt = j)| ≤ φ(k(t) + j) + φ(k(t) + j +1).(14)
It is often possible to obtain an upper bound for negative j, say j ≤−k(t)/2,
directly. In such cases the above elementary renewal theoretic argument leads
to a bound for the ‖ · ‖∞-distance between the probability mass functions of
Nt − k(t) and Zt, for example; note that the latter variable has distribution
Qη(t) where Qη , 0≤ η ≤ 1, is the set of limit distributions along subsequences
that appears in Theorem 2.
The above argument covers the step from (α−mSm)m∈N to (Nt)t≥0. How-
ever, in an application the starting point will usually be the convergence of
the scaled lifetimes in (1), which means that we also need an analogue for
Lemma 1 that gives rates of convergence.
We carry this out in the specific context of digital search trees. The fol-
lowing general bounds will turn out to be useful: If X has density fX and if
P (|Y | ≤ c) = 1, then
dKS(X,X + Y )≤ c‖fX‖∞.(15)
Indeed: For all z ∈R, P (X ≤ z− c)≤ P (X+Y ≤ z)≤ P (X ≤ z+ c), so that
|P (X + Y ≤ z)−P (X ≤ z)|
≤max{P (X ≤ z + c)−P (X ≤ z), P (X ≤ z)−P (X ≤ z − c)},
10 F. DENNERT AND R. GRÜBEL
and, of course, P (X ∈ (a, b]) ≤ (b − a)‖fX‖∞. This bound can easily be
generalized to
dKS(X,X + Y )≤ c‖fX‖∞ +P (|Y |> c) for all c > 0,(16)
where we still assume that X has density fX , but Y may be arbitrary.
Note that X and Y need not be independent in (15) and (16). If they are
independent then it is easy to show that
dKS(X,X + Y )≤ ‖fX‖∞E|Y |.(17)
In (17) boundedness of Y is not needed but the bound obviously makes sense
only if Y has finite first moment. Finally, in connection with density bounds
the interplay with convolution is of interest: We have ‖f ⋆ g‖∞ ≤ ‖f‖∞ for
all probability densities f, g. For example, if a sum of independent random
variables contains a summand with distribution Exp(λ), then the density of
the sum is bounded by λ.
Lemma 3. With (Yk)k∈N and S∞ as in Section 3,
dKS(2
−nSn, S∞) =O(n2
Proof. Let (Zk)k∈N be a sequence of independent random variables, all
exponentially distributed with parameter 1. Then S∞ is equal in distribution
k=1 2
−kZk. We recall that the kth lifetime Yk has a geometric distribu-
tion with parameter 2−k+1. On the basis of (Zk)k∈N we define a sequence
(Ỹk)k∈N by Ỹk := ⌊αkZk⌋+1 for all k ∈N, with
α1 := 0, αk := (− log(1− 2−k+1))−1 for k > 1.
It is easy to check that
(Ỹk)k∈N =distr (Yk)k∈N, 2
2k−1Zk =distr
2−kZk.
Hence, with φ(n) denoting the dKS-distance of 2
−nSn and S∞,
φ(n)≤ φ1(n) + φ2(n) + φ3(n) for all n ∈N,
with φ1, φ2, φ3 defined by
φ1(n) := dKS
Ỹk,2
φ2(n) := dKS
αkZk,2
2k−1Zk
φ3(n) := dKS
2−kZk,
2−kZk
RENEWALS FOR INCREASING LIFETIMES 11
For the random variables in φ1 we have
Vn ≤ 2−n
Ỹk ≤ Vn + n2−n with Vn := 2−n
αkZk.
It is easy to show that the densities of Vn, n ∈N, can be uniformly bounded
for all n by some finite constant C1, hence (15) implies that φ1(n)≤C1n2−n
for all n ∈N.
The elementary bounds
− 1≤−
log(1− x)
for 0< x≤
together with α1 = 0 imply supk∈N |αk − 2k−1|= 1, hence we have
αkZk − 2−n
2k−1Zk
≤ 2−n
The familiar combination of Markov’s inequality and moment generating
functions gives
Zk ≥ (1 + κ)n
=O(2−n)
if κ is chosen large enough, so that we can use (16) with c = c(n) = (1 +
κ)n2−n to obtain that φ2(n)≤C2n2−n for all n ∈N, for some finite constant
For φ3 finally we use (17): For the densities of the finite sums we again
have a finite uniform bound for all n, and
k=n+1
2−kZk
k=n+1
2−kEZk = 2
so that φ3(n) ≤ C32−n for all n ∈ N with some C3 <∞. Putting these to-
gether we arrive at
φ(n)≤Cn2−n for all n ∈N
with some finite constant C. �
In the application under consideration we obtain a rate of convergence
result with respect to the total variation distance, which is stronger than
a result for the supremum norm distance of the corresponding probability
mass functions that we mentioned in connection with (14).
Theorem 4. With (Xn(θ))n∈N and Qη as in Section 3,
dTV(L(Xn(θ)− ⌊log2 n⌋),Q{log2 n}) = o(n
−γ) for all γ < 1.
12 F. DENNERT AND R. GRÜBEL
Proof. We use the abbreviations k(n) := ⌊log2 n⌋ and η(n) := {log2 n}.
Let γ < 1 be given and choose ε > 0 such that ε < 1− γ. Lemma 3 together
with (14) gives
j≥−εk(n)
|P (Nn − k(n) = j)−Qη(n)({j})| ≤C
j≥(1−ε)k(n)
for all n ∈ N with some finite constant C. Our choice of ε implies that the
upper bound has the desired rate o(n−γ).
For the remaining part of the infinite sum in (8) we replace the absolute
difference of the probabilities by their sum, which means that it is now
enough to show that
P (Nn ≤ (1− ε)k(n)) = o(n−γ),(18)
P (− log2(S∞)≤−εk(n) + 1) = o(n−γ).(19)
Here we have used that Qη is the distribution of ⌊− log2(S∞)+ η⌋. It is easy
to show that the moment generating function for S∞ exists in a neighbor-
hood of 0, hence
P (S∞ > x) = o(e
−κx) for all x > 0(20)
with some κ > 0. Straightforward manipulations show that (20) implies (19);
indeed, the probability converges faster to 0 than any negative power of n.
Using once again the relation between the number of renewals and the partial
sums of the lifetimes we further obtain, with m(n, ε) := ⌊(1− ε)k(n)⌋,
P (Nn ≤ (1− ε)k(n)) ≤ P (Sm(n,ε) ≥ n)
= P (2−m(n,ε)Sm(n,ε) ≥ n2−m(n,ε))
≤ dKS(2−m(n,ε)Sm(n,ε), S∞) +P (S∞ ≥ n2−m(n,ε)).
For the Kolmogorov–Smirnov distance we use Lemma 3, for the tail of S∞
the desired rate follows with (20). This gives (18) and hence completes the
proof. �
REFERENCES
[1] Billingsley, P. (1968). Convergence of Probability Measures. Wiley, New York.
MR0233396
[2] Bruss, F. Th. and Grübel, R. (2003). On the multiplicity of the maximum in a
discrete random sample. Ann. Appl. Probab. 13 1252–1263. MR2023876
[3] Feller, W. (1971). An Introduction to Probability Theory and Its Applications II,
2nd ed. Wiley, New York. MR0270403
[4] Flajolet, Ph. (1985). Approximate counting: A detailed analysis. BIT 25 113–134.
MR0785808
[5] Grübel, R. and Reimers, A. (2001). On the number of iterations required by von
Neumann addition. Theor. Inform. Appl. 35 187–206. MR1862462
http://www.ams.org/mathscinet-getitem?mr=0233396
http://www.ams.org/mathscinet-getitem?mr=2023876
http://www.ams.org/mathscinet-getitem?mr=0270403
http://www.ams.org/mathscinet-getitem?mr=0785808
http://www.ams.org/mathscinet-getitem?mr=1862462
RENEWALS FOR INCREASING LIFETIMES 13
[6] Janson, S. (2004). Functional limit theorems for multitype branching processes and
generalized Pólya urns. Stochastic Process. Appl. 110 177–245. MR2040966
[7] Janson, S. (2006). Rounding of continuous random variables and oscillatory asymp-
totics. Ann. Probab. 34 1807–1826.
[8] Knuth, D. E. (1973). The Art of Computer Programming 3. Sorting and Searching.
Addison–Wesley, Reading, MA. MR0445948
[9] Louchard, G. (1987). Exact and asymptotic distributions in digital binary search
trees. Theor. Inform. Appl. 21 479–496. MR0928772
[10] Mahmoud, H. M. (1992). Evolution of Random Search Trees. Wiley, New York.
MR1140708
[11] Sedgewick, R. and Flajolet, Ph. (1996). An Introduction to the Analysis of Al-
gorithms. Addison–Wesley, Reading, MA.
Institut für Mathematische Stochastik
Universität Hannover
Postfach 60 09
D-30060 Hannover
Germany
E-mail: [email protected]
[email protected]
http://www.ams.org/mathscinet-getitem?mr=2040966
http://www.ams.org/mathscinet-getitem?mr=0445948
http://www.ams.org/mathscinet-getitem?mr=0928772
http://www.ams.org/mathscinet-getitem?mr=1140708
mailto:[email protected]
mailto:[email protected]
Introduction
Renewals for increasing lifetimes
An application to digital search trees
Rates of convergence
References
Author's addresses
|
0704.0399 | Hawking radiation of linear dilaton black holes | LAPTH-1178/07
Hawking radiation of linear dilaton black holes
G. Clémenta∗, J.C. Fabrisb†and G.T. Marquesa,b‡
aLaboratoire de Physique Théorique LAPTH (CNRS),
B.P.110, F-74941 Annecy-le-Vieux cedex, France
b Departamento de F́ısica, Universidade Federal do Esṕırito Santo,
Vitória, 29060-900, Esṕırito Santo, Brazil
April 3, 2007
Abstract
We compute exactly the semi-classical radiation spectrum for a
class of non-asymptotically flat charged dilaton black holes, the so-
called linear dilaton black holes. In the high frequency regime, the
temperature for these black holes generically agrees with the surface
gravity result. In the special case where the black hole is massless,
we show that, although the surface gravity remains finite, there is
no radiation, in agreement with the fact that massless objects cannot
radiate.
e-mail: [email protected]
e-mail: [email protected]
e-mail:[email protected]
http://arxiv.org/abs/0704.0399v1
Quantum field theory in curved spacetime predicts new phenomena such
as particle emission by a black hole [1]. This is due to the fact that the vac-
uum for a quantum field near the horizon is different from the observer’s
vacuum at spatial infinity. A distant observer thus receives from a black
hole a steady flux of particles exhibiting, in the high frequency regime, a
black body spectrum with a temperature proportional to the surface grav-
ity [2]. Although Hawking’s original derivation of this black hole evaporation
dealt with realistic collapsing black holes, Unruh [3] showed that the same
results are obtained when the collapse is replaced by appropriate boundary
conditions on the horizon of an eternal black hole. In the semi-classical ap-
proximation, the black hole radiation spectrum may be evaluated by com-
puting the Bogoliubov coefficients relating the two vacua. An equivalent
procedure is to compute the reflection and absorption coefficients of a wave
by the black hole. Usually, the wave equation cannot be solved exactly,
and one must resort to match solutions in an overlap region between the
near-horizon and asymptotic regions [4, 5]. In the special case of the (2+1)-
dimensional BTZ black hole [6], an exact solution of the wave equation is
available, which allows for an exact computation of the radiation spectrum,
leading to the Hawking temperature [7, 8, 9].
In this Letter, we discuss another case of black holes also allowing for
an exact semi-classical computation of their radiation spectrum, that of lin-
ear dilaton black hole solutions to Einstein-Maxwell dilaton (EMD) theory
in four dimensions. Linear dilaton black holes are a special case of the
more general class of non-asymptotically flat black hole solutions to EMD
[10, 11], which we first briefly present. We discuss the evaporation of these
non-asymptotically flat black holes and show that they either collapse to a
naked singularity in a finite time, or evaporate in an infinite time. We then
specialize to linear dilaton black holes, and outline the analytical computa-
tion of their radiation spectrum. For massive black holes, this computation
leads, in the high frequency regime, to the same temperature which is ob-
tained from the surface gravity. However in the case of massless extreme
black holes, we find that, although the surface gravity remains finite, there
is no radiation, in agreement with the fact that a massless object cannot
radiate.
EMD is defined by the following action
R− 2∂µφ∂µφ− e−2αφFµνFµν
, (1)
where Fµν is the electromagnetic field, and φ is the dilatonic field, with cou-
pling constant α. This theory admits static spherically symmetric solutions
representing black holes. Among these black hole solutions there are asymp-
totically flat ones [12, 13] as well as non-asymptotically flat configurations
[10, 11]. In the present work, we are interested in the non-asymptotically
flat black hole solutions
ds2 =
rγ(r − b)
dt2 −
rγ(r − b)
dr2 + r(r − b)dΩ2
, (2)
1 + γ
dr ∧ dt , e2αφ = ν2
. (3)
1− α2
1 + α2
. (4)
The constants b and r0 are related to the mass and to the electric charge of
the black hole through
M = (1− γ)b/4 , Q =
1 + γ
. (5)
The solutions (2),(3) interpolate between the Schwarzschild solution for
γ = −1 (α2 → ∞) and the Bertotti-Robinson solution for γ = +1 (α2 = 0).
For b > 0 the horizon at r = b hides the singularity at r = 0, while in the
extreme black hole case b = 0 the horizon coincides with the singularity.
This is a curious case, with vanishing mass but a finite electric charge. For
−1 < γ < 0 (α2 > 1) the central singularity is timelike and clearly naked
[11]. On the other hand, for 0 ≤ γ < 1 (0 < α2 ≤ 1), the central singularity
is null and marginally trapped [14], so that signals coming from the centre
never reach external observers. Thus in this case, extreme black holes can
be still considered as black holes indeed.
The statistical Hawking temperature of the black holes (2), computed as
usual by dividing the surface gravity by 2π is given by
. (6)
It is finite for all γ if b 6= 0. For b = 0 and −1 < γ < 0 (naked singularity).
the temperature is infinite, while for b = 0 and 0 < γ < 1 (extreme black
hole), the temperature vanishes.
The case b = γ = 0 is intriguing. Although this an extreme black
hole, the situation is different from that of asymptotically flat extreme black
holes. The near-horizon Euclidean extreme Reissner-Nordström geometry
is cylindrical, rather than conical, so that its statistical temperature is ar-
bitrary, contrary to the zero value derived from surface gravity [15]. In the
present case the two-dimensional Euclidean continuation of the metric (2)
with γ = 0 clearly has a conical singularity at r = b for all values of b,
including b = 0, leading for this particular extreme black hole to the finite
temperature TH = 1/4πr0, in agreement with the value (6). However this
result is questionable. A black hole with pointlike horizon and zero mass
clearly cannot radiate, so one should rather expect its temperature to be
zero. We will return to this question presently.
As black holes (2) radiate, they loose mass according to Stefan’s law
= −σAhT 4H , (7)
where σ is Stefan’s constant, and Ah = 4πr
1−γ is the horizon area.
Assuming that only electrically neutral quanta are radiated, (7) implies that
the horizon area decreases according to
(4π)3(1− γ)
−3(1+γ)
1+3γ , (8)
which is solved by
b(t) = r0
t− t0
)−1/3γ
(γ 6= 0) ,
b(t) = r0 exp
t− t0
(γ = 0) , (9)
where c = 3σ/16π3, and t0 is an integration constant. The outcome de-
pends on the sign of γ. For γ < 0, the Hawking temperature increases with
decreasing mass and the black hole collapses to a naked singularity (or evap-
orates away altogether in the Schwarzschild case γ = −1) in a finite time
according to b ∼ (t0 − t)1/3|γ|. On the other hand, for γ ≥ 0, the Hawking
temperature decreases (or is constant for γ = 0) with decreasing mass, and
the black hole evaporates in an infinite time, reaching the extreme black
hole state b = 0 only asymptotically.
We now proceed to a more precise evaluation of the temperature of non-
asymptotically flat black holes from the study of wave scattering in these
spacetimes. The wave equation
∇2φ = 0 (10)
does not generically allow for an exact solution in the spacetimes (2). How-
ever, it can be solved analytically [16] in the case of linear dilaton black
holes with γ = 0 and b 6= 0, with the metric
ds2 =
r − b
dt2 −
r − b
dr2 + r(r − b)dΩ2
, (11)
Considering the harmonic eigenmodes
φ(x) = ψ(r, t)Ylm(θ, ϕ) , ψ(r, t) = R(r)e
−iωt , (12)
we obtain the following radial equation:
r(r − b)∂rR
r − b
− l(l + 1)
R = 0 (13)
(ω̄2 ≡ ω2r20). Putting
, R = yiω̄f , (14)
reduces (13) to the equation
y(1−y)∂2yf+
1+2iω̄−2(1+ iω̄)y
ω̄2− iω̄− λ̄2−1/4
f = 0 , (15)
λ̄2 = ω̄2 − (l + 1/2)2 . (16)
This is a hypergeometric equation
y(1− y)∂2yf +
c− (a+ b+ 1)y
∂yf − abf = 0 , (17)
+ i(ω̄ + λ̄) , b =
+ i(ω̄ − λ̄) , c = 1 + 2iω̄ . (18)
It follows that the general solution of equation (13) is
R = C1
r − b
+ i(ω̄ + λ̄),
+ i(ω̄ − λ̄), 1 + 2iω̄; b− r
r − b
)−iω̄
− i(ω̄ + λ̄), 1
− i(ω̄ − λ̄), 1 − 2iω̄; b− r
.(19)
Putting
r − b
= ex/r0 , (20)
the partial wave near the horizon (x→ −∞) is thus
ψ ≃ C1eiω(x−t) +C2e−iω(x+t) . (21)
To obtain the behavior of the partial wave near spatial infinity, we must
expand the solutions of (15) in hypergeometric functions of argument 1/y.
The relevant transformation is
F (a, b, c; y) =
Γ(c)Γ(b− a)
Γ(b)Γ(c− a)
(−y)−aF (a, a+ 1− c, a+ 1− b; 1/y)
Γ(c)Γ(a − b)
Γ(a)Γ(c − b)
(−y)−bF (b, b+ 1− c, b+ 1− a; 1/y) . (22)
This leads to the asymptotic behavior
)−1/2(
i(λx−ωt) +B2e
−i(λx+ωt)
(λ = λ̄/r0), where the amplitudes of the asymptotic outgoing and ingoing
waves B1 and B2 are related to the amplitudes of the near-horizon outgoing
and ingoing waves C1 and C2 by
B1 = Γ(2iλ̄)
Γ(1 + 2iω̄)
Γ(1/2 + i(ω̄ + λ̄))2
Γ(1− 2iω̄)
Γ(1/2 − i(ω̄ − λ̄))2
B2 = Γ(−2iλ̄)
Γ(1 + 2iω̄)
Γ(1/2 + i(ω̄ − λ̄))2
Γ(1− 2iω̄)
Γ(1/2− i(ω̄ + λ̄))2
. (24)
Hawking radiation can be considered as the inverse process of scattering
by the black hole, with the asymptotic boundary condition B1 = 0 (the
outgoing mode is absent). The coefficient for reflection by the black hole is
then given by
|C1|2
|C2|2
|Γ(1/2 + i(ω̄ + λ̄))2|2
|Γ(1/2 + i(ω̄ − λ̄))2|2
cosh2 π(ω̄ − λ̄)
cosh2 π(ω̄ + λ̄)
. (25)
The resulting radiation spectrum is
= (eω/TH − 1)−1 . (26)
For high frequencies, λ̄ ≃ ω̄ = ω/r0, and we recover from (25) the Hawking
temperature as computed from the surface gravity,
. (27)
The above computation fails in the linear dilaton vacuum case b = 0.
The question of assigning a temperature to such massless black holes might
be evacuated by arguing that they cannot be formed, either through cen-
tral collapse of matter, or (as we have seen above) through evaporation of
massive black holes. Nevertheless, as a matter of principle one should con-
sider the possibility of primordial massless black holes. From the general
temperature law (6) these should have a finite temperature. On the other
hand, being massless they cannot radiate energy away, so their temperature
should vanish.
The question can be settled by solving the massless Klein-Gordon equa-
tion in the metric (11) with b = 0,
ds2 =
dt2 − r0
dr2 − r0rdΩ2 . (28)
This metric can be rewritten as
ds2 = Σ2
dτ2 − dx2 − dΩ2
, (29)
x = ln(r/r0) , τ = t/r0 , Σ = r0e
x/2 , (30)
showing that the linear dilaton vacuum metric is conformal to the product
M2 × S2 of a two-dimensional Minkowski spacetime with the two-sphere.
Performing also the redefinition
φ = Σ−1ψ , (31)
the Klein-Gordon equation (10) is reduced to
∇2φ = Σ−3
∂ττ − ∂xx −∇2Ω +
ψ = 0 , (32)
where ∇2Ω is the Laplacian operator on the two-sphere.
For a given spherical harmonic with orbital quantum number l, the re-
duced Klein-Gordon equation is thus
∇22ψl + (l + 1/2)2ψl = 0 , (33)
with ∇22 the Dalembertian operator on M2. Also, for a given spherical
harmonic the four-dimensional Klein-Gordon norm reduces to theM2 norm:
‖φ‖2 = 1
|g|g0µφ∗
∂µ φ =
dxψ∗l
∂τ ψl . (34)
Thus, the problem of wave propagation in the linear dilaton vacuum reduces
to the propagation of eigenmodes of a free Klein-Gordon field in two dimen-
sions, with effective mass µ = l+1/2. Clearly there is no reflection, so that
the linear dilaton vacuum does not radiate and hence its Hawking temper-
ature vanishes, contrary to the naive surface gravity value (6). A similar
reasoning holds in 2+1 dimensions for the BTZ vacuum [6] (M = L = 0),
which is conformal to M2 × S1.
We have shown that a complete analytical computation of the radia-
tion spectrum is possible for linear dilaton black hole solutions of EMD. For
massive black holes, this leads in the high frequency regime to a Planckian
distribution with a temperature independent of the black hole mass, in ac-
cordance with the surface gravity value. On the other hand, we find that
extreme, massless black holes do not radiate, thereby solving the paradox
presented by apparently hot (if the surface gravity temperature is taken
seriously) yet massless black holes.
Acknowledgements: J.C.F. thanks the LAPTH for the warm hospitality
during the elaboration of this work. He also thanks CNPq (Brazil) for
partial support. J.C.F. and G.T.M. thank the French-Brazilian scientific
cooperation CAPES/COFECUB for partial financial support.
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|
0704.0400 | The S-Matrix of AdS/CFT and Yangian Symmetry | arXiv:0704.0400v4 [nlin.SI] 27 Mar 2008
arxiv:0704.0400
AEI-2007-019
The S-Matrix of AdS/CFT
and Yangian Symmetry
Niklas Beisert
Max-Planck-Institut für Gravitationsphysik
Albert-Einstein-Institut
Am Mühlenberg 1, 14476 Potsdam, Germany
[email protected]
Abstract
We review the algebraic construction of the S-matrix of AdS/CFT. We
also present its symmetry algebra which turns out to be a Yangian of
the centrally extended su(2|2) superalgebra.
1 Introduction and Overview
Bethe’s ansatz [1] for solving a one-dimensional integrable model was and remains a
powerful tool in contemporary theoretical physics: 75 years ago it solved one of the
first models of quantum mechanics, the Heisenberg spin chain [2]; today it provides
exact solutions for the spectra of certain gauge and string theories and thus helps us
understand their duality [3] better. Since the discovery of integrable structures in planar
N = 4 supersymmetric gauge theory [4] and in planar IIB string theory on AdS5×S5 [5]
the tools for computing and comparing the spectra of both models have evolved rapidly.
We now have complete asymptotic Bethe equations [6, 7] which interpolate smoothly
between the perturbative regimes in gauge and string theory and which agree with all
available data.
In this note we will focus on the S-matrix [8] in the excitation picture above a ferro-
magnetic ground state. We start by reviewing the algebraic construction of the S-matrix
in Sec. 2. In Sec. 3 we subsequently show that this S-matrix has indeed a larger symmetry
algebra: a Yangian.
http://arxiv.org/abs/0704.0400v4
2 The Universal Enveloping Algebra U(su(2|2)⋉R2)
In this section the results on the S-matrix of AdS/CFT shall be reviewed from an al-
gebraic point of view. The applicable symmetry is a central extension h of the Lie
superalgebra su(2|2) which we consider first. We continue by presenting the Hopf al-
gebra structure of its universal enveloping algebra and its fundamental representation.
Finally, we comment on the S-matrix and its dressing phase factor.
Lie Superalgebra. The symmetry in the excitation picture for light cone string theory
on AdS5×S5 and for single-trace local operators in N = 4 supersymmetric gauge theory
is given by two copies of the Lie superalgebra [9, 10]
h := su(2|2)⋉ R2 = psu(2|2)⋉ R3. (2.1)
It is a central extension of the standard Lie superalgebras su(2|2) or psu(2|2), see [11].
It is generated by the su(2)× su(2) generators Rab, Lαβ, the supercharges Qαb, Saβ and
the central charges C, P, K. The Lie brackets of the su(2) generators take the standard
[Rab,R
d] = δ
d − δadRcb, [Lαβ ,Lγδ] = δ
Lαδ − δαδ Lγβ,
[Rab,Q
d] = −δadQγb + 12δ
d, [L
d] = +δ
d − 12δ
[Rab,S
δ] = +δ
δ − 12δ
δ, [L
δ] = −δαδ Scβ + 12δ
δ. (2.2)
The Lie brackets of two supercharges yield
{Qαb,Scδ} = δcbLαδ + δαδ Rcb + δcbδαδ C,
{Qαb,Qγd} = εαγεbdP,
{Saβ,Scδ} = εacεβδK. (2.3)
The remaining Lie brackets vanish.
Where appropriate, we shall use the collective symbol JA for the generators. The Lie
brackets then take the standard form
[JA, JB] = fABC J
C . (2.4)
For simplicity of notation, we shall pretend that all generators are bosonic; the general-
isation to fermionic generators by insertion of suitable signs and graded commutators is
straightforward.
Hopf Algebra. Next we consider the universal enveloping algebra U(h) of h. The
construction of the product is standard, and one identifies the Lie brackets (2.4) with
graded commutators. For the coproduct one can introduce a non-trivial braiding [12,13]
∆JA = JA ⊗ 1 + U [A] ⊗ JA (2.5)
∆Rab = R
b ⊗ 1 + 1⊗Rab,
∆Lαβ = L
β ⊗ 1 + 1⊗ Lαβ,
∆Qαb = Q
b ⊗ 1 + U+1 ⊗Qαb,
∆Saβ = S
β ⊗ 1 + U−1 ⊗Saβ,
∆C = C⊗ 1 + 1⊗ C,
∆P = P⊗ 1 + U+2 ⊗P,
∆K = K⊗ 1 + U−2 ⊗ K,
∆U = U ⊗ U .
Table 1: The coproduct of the braided universal enveloping algebra U(h).
with some abelian1 generator U (a priori unrelated to the algebra) and the grading
[R] = [L] = [C] = 0, [Q] = +1, [S] = −1, [P] = +2, [K] = −2. (2.6)
The coproduct is spelt out in Tab. 1 for the individual generators. The above grading
is derived from the Cartan charge of the sl(2) automorphism [11] of the algebra h and
therefore the coproduct is compatible with the algebra relations.
We should define the remaining structures of the Hopf algebra: the antipode S and
the counit ε [12,13]. The antipode is an anti-homomorphism which acts on the generators
S(1) = 1, S(U) = U−1, S(JA) = −U−[A]JA. (2.7)
The counit acts non-trivially only on 1 and U
ε(1) = ε(U) = 1, ε(JA) = 0. (2.8)
Cocommutativity. This coproduct is in general not quasi-cocommutative as can eas-
ily be seen by considering the central charges P, K in Tab. 1. To make it quasi-cocommu-
tative we have to satisfy the constraints [12]
1− U+2
1− U+2
⊗P, K⊗
1− U−2
1− U−2
⊗ K. (2.9)
They are solved by identifying the central charges P, K with the braiding factor U as
follows [13]
P = gα
1− U+2
, K = gα−1
1− U−2
. (2.10)
This leads to the following quadratic constraint
PK− gα−1P− gαK = 0. (2.11)
It was furthermore shown in [14] that the coproduct is quasi-triangular, at least at the
level of central charges, see also [15].
1Curiously, we can include the supersymmetric grading (−1)F in the generator U to manually impose
the correct statistics. This is helpful for an implementation within a computer algebra system. In this
case U would anticommute with fermionic generators.
Fundamental Representation. The algebra h has a four-dimensional representation
[10] which we will call fundamental. The corresponding multiplet has two bosonic states
|φa〉 and two fermionic states |ψα〉. The action of the two sets of su(2) generators has to
be canonical
Rab|φc〉 = δcb |φa〉 − 12δ
b |φc〉,
Lαβ|ψγ〉 = δγβ |ψα〉 − 12δ
β |ψγ〉. (2.12)
The supersymmetry generators must also act in a manifestly su(2)×su(2) covariant way
Qαa|φb〉 = a δba|ψα〉,
Qαa|ψβ〉 = b εαβεab|φb〉,
Saα|φb〉 = c εabεαβ|ψβ〉,
Saα|ψβ〉 = d δβα|φa〉. (2.13)
We can write the four parameters a, b, c, d using the parameters x±, γ and the constants
g, α as
g γ, b =
, c =
, d =
. (2.14)
The parameters x± (together with γ) label the representation and they must obey the
constraint
− x− − 1
. (2.15)
The three central charges C,P,K and U are represented by the values C, P,K and U
which read
1 + 1/x+x−
1− 1/x+x− , P = gα
, K =
, U =
. (2.16)
They furthermore obey the quadratic relation C2−PK = 1
. Note that the corresponding
quadratic combination of central charges C2−PK is singled out by being invariant under
the sl(2) external automorphism.
Fundamental S-Matrix. In [10,14] an S-matrix acting on the tensor product of two
fundamental representations was derived. It was constructed by imposing invariance
under the algebra h
[∆JA,S] = 0. (2.17)
We will not reproduce the result here, it is given in [14]. Note that we have to fix the
parameters ξ = U =
x+/x− in order to make the action of the generators compatible
with the coproduct (2.5).2
2This identification removes all braiding factors from the S-matrix in [14] which will thus satisfy the
standard Yang-Baxter (matrix) equation, see also [10, 16, 17].
This S-matrix has several interesting properties. Firstly, it is not of difference form;
it cannot be written as a function of the difference of some spectral parameters. Sec-
ondly, the S-matrix could be determined uniquely up to one overall function merely by
imposing a Lie-type symmetry (2.17) [10]. This unusual fact is related to an unusual
feature of representation theory of the algebra h: The tensor product of two fundamental
representations is irreducible in almost all cases [14].
Intriguingly this S-matrix is equivalent to Shastry’s R-matrix [18] of the one-dimen-
sional Hubbard model [19]. Furthermore the Bethe equations [10] contain two copies of
the Lieb-Wu equations for the Hubbard model [20]. These observations of [14] estab-
lish a link between an important model of condensed matter physics and string theory
(complementary to the one in [21]).
Finally, let us note that one can derive (asymptotic) Bethe equations from the S-
matrix and thus confirm the conjecture in [6]. So far this step has been performed in
two different ways: by means of the nested coordinate [10] and the algebraic [17] Bethe
ansatz.
Phase Factor. The remaining overall phase factor of the S-matrix clearly cannot be
determined by demanding invariance under h. The phase factor was computed to some
approximation from gauge theory [22] and from string theory [23]. The problem of an
algebraically undetermined phase factor is in fact generic. Usually one imposes a further
crossing symmetry relation to obtain a constraint on it. Indeed the known string phase
factor is consistent with crossing symmetry [24] as was shown in [25]. By substituting a
suitable ansatz [26] for the phase factor into the crossing symmetry relation a conjecture
for the all-orders phase factor at strong coupling was made in [27].
A corresponding all-orders expansion at weak coupling was presented in [7]. The
latter conjecture was obtained by a sort of analytic continuation in the perturbative
order of the series. Let us illustrate this principle by means of a very simple example:
Consider the rational function f(x) = 1/(1−x). It has the following expansions at x = 0
and at x = ∞
n, f(x)
−n (2.18)
with an = 1 and bn = −1. When we consider an and bn as analytic functions of the
index, we can make the observation (“reciprocity”)
an = −b−n. (2.19)
Of course there are various ways in which the two functions +1 and −1 could be related,
but the choice (2.19) appears to work for a surprisingly large class of functions!3 It was
proved in [30] that it does apply for the conjectured expansion of the phase factor. Very
useful integral expressions for the phase have recently appeared in [31]. The analytic
expression of the dressing phase can formally be obtained from the psu(2, 2|4) Bethe
3Among other physical examples, we have identified circular Maldacena-Wilson loops [28] and non-
critical string theory [29] where this reciprocity can be applied. Furthermore, summation by the Euler-
MacLaurin formula (also known as zeta-function regularisation) is consistent with it. I thank Curt
Callan, Marcos Mariño and Tristan McLoughlin for discussions of this principle.
equations [32] (see however appendix D in [33]) in analogy to the covariant approach
of [34, 21, 35]. While this proposal may seem to be encouraging in general, it is at the
same time strange from the Hopf algebra point of view to use an S-matrix which does
not obey the crossing relation [32]. This calls for further investigations.
Several tests of the phase have recently appeared, they are based on four-loop unitary
scattering methods [36], numerical evaluation [37, 38], analytic methods [37, 30, 39] and
on taking a certain highly non-trivial limit [40].
3 The Yangian Y(su(2|2) ⋉ R2)
In the section we investigate Yangian symmetry [41,42] for the above S-matrix. We will
start with a very brief review of Yangian symmetry for generic S-matrices (see [43] for
more extensive reviews), and then we apply the framework to the S-matrix discussed
above.
Yangians and S-Matrices. Typically the symmetries of rational S-matrices are of
Yangian type. The Yangian Y(g) of a Lie algebra g is a deformation of the universal
enveloping algebra of half the affine extension of g.
More plainly, it is generated by the g-generators JA and the Yangian generators ĴA.
Their commutators take the generic form
[JA, JB] = fABC J
[JA, ĴB] = fABC Ĵ
C , (3.1)
and they should obey the Jacobi and Serre relations
J[A, [JB, JC]]
J[A, [JB, ĴC]]
Ĵ[A, [ĴB, JC]]
2fAGD f
F fGHKJ
{DJEJF}. (3.2)
The symbol fABC = gADgBEf
C represents the structure constants f
C with two indices
lowered by means of the inverse of the Cartan-Killing forms gAD and gBE . The brackets
{ } and [ ] at the level of indices imply total symmetrisation and anti-symmetrisation,
respectively. Finally, ~ is a scale parameter whose value plays no physical role. The first
two relations lead to a constraint on the structure constants fABC . The third relation
a deformation of the Serre relation for affine extensions of Lie algebras.
The Yangian is a Hopf algebra and the coproduct takes the standard form
∆JA = JA ⊗ 1 + 1⊗ JA,
∆ĴA = ĴA ⊗ 1 + 1⊗ ĴA + 1
~fABCJ
B ⊗ JC . (3.3)
where fABC = gBDf
C . The antipode S is defined by
S(JA) = −JA, S(ĴA) = −ĴA + 1
~fABCf
D, (3.4)
4For g = su(2) it has to be replaced by a quartic relation.
and the counit ε takes the standard form
ε(1) = 1, ε(JA) = ε(ĴA) = 0. (3.5)
For the study of integrable systems, the evaluation representations of the Yangian
are of special interest. For these the action of the Yangian generators ĴA is proportional
to the Lie generators
ĴA|u〉 = ~uJA|u〉. (3.6)
Here |u〉 is some state of the evaluation module with spectral parameter u. This Yangian
representation is finite-dimensional if the g-representation is. One merely has to ensure
that the Serre relation (3.2) is satisfied. This is indeed not the case for all representations
of all Lie algebras. The power of the Yangian symmetry lies in the fact that tensor
products of evaluation representations are typically irreducible (except for special values
of their spectral parameters). This allows for simple proofs (e.g. for the Yang-Baxter
relation) by representation theory arguments.
Let us finally consider the connection to the S-matrix. The S-matrix is a permutation
operator; it acts by interchanging two modules of the algebra
S : V1 ⊗ V2 → V2 ⊗ V1. (3.7)
In particular, for the tensor product of two evaluation modules one has
S|u1, u2〉 ∼ |u2, u1〉. (3.8)
Invariance of the S-matrix under the Yangian means
[∆JA,S] = [∆ĴA,S] = 0 (3.9)
for all generators JA, ĴA. The existence of such an S-matrix is equivalent to quasi-
cocommutativity of Y(g). Note that only the difference of spectral parameters appears in
the invariance condition: We can write the action of the coproduct of Yangian generators
on the evaluation module |u1, u2〉 as
∆ĴA ≃ (u1 − u2)JA ⊗ 1 + u2∆JA + ~fABCJB ⊗ JC . (3.10)
Here the first equation in (3.9) ensures that the term proportional to u2 drops out from
the second equation. Therefore the S-matrix typically depends on the difference u1 − u2
of spectral parameters only.
Yangians in AdS/CFT. Yangian symmetries for planar AdS/CFT have been inves-
tigated in [44], both for classical string theory and for gauge theory at leading order,
see also [45] Yangian symmetry also persists to higher perturbative orders in both mod-
els [22, 46] and it is likely that it also exists at finite coupling. This Yangian can be
understood as a symmetry of the Hamiltonian on an infinite world sheet or as an expan-
sion of the full monodromy matrix. The Lie symmetry in this picture is psu(2, 2|4) and
the Yangian would be Y(psu(2, 2|4)).
Here we consider a different picture of well-separated excitations on a ferromagnetic
ground state and of their scattering matrix. In this picture the Lie symmetry reduces to
two copies of h and the corresponding Yangian would be Y(h). Our Yangian should arise
as a subalgebra of the full Yangian Y(psu(2, 2|4)) when acting on asymptotic excitation
states.
Hopf Algebra. Let us now consider Y(h). We have already studied the universal
enveloping algebra U(h). All we still need to do is to introduce one generator ĴA for each
JA obeying the relations (3.1,3.2), and define its coproduct, antipode as well as counit.
In (2.5) we have defined a braided coproduct for the universal enveloping algebra.
For consistency with the Serre relations, we also have to apply an analogous braiding to
the standard Yangian coproduct
∆ĴA = ĴA ⊗ 1 + U [A] ⊗ ĴA + ~fABCJBU [C] ⊗ JC . (3.11)
Note that lowering an index requires to use the inverse Cartan-Killing form of the algebra.
In the case of h the Cartan-Killing form is degenerate and we need to extend the algebra
by the sl(2) outer automorphism, see [14]. Effectively, lowering an index leads to an
interchange of the automorphism generators with the central charges. We refrain from
spelling out the Cartan-Killing form or the modified structure constants. Instead we
present the complete set of coproducts of Yangian generators in Tab. 2, where we also
fix the value of ~.
For the sake of completeness we state the antipode5 and the counit
S(ĴA) = −U−[A]ĴA, ε(ĴA) = 0. (3.12)
Cocommutativity. An important question is if this coproduct can be quasi-cocom-
mutative.6 A first step is to consider the central generators Ĉ, P̂, K̂. For that purpose
it is favourable to choose suitable combinations
Ĉ′ = Ĉ+ 1
gα−1P− 1
P̂′ = P̂+ C
P− 2gα
K̂′ = K̂− C
K− 2gα−1
, (3.13)
for whom the coproduct almost trivialises
∆Ĉ′ = Ĉ′ ⊗ 1 + 1⊗ Ĉ′,
∆P̂′ = P̂′ ⊗ 1 + U+2 ⊗ P̂′,
∆K̂′ = K̂′ ⊗ 1 + U−2 ⊗ K̂′. (3.14)
The combination Ĉ′ is already cocommutative, and in order to make the generators P̂′,
K̂′ cocommutative we have to set as above in (2.9,2.10)
P̂′ = iguPP, K̂
′ = iguKK (3.15)
with two universal constants uP and uK. With this choice, Ĉ, P̂, K̂ also become cocom-
mutative because they differ from Ĉ′, P̂′, K̂′ only by central elements.
5Note that fA
= 0 here, so there is no contribution from the Lie generators.
6The braiding factors in (3.11) turn out to be very important for the Yangian. It can easily be
seen that without them the coproduct cannot be quasi-cocommutative. This is in contradistinction
to the universal enveloping algebra where the braided as well as the unbraided coproduct are quasi-
cocommutative.
∆R̂ab = R̂
b ⊗ 1 + 1⊗ R̂ab
Rac ⊗Rcb − 12R
b ⊗Rac
SaγU+1 ⊗Qγb − 12Q
bU−1 ⊗Saγ
δab S
γU+1 ⊗Qγd + 14δ
dU−1 ⊗Sdγ ,
∆L̂αβ = L̂
β ⊗ 1 + 1⊗ L̂αβ
Lαγ ⊗ Lγβ + 12L
β ⊗ Lαγ
QαcU−1 ⊗Scβ + 12S
βU+1 ⊗Qαc
δαβ Q
cU−1 ⊗Scδ − 14δ
δU+1 ⊗Qδc,
∆Q̂αb = Q̂
b ⊗ 1 + U+1 ⊗ Q̂αb
LαγU+1 ⊗Qγb + 12Q
b ⊗ Lαγ
RcbU+1 ⊗Qαc + 12Q
c ⊗Rcb
CU+1 ⊗Qαb + 12Q
b ⊗ C
εαγεbdPU−1 ⊗Sdγ − 12ε
αγεbdS
γU+2 ⊗P,
∆Ŝaβ = Ŝ
β ⊗ 1 + U−1 ⊗ Ŝaβ
RacU−1 ⊗Scβ − 12S
β ⊗Rac
LγβU−1 ⊗Saγ − 12S
γ ⊗ Lγβ
CU−1 ⊗Saβ − 12S
β ⊗ C
εacεβδKU+1 ⊗Qδc + 12ε
acεβδQ
cU−2 ⊗ K,
∆Ĉ = Ĉ⊗ 1 + 1⊗ Ĉ
PU−2 ⊗ K− 1
KU+2 ⊗P,
∆P̂ = P̂⊗ 1 + U+2 ⊗ P̂
− CU+2 ⊗P+P⊗ C,
∆K̂ = K̂⊗ 1 + U−2 ⊗ K̂
+ CU−2 ⊗ K− K⊗ C.
Table 2: The coproduct of the Yangian generators in Y(h).
Fundamental Evaluation Representation. For the fundamental evaluation repre-
sentation we make the ansatz7
ĴA|X 〉 = ig(u+ u0)JA|X 〉. (3.16)
By comparison with (3.13,3.15) we can infer that u has to be related to the parameters
of the fundamental representation by
u = x+ +
= x− +
(x+ + x−)(1 + 1/x+x−) . (3.17)
Furthermore uP and uK in (3.15) have to both coincide with the universal constant
u0 = uP = uK.
As an aside we state the eigenvalue of the quadratic combination
CĈ − 1
PK̂ − 1
KP̂ = 1
ig(u+ u0). (3.18)
Fundamental S-Matrix. Using the coproducts in Tab. 2 we have confirmed that the
S-matrix is also invariant under all of the Yangian generators
[∆ĴA,S] = 0. (3.19)
We have used a computer algebra system to evaluate the action of the Yangian gener-
ators and the S-matrix.9 To show invariance requires heavy use of the identity (2.15).
Superficially it is very surprising to find all these additional symmetries of the S-matrix.
The deeper reason however should be that the coproduct is quasi-cocommutative. We
have thus proved quasi-cocommutativity when acting on fundamental representations.
It is interesting to see that the S-matrix is based on standard evaluation represen-
tations of the Yangian. Nevertheless, it is not a function of the difference of spectral
parameters. This unusual property traces back to the link between the spectral param-
eter u and the h-representation parameters x± in (3.17). The latter is again related to
the braiding in the coproduct (3.11).
As our S-matrix is equivalent [14] to Shastry’s R-matrix, our Yangian is presumably
an extension of the su(2)×su(2) Yangian symmetry of the Hubbard model found in [47].
4 Conclusions and Outlook
In this note we have reviewed the construction of the S-matrix with centrally extended
su(2|2) symmetry that appears in the context of the planar AdS/CFT correspondence
and the one-dimensional Hubbard model. We have furthermore shown that the S-matrix
has an additional Yangian symmetry whose Hopf algebra structure we have presented.
This Yangian is not quite a standard Yangian, but its coproduct needs to be braided in
order to be quasi-cocommutative. This fact is intimately related to the existence of a
7We believe, but we have not verified that this is compatible with the Serre relations (3.2).
8It is conceivable that a further consistency requirement fixes the value of u0, presumably to zero.
9We have also confirmed the invariance of the singlet state found in [10].
triplet of central charges with non-trivial coproduct and leads to the wealth of unusual
features of the S-matrix.
In connection to the Yangian there are many points left to be clarified. Most im-
portantly the representation theory needs to be understood. Which representations of h
lift to evaluation representations of Y(h)? At what values of the spectral parameters do
their tensor products become reducible? This information could be used to prove that
the coproduct is quasi-cocommutative. Also the Yang-Baxter equation for the S-matrix
should follow straightforwardly. It might also give some further understanding of bound
states [48].
Then it would be highly desirable to construct a universal R-matrix for this Yangian
and show that it is quasi-triangular. This would put large parts of the integrable structure
for arbitrary representations of this algebra on solid ground much like for the case of
generic simple Lie algebras.
Some further interesting questions include: Is this Yangian the unique quasi-co-
commutative Hopf algebra based on h? Does the double Yangian [42] exist and what is
its structure? Can the sl(2) automorphism of the algebra be included at the Yangian
level such that the coproduct is quasi-cocommutative? What would the representations
be in this case?
Acknowledgements. I am grateful to C. Callan, D. Erkal, A. Kleinschmidt, P. Ko-
roteev, N. MacKay, M. Mariño, T. McLoughlin, J. Plefka, F. Spill and B. Zwiebel for
interesting discussions.
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|
0704.0401 | Modelling the Near-IR Spectra of Red Supergiant-dominated Populations | Stellar Populations as Building Blocks of Galaxies
Proceedings IAU Symposium No. 241, 2007
A. Vazdekis et alr., eds.
c© 2007 International Astronomical Union
DOI: 00.0000/X000000000000000X
Modelling the Near-IR Spectra of Red
Supergiant-dominated Populations
Ariane Lançon1, Jay S. Gallagher2, Richard de Grijs3,
Peter Hauschildt4, Djazia Ladjal5, Mustapha Mouhcine6,
Linda J. Smith7, Peter R. Wood8, Natascha Förster Schreiber9
1Observatoire de Strasbourg (UMR 7550), 11 rue de l’Université, 67000 Strasbourg, France
email: [email protected]
2Dept. of Astronomy, 5534 Sterling, University of Wisconsin, Madison, WI 53706, USA
3Dept. of Physics & Astronomy, University of Sheffield, Hicks Building, Honusfield Rd.,
Sheffield S3 7RH, UK
4Hamburger Sternwarte, Gojenbergsweg 112, 21029 Hamburg, Germany
5Institute of Astronomy, Katholieke Universiteit, Celestijnenlaan 200B, 3001 Leuven, Belgium
6Astrophysics Research Institute, Liverpool John Moores University, Twelve Quays House,
Egerton Wharf, Birkenhead, CH41 1LD, UK
7Space Telescope Science Institute, 3700 San Martin Drive, Baltimore, MD 21218, USA
8RSAA, Mt Stromlo Observatory, Cotter Road, Weston Creek, ACT 2611, Australia
9MPI für Extraterrestrische Physik, Giessenbachstrasse, 85741 Garching, Germany
Abstract. We report on recent progress in the modelling of the near-IR spectra of young stellar
populations, i.e. populations in which red supergiants (RSGs) are dominant. First, we discuss
the determination of fundamental parameters of RSGs from Phoenix model fits to their near-IR
spectra; RSG-specific surface abundances are accounted for and effects of the microturbulence
parameter are explored. New population synthesis predictions are then described and, as an
example, it is shown that the spectra of young star clusters in M82 can be reproduced very well
from 0.5 to 2.4µm. We warn of remaining uncertainties in cluster ages.
Keywords. galaxies: stellar content, galaxies: starburst, galaxies: star clusters, galaxies: indi-
vidual (M82), infrared: galaxies, infrared: stars, stars: supergiants
1. Introduction
Red supergiant stars (RSGs) provide most of the near-IR light emitted by young stellar
populations, such as those in starburst galaxies. As star forming environments tend to
be dusty, rest-frame optical analyses are incomplete (highly obscured populations are
missed) and it is crucial to improve our understanding of spectra at longer wavelengths. In
the past, the near-IR analysis of young stellar populations has often been restricted to the
determination of the average properties of the dominant stars, such as their spectral types
or abundances. The subsequent interpretation of these results in terms of precise stellar
population ages and star formation histories remains an enormous challenge, as it requires
(i) a good understanding of the near-IR spectra of individual RSGs and (ii) adequate
stellar evolution tracks. We have started a programme that aims at providing state of
the art predictions for the emission of RSG-dominated populations and at characterizing
remaining uncertainties. Currently, the project focuses on wavelengths between 0.81 and
2.4µm and spectral resolutions of order λ/δλ = 103.
http://arxiv.org/abs/0704.0401v1
2 A. Lançon et al.
2. Empirical and synthetic spectra of red supergiants
In principle, synthetic stellar spectra are more practical for the prediction of galaxy
spectra than empirical ones, because theory allows us to sample parameter space without
biases. Lançon et al. (2007) show that modern theoretical spectra can reproduce the near-
IR (+optical) emission of giant stars well down to effective temperatures Teff ≃ 3400K,
but that they are not yet satisfactory at lower temperatures and higher luminosities. They
used new Phoenix models to compute spectra at the necessary resolution (0.1 Å before
smoothing), with solar abundances and with the RSG-specific abundances obtained as
the result of internal mixing along stellar evolution tracks; the models include some 109
individual molecular and atomic lines, assume spherical symmetry, and allow dust to form
if conditions are fulfilled. Model limitations include the assumptions of local thermal
equilibrium (LTE) and of hydrostatic equilibrium. A sample of 101 empirical spectra
covering wavelengths between 0.51, 0.81 or 0.90µm and 2.4µm was used for comparison
(Lançon & Wood 2000, Lançon et al. in preparation). The data were acquired with
CASPIR on the 2.3m ANU Telescope at Siding Spring and with SpeX at IRTF, Hawaii.
Below Teff ∼ 3400K, uncertain input line lists are a problem in the models (especially
for molecular bands around 1µm). At high luminosity (luminosity class Ia and Iab), the
main difficulty is to reproduce simultaneously extremely deep CN bands and the relative
strengths of the CO bandheads around 1.7µm and at 2.3µm. RSG-specific abundances
improve fits to the CN bands. Exploratory calculations show that values near 10 km/s for
the “microturbulence” (a 1D-model parameter that hides poorly understood 3D physical
phenomena) may be able to solve both problems (Fig. 1, top left). The calculation of a
new grid has been launched to explore this further. In the mean time, the study shows
that the population synthesis community still has to rely on empirical spectra for RSGs,
and it warns that the lack of satisfactory stellar models implies large uncertainties on
the derived fundamental parameters of the observed stars.
3. Population synthesis using averaged stellar spectra
In order to compute spectra of synthetic populations, we have constructed three se-
quences of average empirical spectra, corresponding to luminosity classes Ia, Iab and
Ib/II. Each subset was sorted into bins according to the estimated Teff , the spectra were
dereddened (an estimate of the reddening is provided by the model fits), and averages
were computed. The sequences shown in Fig. 1 (top right) account for varying micro-
turbulence in a preliminary way, based on the limited number of high microturbulence
models available to us at the time of this writing. We chose to flag any star with an initial
mass above 7M⊙ as a supergiant, which implies that the new spectra affect predictions
up to the age of about 75Myr (Fig. 1, middle left). We note that predictions vary signif-
icantly depending on the adopted evolutionary tracks; different authors predict different
red supergiant lifetimes, and main sequence rotation affects both the surface abundances
and the final red (and blue) supergiant numbers.
4. Star clusters in M82
The synthetic spectra of single stellar populations (SSPs) at solar metallicity are com-
pared with those of young star clusters in starbursts, such as M82-L and M82-F (Smith
& Gallagher 1999). The selected clusters are massive (well above 105M⊙), i.e. stochastic
effects due to an underpopulated RSG-branch are avoided. A few have well determined
optical ages (based on standard non-rotating evolutionary tracks). Figure 1 (middle right
Modelling Red Supergiant Populations 3
3000 4000 5000
Age = 18 Myr
Av = 1
0.5 1.0 1.5 2.0 2.5
Age=18
Av=3.7
Rv=2.4
chi2=1.58934
1.0 1.5 2.0
Wavelength (micron)
12 CO
[FeII]
1.55 1.60 1.65 1.70
Wavelength (micron)
bad pixels
2.1 2.2 2.3
Figure 1. Top left: Spectrum of an M0Ia RSG compared with models with vmicroturb=2km/s
(top: 4200K, log(g)=-1, AV =4.4) and with vmicroturb=10 km/s (bottom: 4500K, log(g)=0,
AV =4.7; note the improved CN at 1.1µm and CO around 1.6 and 2.3µm). Top right: Param-
eters assigned to the new sequences of average spectra, superimposed on the solar metallicity
tracks of Bressan et al. (1993). Middle left: Comparison of a new SSP spectrum (black) with
the standard predictions of Pegase.2 (differences are largest between 10 and 20Myr). Middle
right: Best near-IR fit to the spectrum of cluster M82-L. The extinction law of Cardelli et al.
(1989) with RV =2.4 allows us to also reproduce the optical spectrum (from Smith & Gallagher
1999). The error spectrum and the χ2 weighting function are shown. Bottom: Zoom-ins of the
H and K windows.
4 A. Lançon et al.
and bottom) shows cluster L, the cluster observed with SpeX with the best signal-to-
noise ratio: an excellent fit is obtained over the whole available range in wavelength.
Such results make the new models valuable tools for purposes such as weak emission line
measurements. The χ2-test restricted to near-IR wavelengths not affected by strong tel-
luric absorption shows that age is formally determined to an accuracy of about ±10Myr.
Because of strong reddening, the optical age of cluster L cannot be determined as well
as that of cluster F: 50-70Myr (Gallagher & Smith 2001, McCrady et al. 2005, Bastian
et al. 2007). For F, our current models provide a near-IR age range of 32 to 46Myr.
This small but nevertheless significant disagreement calls for several comments. (i) Be-
fore accounting for luminosity-dependent microturbulence, we found a near-IR age of 10
to 20Myr; we hope that our next generation of synthetic stellar spectra will significantly
reduce uncertainties originating in uncertain fundamental parameters of stars. (ii) The
spectrum used for optical age-dating and our near-IR spectrum have different slopes in
the region of overlap. This suggests slightly different positions were observed: the ob-
scuration across M82-F is not at all uniform. In addition, a younger cluster located at
very small projected distance might contaminate the near-IR data. (iii) Modified stellar
tracks (e.g. including rotation) might affect optical ages as well as near-IR ones.
5. Conclusions
The spectra of young stellar populations at solar metallicity, observed at R∼ 103, can
now be modelled well from the optical through the near-IR. Nevertheless, ages based on
near-IR spectra remain severely affected by uncertainties. They are due mainly to system-
atic errors, which further work needs to characterize and reduce. Errors are associated on
one hand with the fundamental parameters of red supergiant stars (theoretical spectra,
microturbulence, surface abundances of C, N and O, non-LTE, variability, winds, giant-
supergiant transition), and on the other with evolutionary tracks (convection, opacities,
rotation, binarity, effects of a dense environment). We expect rapid progress in stellar at-
mosphere models to provide us with tools to test stellar tracks further. Complete optical
and near-IR spectra of massive clusters such as those of M82 are useful test cases for the
identification and correction of systematic errors, but even they are not trivial to exploit
(due to inhomogeneous background populations and extinction, mass segregation, etc.).
References
Bastian, N., Konstantopoulos, I., Smith, L.J. & Gallagher, J.S. 2007, MNRAS in press
Cardelli, J.A., Clayton, G.C. & Mathis, J.S. 1989 ApJ 345, 245
Gallagher, J.S. & Smith, L.J. 1999 MNRAS 304, 540
Lançon, A. & Wood, P.R. 2000, A&AS 146, 217
Lançon, A., Hauschildt, P., Ladjal, D. & Mouhcine, M. 2007, A&A in press
McCrady, N., Graham, J.R. & Vacca, W.D. 2005 ApJ 621, 278
Smith, L. J. & Gallagher, J. S. 2001 MNRAS 326, 1027
Discussion
Gustafsson: Do the models with high microturbulence include turbulent pressure in a
consistent way?
Lançon (after discussion with P.H. and H. Lamers): No. But the microturbu-
lent velocities required to reproduce the spectra with 1D models are supersonic, which
suggests that the actual process is not microturbulence... Therefore it is unclear how to
relate this parameter of 1D models to pressure.
Introduction
Empirical and synthetic spectra of red supergiants
Population synthesis using averaged stellar spectra
Star clusters in M82
Conclusions
|
0704.0402 | Locating the peaks of least-energy solutions to a quasilinear elliptic
Neumann problem | LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS TO
A QUASILINEAR ELLIPTIC NEUMANN PROBLEM
YI LI AND CHUNSHAN ZHAO
Abstract. In this paper we study the shape of least-energy solutions to the
quasilinear problem εm∆mu−u
m−1 + f (u) = 0 with homogeneous Neumann
boundary condition. We use an intrinsic variation method to show that as
ε → 0+, the global maximum point Pε of least-energy solutions goes to a
point on the boundary ∂Ω at the rate of o(ε) and this point on the boundary
approaches to a point where the mean curvature of ∂Ω achieves its maximum.
We also give a complete proof of exponential decay of least-energy solutions.
1. Introduction and statement of results
In this paper we study the shape of certain solutions to the following quasilinear
elliptic Neumann problem:
(1.1)
εm∆mu− um−1 + f (u) = 0, u > 0 in Ω,
= 0 on ∂Ω,
where m (2 ≤ m < N) and 0 < ε ≤ 1 are constants and Ω ⊆ RN (N ≥ 3)
is a smooth bounded domain. The operator ∆mu = div(|∇u|m−2 ∇u) is the m-
Laplacian operator, and ν is the unit outer normal to ∂Ω.
Problem (1.1) appears in the study of non-Newtonian fluids, chemotaxis and
biological pattern formation. For example, in the study of non-Newtonian fluids,
the quantity m is a characteristic of the medium: media with m > 2 are called
dilatant fluids, and those with m < 2 are called pseudo-plastics. If m = 2, they
are Newtonian fluids (see [3] and its bibliography for more backgrounds). For
the case m = 2, (1.1) is also known as the stationary equation of the Keller–
Segal system in chemotaxis [14] or the limiting stationary equation of the so-called
Gierer–Meinhardt system in biological pattern formation (see [23]).
First let us recollect some results related to our problem. In a series of remarkable
papers, C.-S. Lin, W.-M. Ni and I. Takagi [14], Ni and Takagi [17], [18] studied the
Neumann problem for certain elliptic equations, including
(1.2)
d∆u− u+ up = 0, u > 0 in Ω,
= 0 on ∂Ω,
where d > 0, p > 1 are constants, and p is subcritical, i.e., p < N+2
. First, Lin,
Ni and Takagi [14] applied the mountain-pass lemma [1] to show the existence of
Key words and phrases. Quasilinear Neumann problem, m-Laplacian operator, least-energy
solution, exponential decay, mean curvature.
http://arxiv.org/abs/0704.0402v1
2 YI LI AND CHUNSHAN ZHAO
a least-energy solution ud to (1.2), by which is meant that ud has the least energy
among all solutions to (1.2) with the energy functional
Id (u) =
|∇u|2 + 1
u2 − 1
defined on W 1,2 (Ω). Hereinafter u+ = max {u, 0} and u− = min {u, 0}. Then in
[17], [18], Ni and Takagi investigated the shape of the least-energy solution ud as
d becomes sufficiently small, and showed that ud has exactly one peak (i.e., local
maximum of ud) at Pd ∈ ∂Ω. Moreover, as d tends to zero, Pd approaches a point
where the mean curvature of ∂Ω achieves its maximum. See [15] for a review in
this field. Also see [16] for the critical case p = N+2
, and [5], [6], [7], [8], [9] for
existence and properties of multiple-peaks solutions to (1.2).
From now on we make some hypotheses on f : R → R, as follows.
(H2) f (t) ≡ 0 for t ≤ 0 and f ∈ C1 (R).
(H3) f(t) = O (t
p) as t→ ∞ with m− 1 < p <
N (m− 1) +m
N −m .
(H4) Let F (t) =
f (s) ds. Then there exists a constant θ ∈
that F (t) ≤ θtf (t) for t > 0.
f (t)
is strictly increasing for t > 0 and f (t) = O
tm−1+δ
as t→ 0+ with
a constant δ > 0.
(H6) Let g (u) =
(m− 1)um−1 − uf ′ (u)
um−1 − f (u)
. Then g (u) is non-increasing on (uc,∞),
where uc is the unique positive solution for f (t) = t
Next we present some preliminary knowledge about least energy solutions of the
following problem:
(1.3)
∆mu− um−1 + f(u) = 0 in RN
u > 0 in RN
As before we define an “energy functional” I:W 1,m(RN ) −→ R associated with
(1.3) by
(1.4) I(ṽ) =
(εm |∇ṽ|m + |ṽ|m)− F (ṽ+)
Next let us give a remark on ground states to the problem 1.3. Here by a
ground state we mean a non-negative nontrivial C1 distribution solution which
tends to zero at ∞. For case m = 2, it is well known that the problem 1.3 has
a unique ground state (up to translations) which is radially symmetric [4]. For
case 2 < m < N uniqueness and radial symmetry of ground states are still open.
But the Steiner symmetrization tells us the least-energy solutions must be radially
symmetric (certainly least-energy solutions are ground states). Our assumptions
guarantee that the uniqueness (up to translations) of radial ground states (see
[20]), which implies the uniqueness of least-energy solutions of the problem (1.3).
Exact exponential decay of radial ground states was given in [11], thus we have the
following proposition about the unique radial least-energy solution to problem 1.3:
Proposition 1.1. Under assumptions (H2)–(H6), there is a unique least energy
solution w(x) for (1.3) satisfying:
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 3
( i ) w is radial, i.e., w(x) = w(|x|) = w(r) and w ∈ C1(RN ) with
w(0) = maxX∈RN w(x), w
′(0) = 0 and w′(r) < 0, ∀r > 0.
(ii) limr−→∞ w(r)r
m(m−1) e(
m r = C0 > 0 for some constant C0 and
limr−→∞
w′(r)
Remark 1.1. A good example for f (t) which satisfies all hypotheses (H2)–(H6) is
f (t) = tp for m− 1 < p < N (m− 1) +m
Next we define an “energy functional” Jε : W
1,m (Ω) → R associated with (1.1)
(1.5) Jε (v) =
(εm |∇v|m + |v|m)− F (v+)
with F (v+) =
f (s) ds. Then the well-known mountain-pass lemma [1] implies
(1.6) cε = inf
t∈[0,1]
Jε (h (t))
is a positive critical value of Jε, where Γ is the set of all continuous paths joining
the origin and a fixed nonzero element e ∈ W 1,m (Ω) such that e ≥ 0 and Jε (e) ≤ 0.
It turns out cε can also be characterized as follows:
cε = inf
Jε (u)
u ∈ W 1,m (Ω) ; u ≥ 0, u 6≡ 0,
(εm |∇u|m + um) dx =
f (u)u dx
(1.7) cε = inf
M [u] | u ∈ W 1,m (Ω) , u 6≡ 0 and u ≥ 0 in Ω
M [u] = sup
Jε (tu) .
Hence cε is the least positive critical value and a critical point uε of Jε with critical
value cε is called a least-energy solution. Notice also that if we let
c∗ = I(w) =
(|∇w|m + wm) dx−
F (w) dx,
where w is the unique least energy solution of (1.3), then c∗ can also be characterized
(1.8) c∗ = inf
M∗ [v] | v ∈ W 1,m
, v 6≡ 0 and v ≥ 0 in RN
M∗ [v] = sup
I (tv) .
We refer to Lemma 2.1 of [13] for the above characterizations.
Next we consider the following problem:
4 YI LI AND CHUNSHAN ZHAO
v ∈W 1,m
with RN+ =
(x1, · · · , xN ) ∈ RN , xN ≥ 0
and satisfies
(1.9)
∆mv − vm−1 + f (v) = 0, v > 0 in RN+ ,
= 0 on xN = 0.
The solutions of (1.9) can be characterized as critical points of the functional defined
over W 1,m
as follows.
(ṽ) =
|∇ṽ|m + ṽm) dx−
F (ṽ+) dx.
Similarly as above the least positive critical value C∗ corresponding to least energy
solutions of (1.9) can be characterized as
(1.10) C∗ = inf
ṽ∈W 1,m(RN+ ),ṽ≥0,ṽ 6≡0
(tṽ)
and moreover
(1.11) C∗ =
due to the boundary condition in (1.9) and the fact that w is radial and hence
= 0. We also refer to Lemma 2.1 of [13] for the above characterization of C∗.
In Theorem 1.3 of [13], we proved the following theorem.
Theorem 1.1. Under hypotheses (H2)–(H6), let uε be a least-energy solution of
(1.1). Then all local maximum points(if more than one) of uε aggregate to a global
maximum point Pε at a rate of o(ε) and dist(Pε, ∂Ω)/ε→ 0 as ε → 0+, where
dist(·, ·) is the general distance function. Moreover, we have the following upper-
bound estimate for cε as ε→ 0+:
(1.12) cε ≤ εN
c∗ − (N − 1) max
H (P ) γε+ o (ε)
where H (P ) denotes the mean curvature of ∂Ω at P , γ > 0 is a positive constant
given by
(1.13) γ =
N + 1
|w′ (|z|)|m zN dz.
Our goal in this paper is to locate the position on ∂Ω where the global maximum
point Pε of uε in Ω approaches, provided ε is sufficiently small. For the case m = 2,
Ni and Takagi [18] located the peak by linearizing the equation d∆u−u+f (u) = 0
around the ground state w. But this method fails for our problem with m 6= 2 due
to the strong nonlinearity of the m-Laplacian operator ∆mu = div(|∇u|m−2 ∇u).
So we have to use the intrinsic variational method created by Del Pino and Felmer
in [2] to attack it. We also give a complete proof of the exponential decay of
the least-energy solution uε. We remark that our proof is complete and does not
require the non-degeneracy of the unique radial least energy solution w as stated
in Proposition 1.1, and hence it is different from Ni’s and Takagi’s work [17]. Now
our results can be stated as follows:
Theorem 1.2. Under hypotheses (H2)–(H6), let uε be a least-energy solution of
(1.1) and P̃ε ∈ ∂Ω with dist(Pε, P̃ε) = dist(Pε, ∂Ω). Then as ε→ 0+, after passing
to a sequence P̃ε approaches P̄ ∈ ∂Ω with
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 5
(ii) H
= max
H (P ), where H (P ) denotes the mean curvature of ∂Ω at P
as stated before, and moreover
(iii) the associated critical value cε can be estimated as ε→ 0+ as follows :
(1.14) cε = ε
c∗ − (N − 1)H
γε+ o (ε)
where c∗, γ are as stated in Theorem 1.1.
The organization of this paper is as follows: In Section 2, we will prove some
lemmas which will be used in proving Theorem 1.2. The proof of Theorem 1.2 will
be given in Section 3.
2. Some lemmas and exponential decay of uε
First we prove the following lemma related to exponential decay of the least-
energy solution uε.
Lemma 2.1. Let ε be sufficiently small and that the least-energy solution uε
achieves its global maximum at some point Pε. Then there exist two positive con-
stants c3 and c4 independent of uε or ε such that
(2.1)
uε (x) ≤ c3 exp {−c4 |x− Pε| /ε}
|∇uε(x)| ≤ c3ε−1 exp{−c4|x− Pε|/ε}.
Before beginning to prove this lemma, we give a remark on it.
Remark 2.1. For the case m = 2, under the assumption of non-degeneracy of the
linearized operator ∆− 1 + f ′ (w), where w is the unique ground state of (1.3), Ni
and Takagi [18] showed that uε (x) can be written as
(2.2) uε (x) = w (x) + εφ1 (x) + o (ε)
and φ1 (x) enjoys the exponential-decay property ([18]). Clearly we cannot derive
exponential decay of uε (x) as stated in Lemma (2.1) from (2.2) even though both
w (x) and εφ1 (x) have exponential decay property.
Proof of Lemma 2.1. Since ∂Ω is a smooth compact submanifold of RN , it follows
from the tubular neighborhood theorem [10] that there exists a constant ω (Ω) > 0
which depends only on Ω such that ΩI =
x ∈ Ω, d (x, ∂Ω) < ω (Ω)
is diffeomor-
phic to the inner normal bundle
I = {(x, y) : x ∈ ∂Ω, y ∈ (−ω (Ω) , 0] νx} ,
here νx is the unit outer normal of ∂Ω at x, and the diffeomorphism is defined as
follows: ∀x ∈ ΩI , there exists an unique x̂ ∈ ∂Ω such that d (x, x̂) = d (x, ∂Ω) ,
then Φ∗ : x −→ (x̂,−d (x, x̂) νx̂) . Moreover this diffeomorphism satisfies Φ∗|∂Ω =
Identity. Similarly, let ΩO =
x ∈ RN \ Ω, d(x, ∂Ω) < ω (Ω)
. Then ΩO is diffeo-
morphic to the outer normal bundle
O = {(x, y) : x ∈ ∂Ω, y ∈ [0, ω (Ω)) νx} ,
and the diffeomorphism is given as follows. ∀x ∈ ΩO, there exists an unique
x̄ ∈ ∂Ω such that d(x, x̄) = d (x, ∂Ω) , and then Φ# : x −→ (x̄, d (x, x̄) νx̄)
and Φ#|∂Ω = Identity. Note that (∂Ω)NI is clearly diffeomorphic to (∂Ω)
O via
the following reflection Φ∗ : (∂Ω)
I −→ (∂Ω)
O defined by Φ
∗ ((x, y)) = (x,−y) .
6 YI LI AND CHUNSHAN ZHAO
Therefore, Φ = Φ−1∗ ◦ Φ∗−1 ◦ Φ# : ΩO −→ ΩI is the desired diffeomorphism and
Φ|∂Ω = Identity. Moreover, if we let x = Φ(z) = (Φ1(z), · · · ,ΦN(z)) , z ∈ ΩO,
and z = Ψ(x) = Φ−1(x) = (Ψ1(x), · · · ,ΨN(x)) , x ∈ ΩI , gij =
gij =
(Φ (z)) , we have gij |∂Ω = gij |∂Ω = δij with δij being the Kro-
necker symbol. Denote G =
and A = G− I with I being the N ×N identity
matrix, g(x) = det (gij) and ûε(x) = uε (Φ (x)) for x ∈ ΩO. Then ûε(x) satisfies the
following equations:
εmLûε −
gûm−1ε +
gf (ûε) = 0, ûε > 0 in ΩO
= 0, on ∂Ω,
where
Lûε =
s,l=1
∇ûεG (∇ûε)T
g (∇ûε)G
where Tr means taking the trace of a square matrix.
For 0 < γ̃ ≤ ω(Ω), let ΩOγ̃ =
x ∈ ΩO, d(x, ∂Ω) < γ̃
. We know ‖A‖C0 can be
made arbitrarily small by making γ̃ sufficiently small. Next we define
ūε =
uε(x), x ∈ Ω
ûε(x), x ∈ ΩO,
g̃ij =
δij , x ∈ Ω
gij , x ∈ ΩO,
g̃ij =
δij , x ∈ Ω
gij , x ∈ ΩO,
and Ã(x, ξ) =
Ã1(x, ξ), · · · , ÃN (x, ξ)
for ξ = (ξ1, · · · , ξN ) with
Ãi(x, ξ) =
s,l=1
g̃slξsξl
g̃ijξj
and g̃ = det(g̃ij), B(x, u) =
−um−1 + f(u)
. Then ūε(x) satisfies
(2.3) εm div
Ã(x,∇ūε)
+B(x, ūε) = 0 in Ω
in the weak sense.
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 7
For any ball Br(x0) ⊂ Ω
ΩO with radius r and center x0 ∈ Ω, let ρ = |x− x0|.
Then for any smooth increasing function φ = φ(ρ) we have
(∇φ)G(∇φ)T
g(∇φ)G
∣∇φ(I +A)(∇φ)T
det(I +A)−1∇φ(I +A)
= |∇φ|m−2 ∇φ+
∣∇φ(I + tA)(∇φ)T
det(I + tA)−1∇φ(I + tA)
= |∇φ|m−2 ∇φ
∣∇φ(I + tA)(∇φ)T
(∇φ)A(∇φ)T
det(I + tA)−1∇φ(I + tA)dt
∣∇φ(I + tA)(∇φ)T
det(I + tA)−1
det(I + tA)−1
∇φ(I + tA)dt
∣∇φ(I + tA)(∇φ)T
det(I + tA)−1(∇φ)Adt.
Therefore
(2.4)
(∇φ)G(∇φ)T
g(∇φ)G
3(N − 1)
φ′ +K
by taking γ̃ sufficiently small, here K > 0 is a constant depending only on Ψ, hence
only on Ω and φ′ =
dφ(ρ)
From now on γ̃ = γ̃(Ω) is fixed such that (i) 3
≤ √g ≤ 5
, (ii) (2.4) holds for
any smooth increasing radial function φ(ρ) and (iii) 3
|ξ|m ≤ Ã(x, ξ) · ξ ≤ 5
|ξ|m for
any ξ = (ξ1, · · · , ξN ). Denote Ωγ̃ = Ω ∪ ΩOγ̃ .
Let Ωε =
(Ω− Pε) and uε(x) = uε(Pε + εx) for x ∈ Ωε. Then uε is a solution
to the following problem:
(2.5)
ε − (uε)m−1 + f(uε) = 0, uε > 0 in Ωε
= 0, on ∂Ωε,
where n is the unit outer normal of ∂Ωε. Similarly, let Ω
Ωγ̃ − Pε
ūε(x) = ūε(Pε + εx) for x ∈ Ωγ̃ε . Since ūε converges to the unique radial least-
energy solution w of (1.3) in C1loc(R
N ) ∩W 1,m(RN ) as ε → 0+ (see the proof of
Theorem 1.2 of [13]) and w satisfies:
(i) w is radial, i.e.,w(x) = w(|x|) = w(r) > 0
(ii) lim
w(r)r
m(m−1) e(
m−1 )
m r = C0 > 0
(see Theorem 1 of [11]) which yields w(r) ≤ κe−µr for a constant κ > 0 and
. First we fix a constant η > 0 such that 1
tm−1 > f(t) for t ∈
(0, η]. From hypothesis (H5) it follows that such an η exists. Then there exist
ε0 > 0 sufficiently small and R0 sufficiently large such that 4κ exp{−µR0} < η and
8 YI LI AND CHUNSHAN ZHAO
‖ūε − w‖C0(BR0(0)∩Ωε) ≤ κ exp{−µR0}, which yields
uε|(∂BR0 (0))∩Ωε ≤ 2κ exp{−µR0}.
Note that
ε − 7
(uε)m−1 = 1
(uε)m−1 − f(uε) > 0 in Ωε \BR0(0),
= 0 on ∂Ωε \BR0(0),
uε ≤ 2κ exp{−µR0} on ∂BR0(0) ∩ Ωε.
Then we have
uε(x) ≤ 2κ exp{−µR0}, for x ∈ Ωε \BR0(0)
due to the strong maximum principle ([22]). We get by scaling back that
uε|Ω\BεR0 (0) ≤ 2κ exp{−µR0}
(2.6)
uε(x) ≤ w
+ κ exp{−µR0} ≤ κ exp{−
}+ κ exp{−µR0}
≤ 2κ exp{−µ|x|
for x ∈ Ω ∩BεR0(0).
From definition of ūε we know
ūε(x) ≤ 2κ exp{−
µ (|x| − 2dist (Pε, ∂Ω))
} ≤ 4κ exp{−µ|x|
} for x ∈ Ωγ̃∩BεR0(0)
for ε ∈ (0, ε0] with ε0 sufficiently small due to the fact dist(Pε, ∂Ω) = o(ε) as
ε→ 0+. Note that
Ωγ̃\BεR0 (0)
ūε ≤ 4κ exp{−µR0}.
Choice of R0 and γ̃ tells us for any 0 < t ≤ 4κ exp{−µR0}
B(x, t) =
−tm−1 + f(t)
tm−1.
∀x0 ∈ Ω \BεR0(0) and Br(x0) ⊂ Ωγ̃ \BεR0(0), define
φ(x) = φ(ρ) = φ(|x− x0|)
Ωγ̃\BεR0 (0)
where λ∗ > 0 is a constant to be determined later. Simple calculations show that
(i) φ′(ρ) > 0 and
3(N−1)
φ′ +K
(m− 1) (λ∗)m
))m−2
tanh(λ∗ρε )
( λ∗ρε )
+ε (λ∗)
))m−1
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 9
for any 0 < λ∗ ≤ λ̂, where λ̂ > 0 is a small constant depending only on m and
Ω through K. We remark that we have used the fact maxr∈[0,∞)
tanhr
< ∞ for
m ≥ 2. From now on we choose λ∗ = λ̂.
Therefore we have
εm div(Ã(x,∇ūε))−
ūm−1ε ≥ 0 in Br(x0),
εm div(Ã(x,∇φ))−
m−1 ≤ 0 in Br(x0).
Clearly
φ|∂Br(x0) ≥ ūε|∂Br(x0).
Then from the Comparison Theorem (Theorem 10.1 of [19]) it follows that
φ(x) ≥ ūε(x) in Br(x0).
In particular, φ(x0) ≥ ūε(x0). Thus we get
uε(x0) ≤
Ωγ̃\BεR0(0)
exp{−λ∗r
Choosing r = d
x0, ∂
Ωγ̃ \BεR0(0)
we get
uε(x0) ≤ 4κ exp{−µR0 −
} ≤ 2κ exp{− λ̃(εR0 + r)
with λ̃ = min{µ, λ∗}. Note that x0 belongs to one of the following two cases:
(i) d
x0, ∂
Ωγ̃ \BεR0(0)
= d (x0, ∂BεR0(0)) ,
(ii) d
x0, ∂
Ωγ̃ \BεR0(0)
x0, ∂Ω
For case (i) we have d(x0, Pε) ≤ εR0 + r and therefore
(2.7) uε(x0) ≤ 4κ exp{−
λ̃d(x0, Pε)
For case (ii) we have r ≥ γ̃ and thus
(2.8)
uε(x0) ≤ 4κ exp{−λ̃
εR0 + r
} ≤ 4κ exp{− λ̃γ̃
≤ 4κ exp{−λ̃ γ̃
diam(Ω)
· d(x0, Pε)
Combining (2.6), (2.7) and (2.8) together and letting c̃3 = 4κ, c̃4 = min{µ, λ̃, λ̃γ̃diam(Ω)}
yields
(2.9) uε(x) ≤ c̃3 exp{−
c̃4|x− Pε|
Next we show the estimate for |∇uε| holds. First from (2.5) it follows that
(2.10) ∆mu
ε = (uε)
m−1 − f(uε), uε > 0 in Ωε
For x ∈ Ωε and dist(x, ∂Ωε) ≥ 1, consider (2.10) in the unit ball centered at x, i.e.,
B1(x). Then by an C
1,α estimate (see [21], for example) there exists two constants
10 YI LI AND CHUNSHAN ZHAO
C > 0 and α∗ ∈ (0, 1) which are independent of ε such that
(2.11)
‖uε‖C1,α∗ (B 1
(x)) ≤ C
‖uε‖L∞(B1(x)) + ‖ (u
m−1 − f(uε)‖
L∞(B1(x))
≤ c∗3 exp{−c∗4|x− Pε|},
where we have used (2.9) and the fact that uε(x) = uε(Pε + εx) for x ∈ Ωε.
Especially we have
(2.12) |∇uε(x)| ≤ c∗3 exp{−c∗4|x− Pε|},
for x ∈ Ωε and dist(x, ∂Ωε) ≥ 1. For x ∈ Ωε with dist(x, ∂Ωε) < 1. Let x0 ∈ ∂Ωε
be a point such that dist(x, x0) = dist(x, ∂Ωε) and consider ū
ε(x) = ūε(Pε + εx)
in B2(x0), the ball of radius 2 centered at x0, then from (2.3) it follows that ū
satisfies
(2.13) div
Ã(Pε + εx,∇ūε)
+B(Pε + εx, ū
ε) = 0 in B2(x0)
in the weak sense. Then applying an C1,α estimate (see [21], for example) again
yields as above that there exists two constants C > 0 and α∗ ∈ (0, 1) which are
independent of ε such that
‖ûε‖C1,α∗(B1(x0)) ≤ C
‖ûε‖L∞(B2(x0)) + ‖B(Pε + εx, û
L∞(B2(x0))
≤ c∗3 exp{−c∗4|x− Pε|}
by adjusting c∗3 and c
4 if it is necessary. Especially we have
(2.14) |∇uε(x)| ≤ c∗3 exp{−c∗4|x− Pε|},
Thus combining (2.11) and (2.14) together and scaling back we have for x ∈ Ω
|∇uε(x)| ≤ c∗3ε−1 exp{−c∗4
|x− Pε|
Proof of Lemma 2.1 is completed by letting c3 = max{c̃3, c∗3} and c4 = min{c̃4, c∗4}.
Remark 2.2. Our proof of the Lemma 2.1 with necessary minor modifications also
works well for elliptic systems.
Next we present a lemma related to extensions of uε.
Lemma 2.2. There exists a C1-extension ũε of uε which has compact support in
N and satisfies
(ii) ‖ũε‖W 1,m(RN ) ≤ c5 ‖uε‖W 1,m(Ω) and ‖ũε‖C1(RN ) ≤ c5 ‖uε‖C1(Ω̄),
(iii) ũε also has the exponential-decay property as stated in Lemma 2.1, i.e.,
there exists an absolute constant λ ≥ 1 such that
(2.15)
0 ≤ ũε ≤ c3λ exp
|x− Pε|
|∇ũε(x)| ≤ c3λε−1 exp{−
|x− Pε|
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 11
(iv) there exists a positive constant δ̃ = δ̃ (Ω) such that for any P ∈ ∂Ω,
ũε|B
(P )\Ω is the reflection of uε through ∂Ω.
Proof. Let d̃ = d
∂Ω, ∂Ωγ̃
and 0 ≤ ̺(x) ≤ 1 be a smooth cut-off function such
that ̺(x) ≡ 1 for x ∈ {x ∈ RN , d(x,Ω) ≤ d̃
} and ̺(x) ≡ 0 for x ∈ RN \
Then ũε = ̺ūε satisfies (ii), (iii) and (iv) automatically. The proof of this lemma
is completed. �
Similar to energy density introduced in [2], we define the energy density associ-
ated with (1.1) as follows:
E (w, y′) =
(|∇w|m + wm)− F (w)
(y′, 0) for y′ ∈ RN−1.
Then we have the following lemma.
Lemma 2.3. Let G be a C2 function in a neighborhood of the origin of RN−1.
i,j=1
Gij (0) yiyjE (w, y
′) dy′ = 2∆G (0) γ,
where γ is the constant defined in (1.13), and y′ = (y1, . . . , yN−1), and
Gij (0) =
∂yi∂yj
(0) .
Proof. In Lemma 2.4 of [13], we showed that
(2.16) γ =
(|∇w|m) + wm − F (w)
zN dz.
Next we introduce the polar coordinates
z1 = r sin θN−1 sin θN−2 · · · sin θ2 sin θ1,
z2 = r sin θN−1 sin θN−2 · · · sin θ2 cos θ1,
z3 = r sin θN−1 sin θN−2 · · · cos θ2,
... ,
zN = r cos θN−1,
and notice that
(r, θ1, . . . , θN−1) | r > 0, 0 ≤ θ1 < 2π,
0 ≤ θj < π for j = 2, . . . , N − 2, and 0 ≤ θN−1 <
and that
dz = rN−1 sin θ2 sin
2 θ3 · · · sinN−2 θN−1 dr dθ1 · · · dθN−1.
After elementary computations one obtains
(2.17) γ =
|w′ (r)|m + wm (r)
− F (w (r))
rN dr · ωN−2,
where ωN−2 is the volume of the unit ball in R
N−2. Here we used the fact that w
is radially symmetric.
12 YI LI AND CHUNSHAN ZHAO
Using the radial symmetry of w again, we obtain
i,j=1
Gij (0) yiyjE (w, y
′) dy′(2.18)
Gii (0) y
iE (w, y
′) dy′
Gii (0) ·
N − 1
|y′|2E (w, y′) dy′
= ∆G (0) ·
E (w, r) rN dr · ωN−2,
where E (w, r) = (1/m)
|w′ (r)|m + wm (r)
− F (w (r)) . Comparing (2.17) and
(2.18) yields
i,j=1
Gij (0) yiyjE (w, y
′) dy′ = 2∆G (0) γ.
The proof of Lemma 2.3 is completed. �
3. Proof of Theorem 1.2
With the help of the lemmas in Section 2, now we can give the proof of Theorem
Proof of Theorem 1.2. Since as ε→ 0+, Pε → ∂Ω at the rate of o(ε), it follows that
d(Pε, P̃ε)/ε → 0, where P̃ε ∈ ∂Ω is the closest point on ∂Ω to Pε. then by passing
to a sequence, P̃ε → P̄ ∈ ∂Ω. After an ε-dependent rotation and translation,
we may assume that P̃ε is at the origin and Ω can be described in a fixed cubic
neighborhood V of P̄ as the set
{ (x′, xN ) | xN > ψε (x′) } with x′ = (x1, . . . , xN−1) ,
where ψε is smooth, ψε (0) = 0, ∇ψε (0) = 0. Furthermore, we may assume that
ψε converges locally in the C
2 sense to ψ, a corresponding parametrization at P̄ .
Note that since P̃ε is the origin, so we have Pε/ε → 0 as ε → 0+. Thus we have
ũε(x) = ũε(εx) = ũε
x− Pε
→ w(x) in C1loc
as ε → 0+. From the
characterization of cε = Jε (uε) in Section 1, we have
ε−NJε (uε) ≥ ε−NJε (tuε) = IΩε (tuε)
for all t > 0. Hereinafter
IΩ∗ (v) =
(|∇v|m + |v|m) dx−
F (v) dx.
IΩε (tu
ε) = IΩε (tũ
ε) ≥ I
(tũε) + I(Ωε∩Vε)\RN+
(tũε)− I(RN+∩Vε)\Ωε (tũ
ε)(3.1)
= I + II− III,
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 13
with Vε =
V. Let us choose t = tε so that IRN
(tũε) maximizes in t. Then from
the definition of C∗ in (1.10), equality (1.11) and Lemma 2.2 it follows that
I = I
(tεũ
ε) ≥ c∗
e−c6/ε
for some constant c6 > 0 independent of ε. Next we give an estimate of tε.
Lemma 3.1. There is a unique tε ∈ (0,∞) such that
tmε (|∇ũε|
+ (ũε)m) dx −
F (tεũ
ε) dx
= sup
tm (|∇ũε|m + (ũε)m) dx−
F (tũε) dx
and moreover
(3.2) tε = 1 + o (1) as ε→ 0+.
Proof. Under assumption (H5), the existence and uniqueness of tε can be proved
similarly to the proof of Lemma 2.1 of [13]. Here we only need show (3.2). Let
(3.3) hε (t) =
(|∇ũε|m + (ũε)m) dx−
F (tũε) dx.
(3.4)
h′ε (t) = t
(|∇ũε|m + (ũε)m) dx−
ũεf (tũε) dx
= tm−1
(|∇w|m + wm) dx−
wf (tw) dx+ o(1),
here we have used the exponential decay of ũε in Lemma 2.2, exponential decay of
w and ũε → w in C1loc
as ε → 0+. Moreover the term o (1) → 0 uniformly
in t on each compact interval as ε → 0+. (3.3) tells us hε(1) = 12c∗ + o(1), which
yields that tε is bounded and away from 0. Also from (3.4) it follows that
(3.5)
h′ε (t) = t
wf (w) dx−
w f (tw) dx+ o (1)
= tm−1
f (w)
− f (tw)
dx+ o (1) .
Therefore at t = tε we have
(3.6)
f (w)
− f (tεw)
(tεw)
dx = o (1) .
Since f(t)/tm−1 is strictly increasing (see (H5)) it follows from (3.6) that tε =
1 + o (1) . The proof of Lemma 3.1 is completed. �
14 YI LI AND CHUNSHAN ZHAO
Proof of Theorem 1.2 continued. Using again the exponential decay of uε in
Lemma 2.1 and the expansion of tε in Lemma 3.1, we obtain
−II = −
(RN−1×{0})∩Vε
dy′(3.7)
(ψε(εy′))
tmε (|∇ũε|
+ (ũε)
)− F (tεũε)
(y′, yN ) dyN
= − (1 + o (1))
(RN−1×{0})∩(Ωε∩Vε)
(ψε(εy′))
(|∇uε|m + (uε)m)− F (uε)
(y′, yN) dyN .
Similarly,
(3.8) III = (1 + o (1))
Vε∩(RN−1×{0})
(ψε(εy
(|∇ũε|m + (ũε)m)− F (ũε)
(y′, yN) dyN .
In above a+ = max{a, 0}, a− = min{a, 0}. Since ψε (0) = 0, ∇ψε (0) = 0 and
ψε converges in the C
2 local sense to ψ, and ũε → w in the C1 local sense in RN
with uniform exponential decay with respect to ε, it follows from the dominated
convergence theorem that
(−II + III)
i,j=1
ψij (0) yiyj
(|∇w|m + wm)− F (w)
(y′, 0) dy′
= ∆ψ (0) γ = (N − 1)H
γ (by Lemma 2.3).
Thus we have
cε ≥ εN
c∗ − (N − 1)H
γε+ o (ε)
But (1.12) in Theorem 1.1 tells us
cε ≤ εN
c∗ − (N − 1) max
H (P ) γε+ o (ε)
Therefore we get
(ii) H
= max
H (P ), which is (ii) of Theorem 1.2,
(iii) cε = ε
c∗ − (N − 1)H
γε+ o (ε)
as ε→ 0+,
which is part (iii) of Theorem 1.2. The proof of Theorem 1.2 is completed. �
Acknowledgement. The authors want to give their thanks to anonymous
referee for some helpful comments.
LOCATING THE PEAKS OF LEAST-ENERGY SOLUTIONS 15
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[19] P. Pucci and J. Serrin, The strong maximum principle revisted, J. Differential Equations 196
(2004), no. 1, 1-66. MR 2004k:35033
[20] J. Serrin and M.-X. Tang, Uniqueness of ground states for quasilinear elliptic equations,
Indiana Univ. Math. J. 49 (2000), no. 3, 897–923. MR 2002d:35072
[21] P. Tolksdorf, Regularity for a more general class of quasilinear elliptic equations, J. Differ-
ential Equations 51 (1984), no. 1, 126–150. MR 85g:35047
[22] J. Vazquez, A strong maximum principle for some quasilinear elliptic equations, Appl. Math.
Optim. 12 (1984), no. 3, 191–202. MR 86m:35018
[23] J. Wei, On the boundary spike layer solutions to a singularly perturbed Neumann problem,
J. Differential Equations 134 (1997), 104–133. MR 98e:35076
Department of Mathematics, The University of Iowa, Iowa City, IA 52242
16 YI LI AND CHUNSHAN ZHAO
Department of Mathematics, Hunan Normal University, Changsha, Hunan
E-mail address: [email protected]
Department of Mathematical Sciences, Georgia Southern University, Statesboro,
GA 30460
E-mail address: [email protected]
1. Introduction and statement of results
2. Some lemmas and exponential decay of u0=x"0122
3. Proof of Theorem ??
References
|
0704.0403 | Review: Semiconductor Quantum Light Sources | Microsoft Word - preprint.doc
Preprint version of Nature Photonics 1, 215 (2007)
Review: Semiconductor Quantum Light Sources
Andrew J Shields
Toshiba Research Europe Limited, 260 Cambridge Science Park, Cambridge CB4 0WE, UK
Abstract
Lasers and LEDs display a statistical distribution in the number of photons emitted in a given time interval. New
applications exploiting the quantum properties of light require sources for which either individual photons, or pairs, are
generated in a regulated stream. Here we review recent research on single-photon sources based on the emission of a
single semiconductor quantum dot. In just a few years remarkable progress has been made in generating
indistinguishable single-photons and entangled photon pairs using such structures. It suggests it may be possible to
realise compact, robust, LED-like semiconductor devices for quantum light generation.
Applications of Quantum Photonics
Applying quantum light states to photonic applications allows functionalities that are not possible using ‘ordinary’
classical light. For example, carrying information with single-photons provides a means to test the secrecy of optical
communications, which could soon be applied to the problem of sharing digital cryptographic keys.1 2 Although secure
quantum key distribution systems based on weak laser pulses have already been realised for simple point-to-point links,
true single-photon sources would improve their performance.3 Furthermore, quantum light sources are important for
future quantum communication protocols such as quantum teleportation. 4 Here quantum networks sharing
entanglement could be used to distribute keys over longer distance or through more complex topologies.5
A natural progression would be to use photons for quantum information processing, as well as communication. In this
regard it is relatively straightforward to encode and manipulate quantum information on a photon. On the other hand,
single-photons do not interact strongly with one-another, a prerequisite for a simple photon logic gate. In linear optics
quantum computing67 (LOQC) this problem is solved using projective measurements to induce an effective interaction
between the photons. Here triggered sources of single-photons and entangled pairs are required as both the qubit
carriers, as well as auxiliary sources to test the successful operation of the gates. Although the component requirements
for LOQC are challenging, they have recently been relaxed significantly by new theoretical schemes. 7 Quantum light
states are also likely to become increasingly important for various types of precision optical measurement.8
For these applications we would ideally like light sources which generate pure single-photon states “on demand” in
response to an external trigger signal. Key performance measures for such a source are the efficiency, defined as the
fraction of photons collected into the experiment or application per trigger, and the second order correlation function at
zero delay, see text box. The latter is essentially a measure of the two-photon rate compared to a classical source with
random emission times of the same average intensity. In order to construct applications involving more than one
photon, it is also important that photons emitted from the source (at different times), as well as those from different
sources, are otherwise indistinguishable.
In the absence of a convenient triggered single-photon source, most experiments in quantum optics rely on non-linear
optical processes for generating quantum light states. Optically pumping a crystal with a χ(2) non-linearity has a finite
probability of generating a pair of lower energy photons via parametric down conversion. This may be used to prepare
photon pairs with time-bin entanglement,9 entangled polarisations,1011 or alternatively single-photon states ‘heralded’ by
the second photon in the pair.12 A χ(3) non-linearity in a semiconductor has also been used to generate entangled
pairs.13 As these non-linear processes occur randomly, there is always a finite probability of generating two pairs that
increases with pump power. As double pairs degrade the fidelity of quantum optical gates, the pump laser power must
be restricted to reduce the rate of double pairs to an acceptable level, which has a detrimental effect upon the efficiency
of the source.14 This means that although down-conversion sources continue to be highly successful in demonstrating
few photon quantum optical gates, scaling to large numbers may be problematic. Solutions have been proposed based
on switching multiple sources,15 or storing photons in a switched fibre loop.16
Ideally we would like a quantum light source that generates exactly one single-photon, or entangled-pair, per excitation
trigger pulse. This may be achieved using the emission of a single quantum system. After relaxation, a quantum
system is by definition no longer excited and therefore unable to re-emit. Photon anti-bunching, the tendency of a
quantum source to emit photons separated in time, was first demonstrated in the resonance fluorescence of a low density
vapour of Na atoms,17 and subsequently for a single ion.18
Quantum dots are often referred to as “artificial atoms”, as their electron motion is quantised in all three spatial
directions, resulting in a discrete energy level spectrum, like that of an atom. They provide a quantum system which
can be grown within robust, monolithic semiconductor devices and can be engineered to have a wide range of desired
properties. In the following we review recent progress towards the realisation of a semiconductor technology for
quantum photonics. An excellent account of the early work can be found in Ref. 19. Space restrictions limit discussion
of work on other quantised systems. For this we refer the reader to the comprehensive review in Ref 20.
Optical Properties of Single Quantum Dots
Nano-scale quantum dots with good optical properties can be fabricated using a natural growth mode of strained layer
semiconductors.21 When InAs is deposited on GaAs it initially grows as a strained two-dimensional sheet, but beyond
some critical thickness, tiny islands like those shown in Fig.1a form in order to minimize the surface strain.
Overgrowth of the islands leads to the coherent incorporation of InxGa1-xAs dots into the crystal structure of the device,
as can be seen in the cross-sectional image of Fig.1c. The most intensively studied are small InAs dots on GaAs
emitting around 900-950nm at low temperatures, which can be conveniently measured with low noise Si single photon
detectors.
A less desirable feature of the self-organising technique is that the dots form at random positions on the growth surface.
However, recently considerable progress has been made on controlling the dot position (Fig.1b) within the device
structure by patterning nanometer sized pits on the growth surface.2223
As InGaAs has a lower energy bandgap than GaAs, the quantum dot forms a potential trap for electrons and holes. If
sufficiently small, the dot contains just a few quantised levels in the conduction and valence bands, each of which holds
two electrons or holes of opposite spin. Illumination by a picosecond laser pulse excites electrons and holes which
rapidly relax to the lowest lying energy states either side of the bandgap. A quantum dot can thus capture two electrons
and two holes to form the biexciton state, which decays by a radiative cascade, as shown schematically in Fig.2a. One
of the trapped electrons recombines with one of the holes and generates a first photon (called the biexciton photon, X2).
This leaves a single electron-hole pair in the dot (the exciton state), which subsequently also recombines to generate a
second (exciton, X) photon. The biexciton and exciton photons have distinct energies, as can be seen in the low
temperature photoluminescence spectrum of Fig.2a, due to the different Coulomb energies of their initial and final
states. Often a number of other weaker lines can also be seen due to recombination of charged excitons which form
intermittently when the dot captures an excess electron or hole.24 Larger quantum dots, with several confined electron
and hole levels, have a richer optical signature due to the large number of exciton complexes that can be confined.
High resolution spectroscopy reveals that the X2 and X transitions of a dot are in fact both doublets with linearly
polarised components parallel to the [110] and [1-10] axes of the semiconductor crystal, labelled here H and V,
respectively.2526 The origin of this polarisation is an asymmetry in the electron-hole exchange interaction of the dot
which produces a splitting of the exciton spin states. The asymmetry derives from an elongation of the dot along one
crystal axis and in-built strain in the crystal. It mixes the exciton eigenstates of a symmetric dot with total z-spin Jz =
+1 and -1 into symmetric and anti-symmetric combinations, which couple to two H or two V polarised photons,
respectively, as shown in Fig.2.
The exciton state of the dot has a typical lifetime of ~1ns, which is due purely to radiative decay. As this is much
longer than the duration of the exciting laser pulse, or the lifetime of the photo-excited carrier population in the
surrounding semiconductor, only one X photon can be emitted per laser pulse. This can be proven, as first reported27 by
Peter Michler, Atac Imamoglu and their colleagues in Santa Barbara, by measuring the second order correlation
function, g(2)(τ) of the exciton photoluminescence,2829 see text box. In fact each of the exciton complexes of the dot
generates at most one photon per excitation cycle, which allows single-photon emission from also the biexciton or
charged exciton transitions.30
Cross-correlation measurements313233 between the X and X2 photons confirm the time correlation expected for the
cascade in Fig.2a, ie the X photon follows the X2 one. Indeed the shape of the cross-correlation function for both CW
and pulsed excitation can be accurately described with a simple rate equation model and the experimentally measured X
and X2 decay rates.
34
Semiconductor Microcavities
A major advantage of using self-assembled quantum dots for single-photon generation is that they can be easily
incorporated into cavities using standard semiconductor growth and processing techniques. Cavity effects are useful for
directing the emission from the dot into an experiment or application, as well as for modifying the photon emission
dynamics. 3536 Purcell37 predicted enhanced spontaneous emission from a source in a cavity when its energy coincides
with that of the cavity mode, due to the greater density of optical states to emit into. For an ideal cavity, in which the
emitter is located at the maximum of the electric field with its dipole aligned with the local electric field, the
enhancement in decay rate is given by Fp = (3/4π
2) (λ/n)3 Q/V, where Q is the quality factor, a measure of the time a
photon is trapped in the cavity, and V is the effective mode volume. Thus high photon collection efficiency, and
simultaneously fast radiative decay, requires small cavities with highly reflecting mirrors and a high degree of structural
perfection. However, without controlling the location of the dot in the cavity, as discussed below, it may be difficult to
achieve the full enhancement predicted by the Purcell formula.
Figure 3 shows images of some of the single quantum dot cavity structures that have proven most successful. Pillar
microcavities, formed by etching cylindrical pillars into semiconductor Bragg mirrors placed either side of the dot layer,
have shown large Purcell enhancements and have a highly directional emission profile, thus making good single-photon
sources.38394041 Purcell factors of around 6 have been measured directly,4041 through the rate of cavity-enhanced radiative
decay compared to that of a dot without cavity, implying a coupling to the cavity mode of β=Fp/(1+Fp)>85%, if we
assume the leaky modes are unaffected by the cavity. However, the experimentally determined photon collection
efficiency, which is a more pertinent parameter for applications, is typically ~10%, due the fact that not all the cavity
mode can be coupled into an experiment and scattering of the mode by the rough pillar edges. We can expect that the
photon collection efficiency will increase with improvements to the processing technology or new designs of
microcavity.
Another means of forming a cavity is to etch a series of holes in a suspended slab of semiconductor, so as to form a
lateral variation in the refractive index which creates a forbidden energy gap for photonic modes in which light cannot
propagate.42 Photons can then be trapped in a central irregularity in this structure: usually an unetched portion of the
slab. Such photonic bandgap defect cavities have been fabricated in Si with Q values approaching 106.4344 High quality
active cavities have also been demonstrated in GaAs containing InAs quantum dots. 45464748 A radiative lifetime of 86 ps,
corresponding to a Purcell factor of Fp~12, has been reported.
47 Very recently a lifetime of 60ps was measured for a
cavity in the strong coupling regeme.48
If the Q-value is sufficiently large, the system enters the strong coupling regime where the excitation oscillates
coherently between an exciton in the dot and a photon in the cavity. The spectral signature of strong coupling, an anti-
crossing between the dot line and the cavity mode, has been observed for quantum dots in pillar microcavities,49
photonic bandgap defect cavities,50 microdisks51 and microspheres.52 It has been demonstrated for atom cavities that
strong coupling allows the deterministic generation of single-photons.5354 Single-photon sources in the strong coupling
regime can be expected to have very high extraction efficiencies and be time-bandwidth limited.55 Encouragingly
single-photon emission has been reported recently for a dot in a strongly coupled pillar microcavity. 56
Another interesting recent development is the ability to locate a single quantum dot within the cavity, as this ensures the
largest possible coupling and removes background emission, as well as other undesirable effects, due to other dots in
the cavity. Above we discussed techniques to control the dot position on the growth surface. The other way is to
position the cavity around the dot. One technique combines micro-photoluminescence spectroscopy to locate the dot
position, with in-situ laser photolithography to pattern markers on the wafer surface.57 An alternative involves growing
a vertical stack of dots so that their location can be revealed by scanning the wafer surface, 58 as shown in Fig.3.
Recently this technique has allowed larger coupling energies for a single dot in a photonic bandgap defect cavity.48
Photon Indistinguishability
Cavity effects are important for rendering different photons from the source indistinguishable, which is essential for
many applications in quantum information. When two identical photons are incident simultaneously on the opposite
input ports of a 50/50 beamsplitter, they will always exit via the same output port, 59 as shown schematically in Fig.4a.
This occurs because of a destructive interference in the probability amplitude of the final state in which one photon exits
through each output port. The amplitude of the case where both photons are reflected exactly cancels with that where
both are transmitted, due to the π/2 phase change upon reflection, provided the two photons are entirely identical.
Two-photon interference of two single-photons emitted successively from a quantum dot in a weakly-coupled pillar
microcavity was first reported by the Stanford group.60 Fig. 4b shows a schematic of their experiment. Notice the
reduction of the co-incidence count rate measured between detectors in either output port, when the two photons are
injected simultaneously (Fig.4c). The dip does not extend completely to zero, indicating that the two photons sometime
exit the beamsplitter in opposite ports. The measured reduction in co-incidence rate at zero delay of 69%, implies an
overlap for the single-photon wavepackets of 0.81, after correcting for the imperfect single-photon visibility of the
interferometer. Two-photon interference dips of 66% and 75% have been reported by Bennett et al61 and Vauroutsis et
al. 62 Similar results have been obtained for a single dot in a photonic bandgap defect cavity.63
This two-photon interference visibility is limited by the finite coherence time of the photons emitted by the quantum
dot,64 which renders them distinguishable. The depth of the dip in Fig.4c depends upon the ratio of radiative decay time
to the coherence time of the dot, ie R=2τdecay/τcoh. When unity, the coherence time is limited by radiative decay and the
source will display perfect 2-photon interference. The most successful approach thus far has been to extend τcoh by
resonant optical excitation of the dot and reduce τdecay using the Purcell effect in a pillar microcavity, to values R~1.5.
the future higher visibilities may be achieved with a larger Purcell enhancement, using a single dot cavity in the strong-
coupling regime or with electrical gating described in the next section.
A source of indistinguishable single-photons was used by Fattal et al to generate entanglement between post-selected
pairs. 65 66 This involves simply rotating the polarisation of one of the photons incident on the final beamsplitter in
Fig.4a by 90o. By post-selecting the results where the two photons arrive at the beamsplitter at the same time and where
there is one photon in each output arm (labelled 1 and 2), the measured pairs should correspond to the Bell state
ψ− = 1/√2 (¦H1 V2 > - ¦V1 H2 >) Eq.1
Note that only if the two photons are indistinguishable and thus the entanglement is only in the photon polarisation, are
the two terms in Eq1 able to interfere. Analysis of the density matrix published by Fattal et al65 reveals a fidelity of the
post-selected pairs to the state in Eq.1 of 0.69, beyond the classical limit of 0.5. This source of entangled pairs has an
importance difference to that based on the biexciton cascade described below. Post-selection implies that the photons
are destroyed when this scheme succeeds. This is a problem for some quantum information applications such as LOQC,
but could be usefully applied to quantum key distribution.65
Single-Photon LEDs
An early proposal for an electrical single-photon source by Kim et al67 was based upon etching a semiconductor
heterostructure displaying Coulomb blockade. However, the light emission from this etched structure was too weak to
allow the second-order correlation function to be studied. Recently encouraging progress has been made towards the
realisation of a single-photon source based on quantising a lateral electrical injection current.6869 However the most
successful approach so far has been to integrate self-assembled quantum dots into conventional p-i-n doped junctions.
In the first report of electrically-driven single-photon emission by Yuan et al,70 the electroluminescence of a single dot
was isolated by forming a micron-diameter emission aperture in the opaque top contact of the p-i-n diode. Fig.5a shows
an improved emission aperture single-photon LED after Bennett et al, 71 which incorporates an optical cavity formed
between a high reflectivity Bragg mirror and the semiconductor/air interface in the aperture. This structure forms a weak
cavity, which enhances the measured collection efficiency 10-fold compared to devices without a cavity. 72
Single-photon pulses are generated by exciting the diode with a train of short voltage pulses. The second order
correlation function g(2)(τ) of either the X or X2 electroluminescence (Fig.5c) shows the suppression of the zero delay
peak indicative of single-photon emission.71 The finite rate of multi-photon pulses is due mostly to background
emission from layers other than the dot, which is also seen for non-resonant optical excitation. Electrical contacts also
allow the temporal characteristics of the single-photon source to be tailored. By applying a negative bias to the diode
between the electrical injection pulses, Bennett et al73 reduced the jitter in the photon emission time <100ps. This
allowed the repetition rate of the single-photon source to be increased to 1.07GHz (Fig.5d) while retaining good single-
photon emission characteristics (Fig.5e). Electrical gating could provide a technique for producing time-bandwidth-
limited single-photons from quantum dots.
Another promising approach is to aperture the current flowing through the device.7475 This is achieved by growing a
thin AlAs layer within the intrinsic region of the p-i-n junction and later exposing the mesa to wet oxidation in a
furnace, converting the AlAs layer around the outer edge of the mesa to insulating Aluminium oxide. By careful
control of the oxidation time, a µm-diameter conducting aperture can be formed within the insulating ring of AlOx.
Such structures have the advantage of exciting just a single dot within the structure, thereby reducing the amount of
background emission. The oxide annulus also confines the optical mode laterally within the structure, potentially
allowing high photon extraction efficiency.
Altering the nanostructure or materials that comprise the quantum dot allows considerable control over the emission
wavelength and other characteristics. Most of the experimental work done so far has concentrated on small InAs
quantum dots emitting around 900-950nm, as these have well understood optical properties and can be detected with
low noise Si single-photon detectors. On the other hand the shallow confinement potentials of this system means they
emit only at low temperatures. At shorter wavelengths optically-pumped single-photon emission has been demonstrated
at ~350nm using GaN/AlGaN,76 500nm using CdSe/ZnSSe77 and 682nm InP/GaInP78 quantum dot. The former two
systems have been shown to operate at 200K.
It is very important for quantum communications to develop sources at longer wavelengths in the fibre optic
transmission bands at 1.3 and 1.55µm. This may be achieved using InAs/GaAs heterostructures by depositing more
InAs to form larger quantum dots. These larger dots offer deeper confinement potentials than those at 900nm and thus
often display room temperature emission.79 Optically pumped single-photon emission at telecom wavelengths has been
achieved using a number of techniques to prepare low densities of longer wavelength dots, including a bimodal growth
mode in MBE to form low densities of large dots,80 ultra-low growth rate MBE81 and MOCVD.82 Recently, the first
electrically-driven single-photon source at a telecom wavelength has been demonstrated.83
Generation of Entangled Photons
By collecting both the X2 and X photons emitted by the biexciton cascade, a single quantum dot may also be used as a
source of photon pairs. Polarisation correlation measurements on these pairs discovered that the two photons were
classically-correlated with the same linear polarisation.848586 This occurs because the cascade can proceed via one of
two intermediate exciton spin states, as described above and shown in Fig.2a, one of which couples to two H- and the
other two V-polarised photons. The emission is thus a statistical mixture of |HX2HX> and |VX2VX>, although exciton
spin scattering during the cascade (discussed below) ensures there are also some cross-polarised pairs.
The spin splitting87,88 of the exciton state of the dot distinguishes the H and V polarised pairs and prevents the emission
of entangled pairs predicted by Benson et al. 89 If this splitting could be removed, the H and V components would
interfere in appropriately designed experiments. The emitted 2-photon state should then be written as a superposition of
HH and VV, which can be recast in either the diagonal (spanned by D, A) or circular (σ+, σ-) polarisation bases, ie
Φ+ = 1/√2 (¦HX2 HX > + ¦VX2 VX >)
= 1/√2 (¦DX2 DX > + ¦ΑX2 ΑX >)
= 1/√2 (¦σ+X2 σ
X > + ¦σ
X >) Eq.2.
Equal weighting of the HH and VV terms assumes the source to be unpolarised, as indicated by experimental
measurements.
Eq.2 suggests that, for zero exciton spin splitting, the biexciton cascade generates entangled photon pairs, similar to
those seen for atoms.90 Entanglement of the X or X2 photons was recently observed experimentally for the first time by
Stevenson, Young and co-workers,9192 using two different schemes to cancel the exciton spin splitting. An alternative
approach by Akopian et al, 93 using dots with finite exciton splitting, post-selects photons emitted in a narrow spectral
band where the two polarisation lines overlap.
The exciton spin splitting depends on the exciton emission energy, tending to zero for InAs dots emitting close to 1.4eV
and then inverting for higher emission energy. 94 95 These correspond to shallow quantum dots for which the carrier
wavefunctions extend into the barrier material reducing the electron-hole exchange. Zero splitting can be achieved by
either careful control of the growth conditions to achieve dots emitting close to the desired energy, or by annealing
samples emitting at lower energy.94 The exciton spin splitting may be continuously tuned by applying a magnetic field
in the plane of the dot.96 It has been observed that the signatures of entanglement then appear only when the exciton
splitting is close to zero.91 Other promising schemes to tune the exciton splitting are now emerging, including
application of strain97 and electric field.9899
Figure 6a plots polarisation correlations reported by Young et al92 for a dot with zero exciton splitting (by control of the
growth conditions). Pairs emitted in the same cascade (ie zero delay) shows a very striking positive correlation (co-
polarisation) measuring in either, rectilinear or diagonal bases and anti-correlation (cross-polarisation) when measuring
in circular basis. This is exactly the behaviour expected for the entangled state of Eq.2. In contrast, a dot with finite
splitting shows polarisation correlation for the rectilinear basis only, with no correlation for diagonal or circular
measurements, see Figure 6b. The strong correlations seen for all three bases in Fig.6a could not be produced by any
classical light source or mixture of classical sources and is proof that the source generates entangled photons. The
measured92 two-photon density matrix (Fig.6c) projects onto the expected 1/√2 (¦HX2 HX > + ¦VX2 VX >) state with
fidelity (ie probability) 0.702 ± 0.022, exceeding the classical limit (0.5) by 9 standard deviations.
Two processes contribute to the ‘wrongly’ correlated pairs which impair the fidelity of the entangled photon source.
The first of these is due to background emission from layers in the sample other than the dot. This background
emission, which is unpolarised and dilutes the entangled photons from the dot, limited the fidelity observed in the first
report91 of triggered entangled photon pairs from a quantum dot and has been subsequently reduced with better sample
design.92 The second mechanism, which is an intrinsic feature of the dot, is exciton spin scattering during the biexciton
cascade. It is interesting that this process does not seem to depend strongly upon the exciton spin splitting. It may be
reduced by suppressing the scattering using resonant excitation or alternatively using cavity effects to reduce the time
required for the radiative cascade.
Outlook
The past several years have seen remarkable progress in quantum light generation using semiconductor devices.
However, despite considerable progress many challenges still remain. The structural integrity of cavities must continue
to improve, thereby enhancing quality factors. This, combined with the ability to reliably position single dots within the
cavity, will further enhance photon collection efficiencies and the Rabi energy in the strong coupling regime. It is also
important to realise all the benefits of these cavity effects in more practical electrically-driven sources. Meanwhile
bandstructure engineering of the quantum dots will allow a wider range of wavelengths to be accessed for both single
and entangled photon sources, as well as structures that can operate at higher temperatures. Techniques for fine tuning
the characteristics of individual emitters will also be important.
One of the most interesting aspects of semiconductor quantum optics is that we may be able to use quantum dots not
only as quantum light emitters, but also as the logic and memory elements which are required in quantum information
processing. Although LOQC is scalable theoretically, quantum computing with photons would be much easier with a
useful single-photon non-linearity. Such non-linearity may be achieved with a quantum dot in a cavity in the strong
coupling regime. Encouragingly strong coupling of a single quantum dot with various type of cavity has already been
observed in the spectral domain. Eventually it may even be possible to integrate photon emission, logic, memory and
detection elements into single semiconductor chips to form a photonic integrated circuit for quantum information
processing.
The author would like to thank Mark Stevenson, Robert Young, Anthony Bennett, Martin Ward and Andy Hudson for
their useful comments during the preparation of the manuscript and the UK DTI “Optical Systems for Digital Age”,
EPSRC and EC Future and Emerging Technologies programmes for supporting research on quantum light sources.
TextBox : Photon Correlation Measurements
The photon statistics of light can be studied via the second order correlation function, g(2)(τ), which describes the
correlation between the intensity of the light field with that after a delay τ and is given by100
This function can be measured directly using the Hanbury-Brown and Twiss101 interferometer, comprising a 50/50
beamsplitter and two single-photon detectors, shown in the figure. For delays much less than the average time between
detection events (ie for low intensities), the distribution in the delays between clicks in each of the two detectors is
proportional to g(2)(τ).
For a continuous light source with random emission times, such as an ideal laser or LED, g(2)(τ)=1. It shows there is no
correlation in the emission time of any two photons from the source. A source for which g(2)(τ=0)>1 is described as
'bunched' since there is an enhanced probability of two photons being emitted within a short time interval. Photons
emitted by quantum light sources are typically 'anti-bunched', (g(2)(τ=0)<1) and tend to be separated in time.
In communication and computing systems, we are interested in pulsed light sources, for which the emission occurs at
times defined by an external clock. In this case g(2)(τ) consists of a series of peaks separated by a clock period. For an
ideal single-photon source, the peak at zero time delay is absent, g(2)(τ=0)=0; as the source cannot produce more than
one photon per excitation period, clearly the two detectors cannot fire simultaneously.
The figure shows g(2)(τ) recorded for resonant pulsed optical excitation of the X emission of a single quantum dot in a
pillar microcavity. Notice the almost complete absence of the peak at zero delay: the definitive signature of a single-
photon source. The weak peak seen at τ=0 demonstrates that the rate of two-photon emission is 50 times less than that
of an ideal laser with the same average intensity. The bunching behaviour observed for the finite delay peaks is
explained by intermittent trapping of a charge carrier in the dot.102 This trace was taken for quasi-resonant laser
excitation of the dot which avoids creating carriers in the surrounding semiconductor. For higher energy laser
excitation, the suppression in g(2)(0) is typically reduced indicating occasional 2-photon pulses due to emission from the
layers surrounding the dot, but can be minimised with careful sample design.
Figure textbox: (a) Schematic of the set-up used for photon correlation measurements, (b) second order correlation
function of the exciton emission of a single dot in a pillar microcavity.
Figure Captions
Figure 1: Self assembled quantum dots (a) Image of a layer of InAs/GaAs self assembled quantum dots recorded in an
Atomic Force Microscope (AFM). Each yellow blob corresponds to a dot with typical lateral diameters of 20-30nm and
a height of 4-8nm. (b) AFM image23 of a layer of InAs quantum dots whose locations have been seeded by a matrix of
nanometer sized pits patterned onto the wafer surface. Under optimal conditions up to 60% of the etch pits contain a
single dot (Courtesy of P Atkinson & D A Ritchie, Cambridge). (c) Cross-sectional STM image of an InAs dot inside a
GaAs device (Courtesy of P. Koenraad, Eindhoven).
Figure 2: Optical spectrum of a quantum dot. (a) Schematic of the biexciton cascade of a quantum dot. (b) Typical
photoluminescence spectrum of a single quantum dot showing sharp line emission due to the biexciton X2 and exciton X
photon emitted by the cascade. The inset shows the polarisation splitting of the transitions originating from the spin
splitting of the exciton level.
Figure 3: SEM images of semiconductor cavities, including pillar microcavities (a)56 and (b), microdisk (c)51 and
photonic bandgap defect cavities (d)47, (e) and (f).48 (Structures fabricated at Univ Wuerzburg (a), CNRS-LPN (UPR-
20), Marcoussis (b, c, e), Univ Cambridge (d), UCSB/ETHZ Zurich (f))
Figure 4: Two Photon Interference. (a) If the two photons are indistinguishable, the two outcomes resulting in one
photon in either arm interfere destructively. This results in the two photons always exiting the beamsplitter together. (b)
Schematic of an experiment using two photons emitted successively from a quantum dot, (c) experimental data showing
suppression of the co-incidence rate in (b) when the delay between input photons is zero due to two-photon
interference.60 (Courtesy of Y Yamamoto, Stanford Univ.)
Figure 5: Electrically driven single-photon emission. (a) Schematic of a single-photon LED. (b) Electroluminescence
spectra of the device. Notice the spectra are dominated by the exciton X and biexciton X2 lines, which have linear and
quadratic dependence on drive current, respectively. Other weak lines are due to charged excitons. (c) second order
correlation function recorded for the exciton (i) and biexciton (ii) emission lines, (d) time-resolved electroluminescence
from a device operate with a 1.07GHz repetition rate, (e) measured (i) and modelled (ii) second order correlation
function of the biexciton electroluminescence at 1.07GHz. (adapted from Refs. 71and 73)
Figure 6: Generation of entangled photons by a quantum dot. (a) Degree of correlation measured for a dot with exciton
polarisation splitting S=0 µeV in linear (i), diagonal (ii) and circular (iii) polarisation bases as a function of the delay
between the X and X2 photons (in units of the repetition cycle). The correlation is defined as the rate of co-polarised
pairs minus the rate of cross-polarised pairs divided by the total rate. Notice that the values at finite delay show no
correlation, as expected for pairs emitted in different laser excitation cycles. More interesting are the peaks close to
zero time delay, corresponding to X and X2 photon emitted from the same cascade. The presence of strong correlations
for all three types of measurement for the dot with zero exciton splitting can only be explained if the X and X2
polarisations are entangled. (b) Degree of correlations measured for the dot in (a) subject to in-plane magnetic field so
as to produce an exciton polarisation splitting of S=25 µeV. Notice that the correlation in diagonal and circular bases
have vanished, indicating only classical correlations at finite splitting. (c) Two-photon density matrix of the device
emission in (a). The strong off-diagonal terms appear due to entanglement. (adapted from Ref 92)
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500 nm 500 nm
Fig. 1
1375 1380 1385
Detection polarisation:
Vertical
Horizontal
Photon Energy (meV)
1378.0 1378.5 1380.0 1380.5
(b)(a)
Fig. 2
ground state
500 nm
500nm
(c)(a) (b)
(e) (f)
Fig. 3
(a) (b)
Fig. 4
substrate/buffer
n-ohmic
contact
InAs QD
insulator
Al p-ohmic
contact
emission
n+ Bragg mirror
cavity layer
contact
metal
p+ GaAs
Semicon/air
interface
905 910 915
X X-X
x100 0.11µA
12.0µA
wavelength (nm)
95.1µ A
X+(b)
-40 -20 0 20 40
Time (ns)
-40 -20 0 20 40
delay (ns)
(ii) X
(c) (i)
-10 -5 5 10
delay (ns)
(i) calculated
(ii) measured
time (ns)
Fig. 5
-0.05
(c) Real Part Imaginary Part
Errors
(magnitude)
S = 0µeV
delay period (/12.5ns)
(a)(i)
(ii)
(iii)
-15 0 15
S = 25µeV
(b)(i)
(iii)
-15 0 15
Fig. 6
detector
beamsplitter
detector
device
emission
-40 -20 0 20 40
delay, τ [ns]
(a) (b)
Fig. textbox
|
0704.0404 | To the origin of the difference of FSI phases in $B\to\pi\pi$ and
$B\to\rho\rho$ decays | arXiv:0704.0404v1 [hep-ph] 3 Apr 2007
To the origin of the difference of FSI phases in
B → ππ and B → ρρ decays
A.B. Kaidalov∗ and M.I. Vysotsky†
ITEP, Moscow, Russia
Abstract
The final state interactions (FSI) model in which soft rescattering
of low mass intermediate states dominates is suggested. It explains
why the strong interaction phases are large in the Bd → ππ channel
and are considerably smaller in the Bd → ρρ one. Direct CP asym-
metries of Bd → ππ decays which are determined by FSI phases are
considered as well.
1 Introduction
There are three reasons to study FSI in B decays: to predict (or explain)
the pattern of branching ratios, to study strong interactions, and to forsee in
what decays direct CPV will be large. In view of this necessity a model for
FSI in B decays to two light mesons is suggested and explored in the present
paper.
The probabilities of three B → ππ and three B → ρρ decays are measured
now with good accuracy. The C-averaged branching ratios of these decays
are presented in Table 1 [1]. Let us look at the ratio of the charge averaged
Bd decay probabilities to the charged and neutral mesons:
Br(Bd → ρ+ρ−)
Br(Bd → ρ0ρ0)
≈ 20 , Rπ ≡
Br(Bd → π+π−)
Br(Bd → π0π0)
≈ 4 . (1)
∗[email protected]
†[email protected]
http://arxiv.org/abs/0704.0404v1
Table 1
Mode Br(10−6) Mode Br(10−6)
Bd → π+π− 5.2± 0.2 Bd → ρ+ρ− 23.1± 3.3
Bd → π0π0 1.3± 0.2 Bd → ρ0ρ0 1.16± 0.46
Bu → π+π0 5.7± 0.4 Bu → ρ+ρ0 18.2± 3.0
C-averaged branching ratios of B → ππ and B → ρρ decays.
The large difference of Rρ and Rπ is due to the difference of FSI phases
in B → ρρ and B → ππ decays (see below). In Section 2 we will determine
the differences of FSI phases of tree amplitudes which describe B → ρρ and
B → ππ decays into the states with isospins zero and two from the data
presented in Table 1. As a next step we will suggest a mechanism which
produces such phases. Once this mechanism is defined it becomes possible to
calculate FSI phases of decay amplitudes into states with a definite isospin
(not only their differences). A central question is: what intermediate states
produce FSI phases in B-meson decays into two light mesons. In the weak
decay b → uū(dd̄)d in the rest frame of a heavy quark (which is B-meson
rest frame as well) three fast light quarks are produced. Their energies are
of the order of MB/3 and momenta are more or less isotropically oriented.
The energy of the fourth (spectator) quark is of the order of ΛQCD. This four
quark state transforms mainly into multi pi-meson final state with the average
pion multiplicity about 9 (this number follows from the experimentally known
charged particles multiplicity in e+e− annihilation at Ecm = 3GeV multiplied
by 1.5 ∗ 1.5 in order to take neutral pions and third quark jet into account).
The total branching ratio of such decays is about 10−2. However such meson
state does not transform into the state composed from two light mesons
moving into opposite directions with momenta MB/2. What meson state
does transform into two light mesons can be understood from the inverse
reaction of two light meson scattering at the center of mass energy equal
to the mass of B-meson. The produced hadronic state consists of two jets
of particles moving in opposite directions. Each jet should originate from
a quark-antiquark pair produced in the weak decay of b-quark. The square
of invariant mass of a jet which contains spectator quark does not exceed
MBΛQCD and is much smaller than M
B. The energy of this jet is determined
by that of a companion quark and is about MB/2. That is why the square
of invariant mass of the second jet also does not exceed MBΛQCD. So for
B-decays the mass of a hadron cluster which transforms into light meson
in the final state should not exceed 1.5 GeV. Following these arguments in
the calculation of the imaginary parts of the decay amplitudes we will take
into account only two (relatively light) particle intermediate states for which
branching ratios of B-meson are maximal.
In Section 3 we will calculate FSI phases of tree amplitudes describing
B → ππ decays taking into account ρρ, ππ and πa1 intermediate states which
by t(u)-channel exchanges are converted into ππ. We will find that large
probability of B → ρ+ρ− decay explains about half of FSI phases of B → ππ
decays. Relatively small probability of B → π+π− decay prevents generation
of noticeable FSI phase of B → ρρ amplitudes through B → π+π− → ρρ
chain.
We will demonstrate that the strong interaction phase of the penguin
amplitude is opposite to the result of quark loop calculation, which is very
important for the value of a direct CPV asymmetry Cπ+π− ≡ C+− discussed
in Section 4. Predictions for CPV asymmetries C00 and S00 will be presented
in Section 4 as well and the value of the unitarity triangle angle α will be
extracted from the experimental data on CPV asymmetry S+−.
Subject of rare B decays is an object of intensive study nowadays and an
interested reader can find extensive list of references in a recent paper [2].
2 Phenomenology; |δπ0 − δπ2 | and |δ
0 − δ
Let us present B → ππ decay amplitudes in the so-called “t-convention”, in
which the penguin amplitude with the intermediate c-quark multiplied by
ud + VcbV
cd + VtbV
td = 0 is subtracted from the decay amplitudes [3]:
MB̄d→π+π− =
|VubV ∗ud|m2Bfπf+(0)
iδπ2 +
+ e−iγ
iδπ0 +
V ∗tdVtb
eiβPei(δ
+δ̃π0 )
, (2)
MB̄d→π0π0 =
|VubV ∗ud|m2Bfπf+(0)
iδπ2 −
− e−iγ 1√
iδπ0 −
V ∗tdVtb
eiβPei(δ
+δ̃π0 )
, (3)
MB̄u→π−π0 =
|VubV ∗ud|m2Bfπf+(0)
e−iγA2e
, (4)
where Vik are the elements of CKM matrix, γ and β are the unitarity triangle
angles and we factor out the product m2Bfπf+(0) which appears when the
decay amplitudes are calculated in the factorization approximation. A2 and
A0 are the absolute values of the decay amplitudes into the states with I = 2
and 0, generated by operators O1 and O2 (tree amplitudes), while P is the
absolute value of QCD penguin amplitude (generated by operators O3 − O6
of effective nonleptonic Hamiltonian which describes b quark decays into the
states without charm and strange quarks). δπ0 , δ
2 and δ̃
0 are FSI phases of
these three amplitudes, and it is very important for what follows that all of
them are different. It is easy to understand why δπ0 is different from δ
2 : strong
interaction depends on the isospin and is different for I = 0 and I = 2. For
example, there are definitely quark-antiquark resonances with I = 0, while
exotic resonances with I = 2 should be made from at least four quarks and
their existence is questionable. The reason why δπ0 differs from δ̃
0 is more sub-
tle. Let us consider the intermediate state made from two charged ρ-mesons
which contributes to FSI phases: Bd → ρ+ρ− → ππ. ρ+ρ− intermediate state
contribution to FSI phases can be large since Br(Bd → ρ+ρ−) is big. Both
tree and penguin induced amplitudes get FSI phases through this chain. Its
contribution to δπ0 is proportional to
(BrBd → ρ+ρ−)T/(BrBd → π+π−)T ≈
(BrBd → ρ+ρ−)/(BrBd → π+π−) ≈ 2.1, while that to δ̃π0 is proportional to
(BrBd → ρ+ρ−)P/(BrBd → π+π−)P .
How can we determine the penguin contributions to the probabilities
of Bd → ρ+ρ− and Bd → π+π−-decays? The most straightforward way
suggested in literature is to extract them from the probabilities of Bu →
K0∗ρ+ and Bu → K0π+ decays to which tree amplitudes almost do not
contribute [4, 5]1:
Br(Bd → ρ+ρ−)P =
η2 + (1− ρ)2
]2 τBd
Br(K0∗ρ+) ≈ 0.34·10−6 ,
1Contribution of tree amplitudes to these decays comes from the rescattering (Bu →
K+π0)T , K
+π0 → K0π+, and taking into account CKM suppression of the tree am-
plitudes of B → Kπ(K∗ρ) decays relative to the penguin amplitudes we can cautiously
estimate tree contribution as not more than 10% of penguin one .
Br(Bd → π+π−)P =
η2 + (1− ρ)2
]2 τBd
Br(K0π+) ≈ 0.59·10−6 ,
where fρ = 209 MeV and fK∗ = 218 MeV are the vector meson decay
constants, λ = 0.23, η = 0.34 and ρ = 0.20 are the CKM matrix parameters
in Wolfenstein parametrization [6], fK/fπ = 1.2 and the central values of
Br(Bu → K0∗ρ+) = (9.2±1.5)·10−6 and Br(Bu → K0π+) = (23.1±1.0)·10−6
[1] were used. The accuracy of equations (5) and (6) depends on the accuracy
of d ↔ s interchange symmetry (U -spin symmetry) of b → d(s) transition
amplitudes described by QCD penguin, however when the ratio of (5) to (6)
is calculated uncertainty factors partially cancel out and we obtain rather
stable result: instead of being enchanced as in the case of the tree amplitude
intermediate vector mesons contribution into penguin Bd → π+π− amplitude
is suppressed, (δ̃π0 )ρρ ≈ 1/2.8(δπ0 )ρρ. Taking into account that fraction of
longitudinally polarized vector mesons produced in Bu → K0∗ρ+ decays is
about 50% we get additional suppression of (δ̃π0 )ρρ by factor
Finally, phase δπP comes from the imaginary part of the penguin loop
with c-quark propagating in it [8]. In order to calculate δπP let us consider
corresponding quark diagram. The charm penguin contribution is given by
the following expression:
P = −Pc(k2) =
) + i
1− 4m
, (7)
where k is the sum of momenta of two quarks to which gluon radiated from
penguin decays: k = p1 + p2. One of these quarks forms π-meson with the
spectator quark, so neglecting spectator quark momentum in the rest frame
of B-meson we have p1 = (
). The second quark forms another π-meson
with d̄-quark radiated from penguin: p2 = x(
) where 0 < x < 1 is
the fraction of π+ momentum carried by u-quark. Substituting k2 = xm2b
into (7) and integrating it with the asymptotic quark distribution function
in π-meson ϕπ(x) = x(1−x) we obtain the value of δπP which depends on the
ratio 4m2c/m
b . In particular, for mb = 5.3 GeV and mc = 1.9 GeV (which
correspond to the masses of physical states) we obtain δπP ≈ 10o, a small
positive value. A nonperturbative calculation of δπP described in Section 3
demonstrates that the sign of δπP can be negative.
Our next task is to determine the difference of FSI phases δπ0 − δπ2 (the
large value of it is responsible for a relatively small value of Rπ). If we neglect
the penguin contribution, then from (2) - (4) we get the following expression:
cos(δπ0 − δπ2 ) =
B+− − 2B00 + 23
B+− +B00 − 23
, (8)
where Bik’s are the C-averaged branching ratios, while τ0/τ+ ≡ τ(Bd)/τ(Bu) =
0.92. Substituting the central values from Table 1 we get |δπ0 − δπ2 | = 48o.
Penguin contributions to Bik do not interfere with tree ones because α =
π − β − γ is almost equal to π/2. Taking P 2 terms into account with the
help of (6) (subtracting 0.59 and 0.30 from the first and the second lines of
Table 1 numbers describing B → ππ data correspondingly) we get:
|δπ0 − δπ2 | = 37o ± 10o . (9)
The accuracy of this 11o decrease of the absolute value of the phases difference
is determined by the accuracy of (6) and is not high. In recent paper [2] the
global fit of B → ππ and B → πK decay data was made. The tree amplitudes
of B → ππ decays were designated in [2] by T for B → π+π− and by C for
B → π0π0. According to [2] the difference of FSI phases between C and T
equals δC = −58o ± 10o, |C| = 0.37 ± 0.05, |T | = 0.57 ± 0.05 in the units
of 104 eV. The phase shift between the isospin amplitudes is determined by
these quantities:
tan(δ0 − δ2) =
3TC sin(−δC)
2T 2 + TC cos δC − C2
, (10)
and substituting the numbers we obtain:
δ0 − δ2 = 40o ± 7o , (11)
the result very close to (9). However, the same d ↔ s interchange symmetry
was used in [2] when relating B → ππ and B → Kπ decays. Fit [2] was
made in the same “t-convention” which we use (see the statement at the end
of page 3 of the paper [2]: “for simplicity, we will assume ... Ptc = Ptu”),
therefore the obtained results can be directly compared with ours.
Now let us consider B → ρρ decays. According to BABAR and BELLE
results ρ mesons produced in B decays are almost entirely longitudinally
polarized (fL(ρ+ρ−) = 0.98± 0.03[9], fL(ρ+ρ0) = 0.91± 0.4 [10], fL(ρ0ρ0) =
0.86 ± 0.12 [11]). For B decays into the longitudinally polarized ρ-mesons
we can write formulas analogous to (2) - (4) and we can find FSI phases
difference with the help of analog of (8). Substituting the central values of
branching ratios of B → ρρ decays from Table 1 we obtain: |δρ0−δ
2 | = 21o. In
order to subtract the penguin contribution with the help of (5) we should take
into account that in Bu → K0∗ρ+ decays the fraction of the longitudinally
polarized vector mesons equals approximately 50% [12], so we should subtract
0.17 · 10−6 in case of decay to ρ+ρ− and 0.08 · 10−6 for decay into ρ0ρ0. In
this way we obtain:
|δρ0 − δ
2 | = 20o+8
−20o , (12)
and the factor 2 difference between (12) and (9) or (11) is responsible for
the different patterns of B → ρρ and B → ππ decay probabilities. Let us
emphasize that while |δρ0 − δ
2 | being only one standard deviation from zero
can be very small this is not so for |δπ0 − δπ2 |.
3 Calculation of the FSI phases of B → ππ
and B → ρρ decay amplitudes
Among three amplitudes of B → ππ decays (2)–(4) only two are independent.
We will calculate FSI phases of B → π+π0 and B → π+π− amplitudes and
extract from them FSI phases of amplitudes with a definite isospin.
Our task is to take into account the intermediate state contributions
into FSI phases. As it was argued in Introduction we should consider only
two particle intermediate states with positive G-parity to which B-mesons
have relatively large decay probabilities. Alongside with ππ and ρρ there is
only one such state: πa1. So we will consider ρρ intermediate state which
transforms into ππ by π exchange in t-channel, πa1 intermediate state which
transforms into ππ by ρ exchange in t-channel and will take into account the
elastic channel B → ππ → ππ as well. This approach is analogous to the
FSI consideration performed in paper [13]. However in [13] 2 → 2 scattering
amplitudes were considered to be due to elementary particle exchanges in
t-channel. For vector particles exchanges s-channel partial wave amplitudes
behave as sJ−1 ∼ s0 and thus do not decrease with energy (decaying meson
mass). However it is well known that the correct behavior is given by Regge
theory: sαi(0)−1. For ρ-exchange αρ(0) ≈ 1/2 and the amplitude decrease
with energy as 1/
s. This effect is very spectacular for B → DD → ππ
chain with D∗(D∗2) exchange in t-channel: αD∗(0) ≈ −1 and reggeized D∗
meson exchange is damped as s−2 ≈ 10−3 in comparison with elementary
D∗ exchange (see for example [14]). For π-exchange, which gives a dominant
contribution to ρρ → ππ transition (see below), in the small t region the pion
is close to mass shell and its reggeization is not important.
We will use Feynman diagram approach to calculate FSI phases from
the triangle diagram with the low mass intermediate states X and Y (see
Figure 1). Integrating over loop momenta d4k we assume that integrals over
masses of intermediate states X and Y decrease rapidly with increase of these
masses. Then choosing z axis in the direction of momenta of the produced
meson M1 we can transform the integral over k0 and kz into the integral over
the invariant masses of clusters of intermediate particles X and Y
dk0dkz =
dsXdsY (13)
and deform integration contours in such a way that only low mass interme-
diate states contributions are taken into account while the contribution of
heavy states being small is neglected. In this way we get:
M Iππ = M
XY (δπXδπY + iT
XY→ππ) , (14)
where M
XY are the decay matrix elements without FSI interactions and
T J=0XY→ππ is the J = 0 partial wave amplitude of the process XY → ππ (T J =
(SJ − 1)/(2i)) which originates from the integral over d2k⊥.
For real T (14) coincides with the application of the unitarity condition
for the calculation of the imaginary part of M while for the imaginary T the
corrections to the real part of M are generated.
Let us calculate the imaginary parts of B → ππ decay amplitudes which
originate from B → ρρ → ππ chain with the help of unitarity condition 2:
ImM(B → ππ) =
d cos θ
M(ρρ → ππ)M∗(B → ρρ) , (15)
where θ is the angle between ρ and π momenta. For small values of θ or t
π-exchange in t-channel dominates and the calculation of Feynman diagram
2in this section the phases which originate from CKM matrix elements are omitted.
Figure 1: Diagram which describes FSI in the decay of heavy meson MQq
into two light mesons M1 and M2. X and Y are the clusters of particles with
small invariant masses sX , sY ≤ MQΛQCD, k is 4-momentum of a virtual
particle propagating in t-channel.
for ρρ → ππ amplitude with the elementary virtual π-meson exchange can
be trusted, as it was noted above. It was already stressed that ρ-mesons pro-
duced in B-decays are almost entirely longitudinally polarized. That is why
we will take into account only longitudinal polarization for the intermediate
ρ-mesons and amplitudes of B-decays into ππ and ρLρL are simply related
MB+→ρ+ρ0 = −
MB+→π+π0 , MBd→ρ+ρ− = −
MBd→π+π− . (16)
For the amplitude of ρ+ρ0 → π0π+ transition we have:
iM(ρ+ρ0 → π0π+) = −i
g2ρππ
(p1 − k1)2 −m2π
+)(k2ρ
0) , (17)
where p1, k1 and k2 are ρ
+, π0 and π+ momenta. From the width of ρ-meson
we get g2ρππ/16π = 2.85. For the longitudinally polarized ρ-mesons in their
center of mass system we have:
+ = k2ρ
0 = − 1
(t−m2π)(1 +
) +m2ρ
, (18)
where t = (p1 − k1)2. Changing the integration variable in (15) to t with the
help of dt =
(1− 2 m
)d cos θ and introducing formfactor exp(t/µ2) with
3relative negative sign of the amplitudes follows from the expressions for transition
formfactors in the factorization approximation, see for example [15].
the parameter µ2 ∼ 1 GeV 2 we obtain:
ImMB→π+π0 = +
(m2ρ−m
g2ρππdt
16πM2B ∗ 4m2ρ
(t−m2π)(1 +
+ 2m2ρ(1 +
t−m2π
exp(t/µ2)
MB→π+π0 . (19)
For µ2 = 2m2ρ the contributions of the first two terms in square brackets
cancel, while the third term gives:
ImMB→π+π0 = −
g2ρππ
3.1MB→π+π0 , (20)
and from (4) we get:
δπ2 (ρρ) = −4.9o . (21)
Let us note that in the limit MB → ∞ the ratio Br(Bd → ρρ)/Br(Bd → ππ)
grows as M2B, that is why FSI phase δ
2 (ρρ) (and δ
0 (ρρ)) diminishes as 1/MB.
The analogous consideration of ρ+ρ− intermediate state leads to the pos-
itive FSI phase of Bd → π+π− amplitude which is enhanced relatively to
δπ2 (ρρ) according to (16):
δπ+−(ρρ) = +5.7
o , (22)
and for FSI phase of the amplitude with isospin zero in the linear approxi-
mation we get:
δπ0 (ρρ) = δ
+−(ρρ) +
δπ+−(ρρ)− δπ2 (ρρ)
. (23)
We are able to extract the ratio A2/A0 from that of C-averaged Br(Bd →
π+π−), Br(Bd → π0π0) and Br(Bu → π+π0), subtracting penguin contribu-
tion as we did deriving (9):
B̃+− + B̃00
− 1 , (24)
= 0.80± 0.09 , (25)
and, finally:
δπ0 (ρρ) = 15
o , δπ0 (ρρ)− δπ2 (ρρ) = 20o . (26)
In this way we see that B → ρρ → ππ chain generates half of the experi-
mentally observed FSI phase difference of B → ππ tree amplitudes.
It is remarkable that FSI phases generated by B → ππ → ρρ chain are
damped by Br(B → ρ+ρ−, ρ+ρ0)/Br(B → π+π, π+π0) ratios and are a few
degrees:
2(ππ) =
∗ δπ2 (ρρ) = −1.4o , δ
+−(ππ) =
∗ δπ+−(ρρ) = 1.2o ,
(A0/A2)ρρ = 1.1 , δ
0(ππ) = 2.9
o , δ
0(ππ)− δ
2(ππ) ≈ 4o . (27)
Next we will take into account ππ intermediate state. From Regge analy-
sis of ππ elastic scattering we know that good description of the experimental
data is achieved when the exchanges of pomeron, ρ and f trajectories in t-
channel are taken into account [16]. Pomeron exchange dominates in elastic
ππ → ππ scattering at high energies. For αP (0) = 1 the corresponding
amplitude T is purely imaginary and the phases of matrix elements do not
change [3]. However taking into account that pomeron is ”supercritical”,
αP (0) ≈ 1.1, we obtain the phase of the amplitude generated by pomeron
exchange 4 which cancels the phases generated by ρ and f exchanges for
I = 2. For I = 0 the sum of ρ and f exchanges produces the purely imag-
inary amplitude T and the phase of the amplitude M is due to pomeron
”supercriticallity”:
δπ0 (ππ) = 5.0
o , δπ2 (ππ) = 0
o . (28)
In paper [3] the pomeron exchange amplitude was considered as purely
imaginary. As a result though important for branching ratios phase difference
δπ0 (ππ)−δπ2 (ππ) was the same (pomeron contribution being universal cancels
in the difference of phases) it came mainly from δπ2 (ππ) negative value. In
4The amplitude of 2 → 2 process due to supercritical pomeron exchange is T ∼
(s/s0)
αP (t)(1 + exp(−iπαP (t)))/(− sin(παP (t))) = (s/s0)(1+∆)(i + ∆π/2), where in the
last expression t = 0 was substituted and αP (0) = 1 +∆ was used (∆ ≈ 0.1).
this way result for the absolute value of direct CP-asymmetry Cπ+π− was
underestimated, see below.
Finally πa1 intermediate state should be accounted for. Large branching
ratio of Bd → π±a∓1 -decay ( Br(Bd → π±a∓1 ) = (40 ± 4) ∗ 10−6) is partially
compensated by small ρπa1 coupling constant (it is 1/3 of ρππ one). As a
result the contributions of πa1 intermediate state (which transforms into ππ
by ρ-trajectory exchange in t-channel) to FSI phases equal approximately
that part of ππ intermediate state contributions which is due to ρ-trajectory
exchange. Assuming that the sign of the πa1 intermediate state contribution
into phases is the same as that of elastic channel we obtain:
δπ0 (πa1) = 4
o , δπ2 (πa1) = −2o . (29)
Summing the imaginary parts of the amplitudes which follow from (21),
(26), (28) and (29) we finally obtain:
δπ0 = 23
o , δπ2 = −7o , δπ0 − δπ2 = 30o , (30)
and the accuracy of these numbers is not high, at the level of 50%.
The analogous consideration of the real parts of the loop corrections to
B → ππ decay amplitudes leads to the diminishing of the (real) tree am-
plitudes by ≈ 30%, and we can explain the experimentally observed value
δπ0 − δπ2 ≈ 40o in our model while for ρρ final state the analogous difference
is about three times smaller, δ
0 − δ
2 ≈ 15o.
Let us estimate the phase of the penguin amplitude δπP considering charmed
mesons intermediate states: B → D̄D, D̄∗D, D̄D∗, D̄∗D∗ → ππ 5. In Regge
model all these amplitudes are described at high energies by exchanges of
D∗(D∗2)-trajectories. An intercept of these exchange-degenerate trajectories
can be obtained using the method of [17] or from masses of D∗(2007) – 1−
andD∗2(2460) – 2
+ resonances, assuming linearity of these Regge-trajectories.
Both methodes give αD∗(0) = −0.8÷−1 and the slope α′D∗ ≈ 0.5GeV −2.
The amplitude of D+D− → π+π− reaction in the Regge model proposed
in papers [18, 19] can be written in the following form:
TDD̄→ππ(s, t) = −
e−iπα(t)Γ(1− αD∗(t))(s/scd)αD∗(t) , (31)
5These amplitudes are considered as penguin due to the proper combination of CKM
matrix elements.
where Γ(x) is the gamma function.
The t-dependence of Regge-residues is chosen in accord with the dual
models and is tested for light (u,d,s) quarks [18]. According to [19] scd ≈
2.2GeV 2.
Note that the sign of the amplitude is fixed by the unitarity in the t-
channel (close to the D∗-resonance). The constant g20 is determined by the
width of the D∗ → Dπ decay: g20/(16π) = 6.6. Using (14), analog of (15),
(31) and the branching ratio Br(B → DD̄) ≈ 2 · 10−4 [20] we obtain the
imaginary part of P and comparing it with the contribution of P in B →
π+π− decay probability (6) we get δπP ≈ −3.5o6. A smallness of the phase
is due to the low intercept of D∗-trajectory. The sign of δP is negative -
opposite to the positive sign which was obtained in perturbation theory (7).
Since DD̄-decay channel constitutes only ≈ 10% of all two-body charm-
anticharm decays of Bd-meson [20] taking these channels into account we can
easily get
δP ∼ −10o , (32)
which may be very important for the interpretation of the experimental data
on direct CP asymmetry C+− discussed in the next section.
4 CP asymmetries of Bd(B̄d) → ππ decays
The CP asymmetries are given by :
Cππ ≡
1− |λππ|2
1 + |λππ|2
, Sππ ≡
2Im(λππ)
1 + |λππ|2
, λππ ≡ e−2iβ
MB̄→ππ
MB→ππ
, (33)
where ππ is π+π− or π0π0.
From (2) for direct CP asymmetry in Bd(B̄d) → π+π− decays we readily
obtain:
C+− = −
sinα[
2A0 sin(δ0 − δ̃0 − δP ) + A2 sin(δ2 − δ̃0 − δP )]/
cos(δ0 − δ2)−
A0P̃ cosα cos(δ0 − δ̃0 − δP )−
6In integration over cos θ the region θ ≪ 1 dominates. In this region representation
(31) is valid.
− A2P̃√
cosα cos(δ2 − δ̃0 − δP ) + P̃ 2] , (34)
where
V ∗tdVtb
P . (35)
In order to make a numerical estimate we should know the ratios A0/A2
and P/A2. The first one is given by (25) while the second one can be extracted
from the ratio Br(Bu → K0π+)/Br(Bu → π0π+) assuming d ↔ s invariance
of the strong interactions:
Br(Bu → K0π+)
Br(Bu → π0π+)
f 2KP
2|V ∗tsVtb|2
A22|V ∗udVub|2
, (36)
= 0.092(0.009) . (37)
The numerical values of A0 andA2 are given with good accuracy by factor-
ization calculation, while P appears to be 2.5 times larger than factorization
result [3]. In view of this the validity of factor fK in (36) which originates
from factorization calculation of the penguin amplitude is questinable. If
factorization of the penguin amplitudes is not assumed then the ratio fK/fπ
in (36) should be replaced by unity. In this way we get 20% larger value of
P/A2 in (37) and we will take this value of uncertainty as an estimate of the
theoretical accuracy of the determination of P :
= 0.21(0.04) , (38)
Taking into account that unitarity triangle angle α ≈ 90o and angles δ̃0 and
δP are of the order of few degrees from (34) we obtain:
C+− ≈ −0.28[1.1 sin(δ0 − δ̃0 − δP ) + sin(δ2 − δ̃0 − δP )] ≈
≈ −0.56 sin((δ0 + δ2)/2− δ̃0 − δP ) . (39)
In order to determine the lower bound on the value of C+− let us suppose
that δ0 = 37
o, δ2 = 0
o (we keep the difference δ0− δ2 = 37o, as it follows from
the data on B → ππ decay probabilities (9)), and neglect small values of δ̃0
and δP :
C+− > −0.18 . (40)
Concerning experimental number it could well happen that finally it will
be considerably below our bound. In this case the result of nonperturbative
calculation of penguin phase will be confirmed. Substituting in (39) δ0 =
30o, δ2 = −7o and δP from (32) we obtain the following central value:
C+− = −0.21 . (41)
It is instructive to compare the obtained numbers with the value of C+−
which follows from the asymmetry ACP (K
+π−) if d ↔ s symmetry is sup-
posed [21]:
C+− =
ACP (K
Γ(B → K+π−)
Γ(B → π+π−)
sin(β + γ)
sin(γ)
= 1.2(−2)(−0.093± 0.015)19.8
sin 82o
sin 60o
0.87 = −0.24± 0.04 . (42)
Let us note that one factor fπ/fK in the last equation appears from the
matrix element of the tree operator, the second one - from the matrix element
of the penguin operator. If because of nonfactorization of penguin amplitudes
we will omit the factor which appears from the penguin [5], then the numbers
in the right-hand sides of (40, 41) and (42) will become 20% smaller.
The experimental results obtained by Belle [22] and BABAR [23] are
contradictory
CBelle+− = −0.55(0.09) , CBABAR+− = −0.21(0.09), (43)
Belle number being far below (40) and (41).
For direct CP asymmetry in Bd(B̄d) → π0π0 decay from (3) we readily
obtain:
C00 = −
P̃ sinα[A0 sin(δ0 − δ̃0 − δP )−
2A2 sin(δ2 − δ̃0 − δP )]/
A0A2 cos(δ0 − δ2)−
A0P̃ cosα cos(δ0 − δ̃0 − δP ) +
A2P̃ cosα cos(δ2 − δ̃0 − δP ) + P̃ 2] , (44)
C00 ≈ −1.06[0.8 sin(δ0 − δ̃0 − δP )− 1.4 sin(δ2 − δ̃0 − δP )] ≈ −0.6 , (45)
considerably smaller than C+−. This unusually large direct CPV (measured
by |C00|) is intriguing task for future measurements since the present exper-
imental error is too big:
exper
00 = −0.36(0.32) . (46)
Belle and BABAR agree now on the value of another CPV asymmetry
measured in Bd(B̄d) → π+π− decays: Sexper+− = −0.62 ± 0.09 [22, 23]. From
this measurement the value of unitarity triangle angle α can be extracted.
Neglecting the penguin contribution we get:
sin 2αT = S+− , (47)
αT = 109o ± 3o . (48)
Penguin shifts the value of α. The accurate formula looks like:
S+− = [sin 2α(
cos(δ0 − δ2))−
A2P̃√
sinα cos(δ2 − δ̃0 − δP )−
A0P̃ sinα cos(δ0 − δ̃0 − δP )]/
cos(δ0 − δ2)−
A0P̃ cosα cos(δ0 − δ̃0 − δP )−
A2P̃√
cosα cos(δ2 − δ̃0 − δP ) + P̃ 2] , (49)
and since all the phase shifts are not big the values of cosines in (49) are
rather stable relative to their variations. For numerical estimates we take
δ0 = 30
o, δ2 = −7o and neglect δ̃0 and δP . In this way we get:
(α)ππ = 88
o ± 40(exper)± 50(theor) , (50)
where the first error comes from uncertainty in S
exper
+− while the second one
comes from that in the value of penguin amplitude, (38). Relatively large
theoretical uncertainty in the value of P̃ does not prevent to determine α
with good precision.
The relative smallness of penguin contribution to B → ρρ decay am-
plitudes allow us to determine α with better theoretical accuracy from the
experimental measurement of (S+−)ρρ just as it was done in [24]. With the
help of (5) we obtain:
)ρρ = 0.060(0.012) , (51)
where the same 20% uncertainty in extracting penguin amplitude is supposed.
Using the ratio (A0/A2)ρρ determined in (27) from the (49) neglecting strong
phases (which are much smaller than in the case of B → ππ decays) and
taking into account the recent experimental result (S
exper
+− )ρρ = −0.06± 0.18
[1] we obtain:
(α)ρρ = 87
o ± 50(exper)± 10(theor) . (52)
Let us point out that considerably larger theoretical error quoted in [4]
follows from the larger theoretical uncertainty in the value of penguin ampli-
tude assumed in that paper.
Our results for α should be compared with the numbers which follow from
the global fit of unitarity triangle [6, 7]:
αCKMfitter = (99.0+4.0−9.4)
o , αUTfit = (93± 4)o . (53)
We conclude this section with the prediction for the value of CPV asym-
metry S00:
S00 = [sin 2α(
2A0A2
cos(δ0 − δ2)) +
2A2P̃√
sinα cos(δ2 − δ̃0 − δP )−
A0P̃ sinα cos(δ0 − δ̃0 − δP )]/
2A0A2
cos(δ0 − δ2)−
A0P̃ cosα cos(δ0 − δ̃0 − δP ) +
2A2P̃√
cosα cos(δ2 − δ̃0 − δP ) + P̃ 2] = 0.70± 0.15 , (54)
a large asymmetry with the sign opposite to that of S+−.
5 Conclusions
FSI appeared to be very important in B → ππ decays.
The description of these interactions presented in the paper allows to
explain the experimentally observed difference of the ratios of decay proba-
bilities to the neutral and charged modes in B → ππ and B → ρρ decays.
Rather large absolute value of direct CP asymmetry C+− (if confirmed
experimentally) will be a manifestation of the negative sign of penguin FSI
phase in accord with nonperturbative calculation and opposite to perturba-
tive result.
We are grateful to L.V.Akopyan for checking formulas, Jose Ocariz for
recommendation to include the result for angle α which follows from CP
asymmetry (S+−)ρρ and M.B.Voloshin for useful discussion.
This work was supported by Russian Agency of Atomic Energy;
A.K. was partly supported by grants RFBR 06-02-17012, RFBR 06-02-
72041-MNTI, INTAS 05-103-7515 and state contract 02.445.11.7424;
M.V. was partly supported by grants RFBR 05-02-17203 and
NSh-5603.2006.2.
References
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(2006) 68.
[5] M. Gronau, J.L. Rosner, Phys. Lett. B595 (2004) 339.
[6] CKM fitter, http://ckmfitter.in2p3.fr.
[7] UTfit, http://utfit.roma1.infn.it.
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G.M. Gérard and W.-S. Hou, Phys. Rev. D43, (1991) 2909.
[9] B. Aubert et al., BABAR Collaboration, hep-ex/0607098 (2006).
[10] B. Aubert et al., BABAR Collaboration, Phys. Rev. Lett. 97 (2006)
261801.
[11] B. Aubert et al., BABAR Collaboration, hep-ex/0607097 (2006).
[12] J. Zhong et al. Belle Collaboration, Phys. Rev. Lett. 95 (2005) 141801;
B. Aubert et al., BABAR Collaboration, Phys. Rev. Lett. 97 (2006)
201801.
[13] H-Y. Cheng, C-K. Chua and A.Soni, Phys. Rev. D71 (2005) 014030.
[14] A.Deandrea et al., Int. J. Mod. Phys. (2006) 4425.
[15] R.Aleksan et al., Phys. Lett. B356 (1995) 95.
[16] K.G.Boreskov, A.A.Grigoryan, A.B.Kaidalov, I.I.Levintov, Yad. Fiz.
27, (1978) 813.
[17] A.B.Kaidalov, Zeit. fur Phys. C12, (1982) 63.
[18] P.E.Volkovitsky, A.B.Kaidalov, Sov.J.Nucl.Phys. 35, (1982) 909.
[19] K.G.Boreskov, A.B.Kaidalov, Sov.J.Nucl.Phys. 37, (1983) 100.
[20] Review of Particle Physics, W.-M. Yao et al., Journal of Physics G 33,
(2006) 1.
[21] R.Fleischer, Phys. Lett. B459, (1999) 306.
[22] H.Ishino, Belle, talk at ICHEP06, Moscow (2006).
[23] B.Aubert et al, BABAR Collaboration, hep-ex/0703016 (2007).
[24] M.I.Vysotsky, Yad. Fiz. 69, (2006) 703.
|
0704.0405 | An invariance principle for semimartingale reflecting Brownian motions
in domains with piecewise smooth boundaries | An invariance principle for semimartingale reflecting Brownian motions in domains with piecewise smooth boundaries
The Annals of Applied Probability
2007, Vol. 17, No. 2, 741–779
DOI: 10.1214/105051606000000899
c© Institute of Mathematical Statistics, 2007
AN INVARIANCE PRINCIPLE FOR SEMIMARTINGALE
REFLECTING BROWNIAN MOTIONS IN DOMAINS WITH
PIECEWISE SMOOTH BOUNDARIES1
By W. Kang and R. J. Williams
Carnegie Mellon University and University of California, San Diego
Semimartingale reflecting Brownian motions (SRBMs) living in
the closures of domains with piecewise smooth boundaries are of in-
terest in applied probability because of their role as heavy traffic ap-
proximations for some stochastic networks. In this paper, assuming
certain conditions on the domains and directions of reflection, a per-
turbation result, or invariance principle, for SRBMs is proved. This
provides sufficient conditions for a process that satisfies the definition
of an SRBM, except for small random perturbations in the defining
conditions, to be close in distribution to an SRBM. A crucial ingredi-
ent in the proof of this result is an oscillation inequality for solutions
of a perturbed Skorokhod problem. We use the invariance principle
to show weak existence of SRBMs under mild conditions. We also
use the invariance principle, in conjunction with known uniqueness
results for SRBMs, to give some sufficient conditions for validating
approximations involving (i) SRBMs in convex polyhedrons with a
constant reflection vector field on each face of the polyhedron, and
(ii) SRBMs in bounded domains with piecewise smooth boundaries
and possibly nonconstant reflection vector fields on the boundary
surfaces.
1. Introduction. Semimartingale reflecting Brownian motions (SRBMs)
living in the closures of domains with piecewise smooth boundaries are of
interest in applied probability because of their role as heavy traffic diffusion
approximations for some stochastic networks. The nonsmoothness of the
boundary for such a domain, combined with discontinuities in the oblique
directions of reflection at intersections of smooth boundary surfaces, present
Received May 2006; revised November 2006.
1Supported in part by NSF Grants DMS-03-05272 and DMS-06-04537.
AMS 2000 subject classifications. 60F17, 60J60, 60K25, 90B15, 93E03.
Key words and phrases. Semimartingale reflecting Brownian motion, piecewise smooth
domain, invariance principle, oscillation inequality, Skorokhod problem, stochastic net-
works.
This is an electronic reprint of the original article published by the
Institute of Mathematical Statistics in The Annals of Applied Probability,
2007, Vol. 17, No. 2, 741–779. This reprint differs from the original in pagination
and typographic detail.
http://arxiv.org/abs/0704.0405v1
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http://dx.doi.org/10.1214/105051606000000899
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http://www.imstat.org
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http://dx.doi.org/10.1214/105051606000000899
2 W. KANG AND R. J. WILLIAMS
challenges in the development of a rigorous theory of existence, uniqueness
and approximation for such SRBMs.
When the state space is an orthant and the direction of reflection is con-
stant on each boundary face, a necessary and sufficient condition for weak
existence and uniqueness of an SRBM is known [14]. This condition involves
a so-called completely-S condition on the matrix formed by the reflection
directions. An invariance principle for such SRBMs was established in [15]
and used in [16] to justify heavy traffic diffusion approximations for cer-
tain open multiclass queueing networks. Loosely speaking, the invariance
principle shows that, assuming uniqueness in law for the SRBM, a process
satisfying the definition of the SRBM, except for small perturbations in the
defining conditions, is close in distribution to the SRBM.
For more general domains with piecewise smooth boundaries, some con-
ditions for existence and uniqueness of SRBMs are known. In particular, for
convex polyhedrons with a constant direction of reflection on each boundary
face, necessary and sufficient conditions for weak existence and uniqueness of
SRBMs are known for simple convex polyhedrons (where precisely d faces
meet at each vertex in d-dimensions) and sufficient conditions are known
for nonsimple convex polyhedrons, see [4]. For a bounded domain that can
be represented as a finite intersection of domains, each of which has a C1-
boundary and an associated uniformly Lipschitz continuous reflection vector
field, sufficient conditions for strong existence and uniqueness were provided
by Dupuis and Ishii [6]; in fact, these authors study stochastic differential
equations with reflection which include SRBMs. Despite these existence and
uniqueness results, a general invariance principle for SRBMs living in the
closures of domains with piecewise smooth boundaries has not been proved
to date. (We note that for the special case when the directions of reflection
are normal, that is, perpendicular to the boundary, there are a number of
perturbation results for reflecting Brownian motions. Our emphasis here is
on treating a wide range of oblique reflection directions.)
Motivated by its potential for use in approximating heavily loaded stochas-
tic networks that are more general than open multiclass queueing networks,
in this paper, we formulate and prove an invariance principle for SRBMs
living in the closures of domains with piecewise smooth boundaries with
possibly nonconstant directions of reflection on each of the smooth bound-
ary surfaces. An application of the results of this paper to the analysis of
an internet congestion control model can be found in [13]. An outline of the
current paper is as follows.
The definition of an SRBM and assumptions on the domains and direc-
tions of reflection are given in Sections 2 and 3, respectively. Some sufficient
conditions for these assumptions to hold are provided in Section 3. Section
4 is devoted to proving the main result of this paper, namely, the invari-
ance principle. A key element for our proof of this result is an oscillation
INVARIANCE PRINCIPLE FOR SRBMS 3
inequality for solutions of a perturbed Skorokhod problem; this inequality,
which may be of independent interest, is proved in Section 4.1. In Section 5
we give some applications of the invariance principle. We prove weak ex-
istence of SRBMs under the conditions specified in Section 3. We also use
the invariance principle, in conjunction with known uniqueness results for
SRBMs, to give sufficient conditions for validating approximations involving
(i) SRBMs in convex polyhedrons with a constant reflection vector field on
each face of the polyhedron, and (ii) SRBMs in bounded domains with piece-
wise smooth boundaries and possibly nonconstant reflection vector fields on
the boundary surfaces.
Beyond its possible use in justifying SRBM approximations for stochastic
networks, the invariance principle might be used to justify numerical ap-
proximations to SRBMs. A further possible extension of the results stated
here would involve an invariance principle for stochastic differential equa-
tions with reflection. The oscillation inequality for the perturbed Skorokhod
problem and associated criteria for C-tightness described in Sections 4.1
and 4.2 are likely to be useful for this. We have not developed such an ex-
tension here as that would involve introduction of extra assumptions that
would make the result less relevant for potential applications to stochastic
networks. In particular, the approximating processes would involve stochas-
tic integrals driven by a Brownian motion, whereas in stochastic network
applications, the Brownian motion typically only appears in the limit.
1.1. Notation, terminology and preliminaries. Let N denote the set of all
positive integers, that is, N = {1,2, . . .}, R denote the set of real numbers,
which is also denoted by (−∞,∞), R+ denote the nonnegative half-line,
which is also denoted by [0,∞). For x ∈ R, we write |x| for the absolute
value of x, [x] for the largest integer less than or equal to x, x+ for the
positive part of x. For any positive integer d, we let Rd denote d-dimensional
Euclidean space, where any element in Rd is denoted by a column vector.
Let ‖ · ‖ denote the Euclidean norm on Rd, that is, ‖x‖ = (
i=1 x
1/2 for
x ∈Rd, and 〈·, ·〉 denote the inner product on Rd, that is, 〈x, y〉=
i=1 xiyi,
for x, y ∈ Rd. We note that for any x ∈ Rd, ‖x‖ ≤
i=1 |xi|. Let R
+ denote
the positive orthant in Rd, that is, Rd+ = {x ∈ R
d :xi ≥ 0,1 ≤ i ≤ d}. Let
B(S) denote the Borel σ-algebra on S ⊂ Rd, that is, the collection formed
by intersecting all Borel sets in Rd with S. Let dist(x,S) denote the distance
between x ∈Rd and S ⊂Rd, that is, dist(x,S) = inf{‖x−y‖ :y ∈ S}, with the
convention that dist(x,∅) =∞ for x ∈Rd. Let Ur(S) denote the closed set
{x ∈Rd : dist(x,S)≤ r} for any r > 0 and S ⊂Rd, where if S =∅, Ur(S) =∅
for all r > 0. Let Br(x) denote the closed ball {y ∈R
d :‖y− x‖ ≤ r} for any
x ∈Rd and r > 0. For any set S ⊂Rd, we write S for the closure of S, So for
the interior of S and ∂S = S \So. For a finite set S, |S| denotes the number
4 W. KANG AND R. J. WILLIAMS
of elements in S. For any v ∈Rd, v′ denotes the transpose of v. Inequalities
between vectors in Rd should be interpreted componentwise, that is, if u, v ∈
d, then u ≤ (<)v means that ui ≤ (<)vi for each i ∈ {1, . . . , d}. For any
matrix A, let A′ denote the transpose of A. For any function x :R+ → R
x(t−) denotes the left limit of x at t > 0 whenever x has a left limit at t;
unless explicitly stated otherwise, x(0−)≡ 0, where 0 is the zero vector in Rd.
For any function x :R+ → R
d, we let ∆x(t) = x(t)− x(t−) for t ∈ R+ when
x(t−) exists. We let 0 be the constant deterministic function x :R+ → R
such that x(t) = 0 for all t ∈R+.
A domain in Rd is an open connected subset of Rd. For each continuously
differentiable function f defined on some nonempty domain S ⊂Rd, ∇f(x)
is the gradient of f at x ∈ S. For each x ∈ Rd, a neighborhood Vx of x is a
bounded domain in Rd that contains x. For any nonempty domain S ⊂Rd,
we say that the boundary ∂S of S is C1, if for each x ∈ ∂S there exists a
Euclidean coordinate system Cx for R
d centered at x, an rx > 0, and a once
continuously differentiable function ϕx :R
d−1 →R such that ϕx(0) = 0 and
S ∩Brx(x) = {z = (z1, . . . , zd)
′ in Cx : zd >ϕx(z1, . . . , zd−1)} ∩Brx(x).
Then, for x ∈ ∂S, the inward unit normal to ∂S at z ∈ ∂S ∩Brx(x) is given
in the coordinate system Cx by
n(z) =
(1 + ‖∇ϕx(z1, . . . , zd−1)‖
2)1/2
(−∇ϕx(z1, . . . , zd−1)
where ∇ϕx(z1, . . . , zd−1) = (
, . . . ,
∂zd−1
)′(z1, . . . , zd−1). For any nonempty
convex set S ⊂Rd, we call a vector n ∈Rd \{0} an inward unit normal vector
to S at x ∈ ∂S if ‖n‖= 1 and 〈n, y − x〉 ≥ 0 for all y ∈ S. Note that such a
vector need not be unique.
All stochastic processes used in this paper will be assumed to have paths
that are right continuous with finite left limits (abbreviated henceforth as
r.c.l.l.). A process is called continuous if almost surely its sample paths
are continuous. We denote by D([0,∞),Rd) the space of r.c.l.l. functions
from [0,∞) into Rd and we endow this space with the usual Skorokhod
J1-topology (cf. Chapter 3 of [7]). We denote by C([0,∞),R
d) the space
of continuous functions from [0,∞) into Rd. The Borel σ-algebra on either
D([0,∞),Rd) or C([0,∞),Rd) will be denoted by Md. The abbreviation
u.o.c. will stand for uniformly on compacts and will be used to indicate
that a sequence of functions in D([0,∞),Rd) (or C([0,∞),Rd)) is converg-
ing uniformly on compact time intervals to a limit in D([0,∞),Rd) (or
C([0,∞),Rd)). Consider W 1,W 2, . . . ,W , each of which is a d-dimensional
process (possibly defined on different probability spaces). The sequence
{W n}∞n=1 is said to be tight if the probability measures induced by the
W n on the measurable space (D([0,∞),Rd),Md) form a tight sequence,
INVARIANCE PRINCIPLE FOR SRBMS 5
that is, they form a weakly relatively compact sequence in the space of
probability measures on (D([0,∞),Rd),Md). The notation “W n ⇒W” will
mean that, as n → ∞, the sequence of probability measures induced on
(D([0,∞),Rd),Md) by {W n} converges weakly to the probability measure
induced on the same space by W . We shall describe this in words by saying
that W n converges weakly (or in distribution) to W as n → ∞. The se-
quence of processes {W n}∞n=1 is called C-tight if it is tight, and if each weak
limit point, obtained as a weak limit along a subsequence, almost surely has
sample paths in C([0,∞),Rd). The following proposition provides a useful
criterion for checking C-tightness.
Proposition 1.1. Suppose that, for each n ∈N, W n is a d-dimensional
process defined on the probability space (Ωn,Fn, Pn). The sequence {W n}∞n=1
is C-tight if and only if the following two conditions hold:
(i) For each η > 0 and T ≥ 0, there exists a finite constant Mη,T > 0
such that
lim inf
0≤t≤T
‖W n(t)‖ ≤Mη,T
≥ 1− η.(1)
(ii) For each ε > 0, η > 0 and T > 0, there exists λ ∈ (0, T ) such that
lim sup
Pn{wT (W
n, λ)≥ ε} ≤ η,(2)
where for x ∈D([0,∞),Rd),
wT (x,λ) = sup
u,v∈[t,t+λ]
‖x(u)− x(v)‖ : 0≤ t < t+ λ≤ T
Proof. See Proposition VI.3.26 in [12]. �
A d-dimensional process W is said to be locally of bounded variation if
all sample paths of W are of bounded variation on each finite time interval.
For such a process W , we define V(W ) = {V(W )(t), t≥ 0} such that for each
t≥ 0,
V(W )(t) = ‖W (0)‖
+ sup
‖W (ti)−W (ti−1)‖ : 0 = t0 < t1 < · · ·< tl = t, l≥ 1
A triple (Ω,F ,{Ft, t ≥ 0}) will be called a filtered space if Ω is a set,
F is a σ-algebra of subsets of Ω, and {Ft, t≥ 0} is an increasing family of
sub-σ-algebras of F , that is, a filtration. Henceforth, the filtration {Ft, t≥ 0}
will be simply written as {Ft}. If P is a probability measure on (Ω,F), then
6 W. KANG AND R. J. WILLIAMS
(Ω,F ,{Ft}, P ) is called a filtered probability space. A d-dimensional process
X = {X(t), t ≥ 0} defined on (Ω,F , P ) is called {Ft}-adapted if for each
t ≥ 0, X(t) :Ω→ Rd is measurable when Ω is endowed with the σ-algebra
Given a filtered probability space (Ω,F ,{Ft}, P ), a vector µ ∈R
d, a d× d
symmetric, strictly positive definite matrix Γ, and a probability distribution
ν on (Rd,B(Rd)), an {Ft}-Brownian motion with drift vector µ, covariance
matrix Γ, and initial distribution ν, is a d-dimensional {Ft}-adapted process
defined on (Ω,F ,{Ft}, P ) such that the following hold under P :
(a) X is a d-dimensional Brownian motion whose sample paths are almost
surely continuous and that has initial distribution ν,
(b) {Xi(t)−Xi(0)− µit,Ft, t≥ 0} is a martingale for i= 1, . . . , d, and
(c) {(Xi(t)−Xi(0)−µit)(Xj(t)−Xj(0)−µjt)−Γijt,Ft, t≥ 0} is a mar-
tingale for i, j = 1, . . . , d.
In this definition, the filtration {Ft} may be larger than the one generated
by X ; however, for each t≥ 0, under P , the σ-algebra Ft is independent of
the increments of X from t onward. The latter follows from the martingale
properties of X . If ν = δx, the unit mass at x ∈ R
d, we say that X starts
from x.
2. Definition of an SRBM. Let G=
i∈I Gi be a nonempty domain in
d, where I is a nonempty finite index set and for each i ∈ I , Gi is a
nonempty domain in Rd. For simplicity, we assume that I = {1,2, . . . , I}
and then |I|= I. For each i ∈ I , let γi(·) be a vector valued function defined
from Rd into Rd. Fix µ ∈Rd, Γ a d × d symmetric and strictly positive
definite covariance matrix and ν a probability measure on (G,B(G)), where
B(G) denotes the σ-algebra of Borel subsets of the closure G of G.
Definition 2.1 (Semimartingale reflecting Brownian motion). A semi-
martingale reflecting Brownian motion (abbreviated as SRBM) associated
with the data (G,µ,Γ,{γi, i ∈ I}, ν) is an {Ft}-adapted, d-dimensional pro-
cess W defined on some filtered probability space (Ω,F ,{Ft}, P ) such that:
(i) P -a.s., W (t) =X(t) +
(0,t] γ
i(W (s))dYi(s) for all t≥ 0,
(ii) P -a.s., W has continuous paths and W (t) ∈G for all t≥ 0,
(iii) under P , X is a d-dimensional {Ft}-Brownian motion with drift
vector µ, covariance matrix Γ and initial distribution ν,
(iv) for each i ∈ I , Yi is an {Ft}-adapted, one-dimensional process such
that P -a.s.,
(a) Yi(0) = 0,
(b) Yi is continuous and nondecreasing,
INVARIANCE PRINCIPLE FOR SRBMS 7
(c) Yi(t) =
(0,t] 1{W (s)∈∂Gi∩∂G} dYi(s) for all t≥ 0.
We shall often refer to Y = {Yi, i ∈ I} as the “pushing process” associated
with the SRBM W . When ν = δx, we may alternatively say that W is an
SRBM associated with the data (G,µ,Γ,{γi, i ∈ I}) that starts from x. We
will call (W,X,Y ) satisfying Definition 2.1 an extended SRBM associated
with the data (G,µ,Γ,{γi, i ∈ I}, ν).
Loosely speaking, an SRBM behaves like a Brownian motion in the inte-
rior of the domain G and it is confined to G by instantaneous “reflection” (or
“pushing”) at the boundary, where the allowed directions of “reflection” at
x ∈ ∂G are convex combinations of the vectors γi(x) for i such that x ∈ ∂Gi.
Under the assumptions imposed on G and {γi, i ∈ I} in Sections 3.1 and 3.2
below, at each point on the boundary of G there is an allowed direction of
reflection that can be used there which “points into the interior of G.” We
end this section by introducing a related set-valued function I(·) and show
a key property of it.
Definition 2.2. For each x∈Rd, let I(x) = {i ∈ I :x ∈ ∂Gi}.
The set-valued function I(·) has the following property called upper semi-
continuity on ∂G.
Lemma 2.1. For each x ∈ ∂G, there is an open neighborhood Vx of x in
d such that
I(y)⊂ I(x) for all y ∈ Vx.(4)
Proof. We prove this lemma by contradiction. Suppose that the func-
tion I(·) does not satisfy (4). Then there is a point x ∈ ∂G such that there
is no open neighborhood Vx of x such that I(y)⊂ I(x) for all y ∈ Vx. Since
the index set I is finite, there is an index k ∈ I \ I(x) and a sequence of
points {yn} ⊂R
d such that ‖yn−x‖<
and k ∈ I(yn) for each n≥ 1. Hence
yn ∈ ∂Gk for all n≥ 1. Since ∂Gk is closed and yn → x as n→∞, we con-
clude that x ∈ ∂Gk. This implies that k ∈ I(x), which is a contradiction, as
desired. �
3. Assumptions on the domain G and the reflection vector fields {γi}.
3.1. Assumptions on the domain G. We henceforth assume that the
domain G satisfies assumptions (A1)–(A3) below. In the case when G is
bounded, assumptions (A2)–(A3) follow from assumption (A1) (see Lem-
mas A.1 and A.2 in the Appendix for details). If the domain G is a convex
polyhedron satisfying assumption (A1), then assumptions (A2)–(A3) hold
by Lemma A.3 in the Appendix.
8 W. KANG AND R. J. WILLIAMS
(A1) G is a nonempty domain in Rd with representation
Gi,(5)
where for each i ∈ I , Gi is a nonempty domain, Gi 6=R
d, and the boundary
∂Gi of Gi is C
1. For each i ∈ I , we let ni(·) be the unit normal vector field
on ∂Gi that points into Gi.
(A2) For each ε ∈ (0,1) there exists R(ε) > 0 such that for each i ∈ I ,
x ∈ ∂Gi ∩ ∂G and y ∈G satisfying ‖x− y‖<R(ε), we have
〈ni(x), y − x〉 ≥−ε‖x− y‖.(6)
(A3) The function D : [0,∞)→ [0,∞] defined such that D(0) = 0 and
D(r) = sup
J 6=∅
(∂Gj ∩ ∂G)
Ur(∂Gj ∩ ∂G)
for r > 0, satisfies
D(r)→ 0 as r→ 0.(8)
Remark. Assumption (A2) is reminiscent of the uniform exterior cone
condition (cf. [9], page 195). We say that a region G ⊂ Rd satisfies a uni-
form exterior cone condition if for each x0 ∈ ∂G, there is a truncated closed
right circular cone Vx0 , with nonempty interior and vertex x0, satisfying
Vx0 ∩G= {x0}, and the truncated cones Vx0 are all congruent to some fixed
truncated closed right circular cone V . By comparing assumption (A2) with
the uniform exterior cone condition, we see that assumption (A2) implies
the uniform exterior cone condition. On the other hand, under assumption
(A1), assumption (A2) is implied by a family of uniform exterior cone condi-
tions where for each ε ∈ (0,1), the axis of the truncated closed right circular
cone at x ∈ ∂G is along the vector −ni(x) and all of the truncated closed
right circular cones are congruent to a truncated closed right circular cone
whose height and base radius are R(ε) and R(ε)( 1
− 1)1/2 respectively.
Assumption (A2) holds automatically if G is convex. We also note that as-
sumption (A2) is strictly weaker than the uniform exterior sphere condition.
The definition of the uniform exterior sphere condition is similar to that for
the uniform exterior cone condition where a closed ball with x0 on its bound-
ary takes the place of the truncated closed right circular cone Vx0 . It can
be checked that for the domain G = {(x, y) ∈ R2 :y < |x|α} with α ∈ (1,2),
the uniform exterior sphere condition fails to hold, but assumption (A2)
holds. In fact, at the point (0,0) ∈R2, there is no r > 0 and y ∈R2 such that
Br(y)∩ ∂G= {(0,0)}.
INVARIANCE PRINCIPLE FOR SRBMS 9
Remark. For the definition of D(·) in (A3), we adopt the convention
that the supremum over an empty set is zero and dist(x,∅) =∞. Since ∂Gi∩
∂G 6=∅ for at least one i ∈ I , the function D(·) satisfies limr→∞D(r) =∞.
Furthermore, D(r1)≤D(r2) whenever r1, r2 ∈ [0,∞) and r1 ≤ r2. Assump-
tion (A3) requires that for any nonempty subset J ⊂ I , the intersection of
tubular neighborhoods of the boundaries {∂Gj ∩ ∂G : j ∈ J } given by the
j∈J Ur(∂Gj ∩ ∂G) “converges” to the intersection of the boundaries
given by the set
j∈J (∂Gj ∩ ∂G) as r approaches 0. Property (8) need
not always hold. For example, let G1 = {(x, y) ∈ R
2 :y < e−x
2/2, x ∈ R} and
G2 = {(x, y) ∈ R
2 :y > 0, x ∈ R}. Then ∂G1 ∩ ∂G2 =∅. But for each r > 0,
Ur(∂G1)∩Ur(∂G2) 6=∅. Hence D(r) =∞ for each r > 0.
3.2. Assumptions on the reflection vector fields {γi}. We henceforth as-
sume that there are vector fields {γi(·), i ∈ I} satisfying assumptions (A4)–
(A5) below.
(A4) There is a constant L > 0 such that for each i ∈ I , γi(·) is a uni-
formly Lipschitz continuous function from Rd into Rd with Lipschitz con-
stant L and ‖γi(x)‖= 1 for each x ∈Rd.
(A5) There is a constant a ∈ (0,1), and vector valued functions b(·) =
(b1(·), . . . , bI(·)) and c(·) = (c1(·), . . . , cI(·)) from ∂G into R
+ such that for
each x ∈ ∂G,
i∈I(x) bi(x) = 1,
j∈I(x)
i∈I(x)
bi(x)n
i(x), γj(x)
≥ a,(9)
i∈I(x) ci(x) = 1,
j∈I(x)
i∈I(x)
ci(x)γ
i(x), nj(x)
≥ a.(10)
We note here for future use that by (A4), if we set ρ0 =
, then for
any x, y ∈ Rd satisfying ‖x − y‖ < ρ0, we have ‖γ
i(x) − γi(y)‖ < a/4 for
each i ∈ I . So for each 0< ρ < ρ0/4, by (9)–(10) and the normalization of
b(·), c(·), γi(·), nj(·) for i, j ∈ I , we obtain
j∈I(x)
y∈B4ρ(x)
i∈I(x)
bi(x)n
i(x), γj(y)
≥ a/2(11)
j∈I(x)
y∈B4ρ(x)
i∈I(x)
ci(x)γ
i(y), nj(x)
≥ a/2.(12)
10 W. KANG AND R. J. WILLIAMS
The use of B4ρ(x) here is related to the form in which this is used in Section
Remark. Assumption (A4) is equivalent to (3.4) in [6] whenG is bounded.
Property (10) means that, at each point x ∈ ∂G, there is a convex combi-
nation γ(x) =
i∈I(x) ci(x)γ
i(x) of the vectors {γi(x), i ∈ I(x)} that can
be used there such that γ(x) “points into” G. Property (9) is in a sense a
dual condition to property (10), where the roles of γi and ni are reversed
for i ∈ I(x). This property (9) is used in showing the oscillation inequality
in Theorem 4.1 below. Assumption (A5) is an analogue of Assumption 1.1
in [4]. When G is bounded, (10) is similar to condition (3.6) in [6] (we as-
sume some additional uniformity through the lack of dependence of a on
It is straightforward to see using the triangle inequality that the following
condition (A5)′ implies (A5).
(A5)′ There is a ∈ (0,1) and vector valued functions b, c from ∂G into RI+
such that for each x ∈ ∂G,
i∈I(x) bi(x) = 1, and for each i ∈ I(x),
bi(x)〈n
i(x), γi(x)〉 ≥ a+
j∈I(x)\{i}
bj(x)|〈n
j(x), γi(x)〉|,(13)
i∈I(x) ci(x) = 1, and for each i ∈ I(x),
ci(x)〈γ
i(x), ni(x)〉 ≥ a+
j∈I(x)\{i}
cj(x)|〈γ
j(x), ni(x)〉|.(14)
Condition (A5)′(ii) is similar to condition (3.8) in [6], although here we
assume additional uniformity through the lack of dependence of a on x. As
noted in [6], their condition (3.8) can be phrased in terms of a nonsingu-
lar M-matrix requirement [2]. (This is sometimes also called a generalized
Harrison–Reiman type of condition [10].) Since that nonsingular M-matrix
property is invariant under transpose, and this property for the transpose
corresponds to a local form of (A5)′(i), one might conjecture that there is an
equivalence between the existence of a nonnegative vector valued function b
such that (A5)′(i) holds for each x ∈ ∂G and the existence of a nonnegative
vector valued function c such that (A5)′(ii) holds for each x ∈ ∂G. Indeed
we have the following lemma. We have stated the two (equivalent) condi-
tions (i) and (ii) in specifying (A5)′ to preserve a parallel with (A5) and
since both properties can be useful in proofs. Furthermore, in light of the
following lemma, verifying either condition suffices for both to hold.
INVARIANCE PRINCIPLE FOR SRBMS 11
Lemma 3.1. There is a constant a ∈ (0,1) and a vector valued function
b :∂G → RI+ such that (A5)
′(i) holds for each x ∈ ∂G if and only if there
is a constant a ∈ (0,1) and a vector valued function c :∂G→ RI+ such that
(A5)′(ii) holds for each x ∈ ∂G.
Proof. We just prove the “if” part; the “only if” part can be proved
in a similar manner.
We suppose that there is a constant a ∈ (0,1) and a vector valued function
c :∂G → RI+ such that (A5)
′(ii) holds for each x ∈ ∂G. For fixed x ∈ ∂G,
consider the square matrix A(x) whose diagonal entries are given by the
(positive) elements 〈ni(x), γi(x)〉 for i ∈ I(x) and whose off-diagonal entries
are given by−|〈ni(x), γj(x)〉| for i ∈ I(x), j ∈ I(x), j 6= i. Let E be the square
matrix having the same dimensions as A(x) and whose entries are all equal
to one. By the theory of M-matrices (see [2], Chapter 6, especially condition
(M35)), condition (ii) of (A5)
′ implies that A(x)− a
E is a nonsingular M-
matrix, that is, A(x)− a
E has nonnegative diagonal entries and nonpositive
off-diagonal entries and it can be written in the form s(x)I −B(x) where
B(x) is a matrix with nonnegative entries and s(x)> 0 is a constant that is
strictly larger than the spectral radius of B(x).
Since the nonsingular M-matrix property is invariant under transpose (cf.
(G21) in Chapter 6 of [2]), then A
′(x)− a
E is also a nonsingular M-matrix.
Hence, there is a vector b̃(x) = (b̃i(x) : i ∈ I(x)) with nonnegative entries such
that (A′(x)− a
E)b̃(x)> 0 (cf. (I27) in Chapter 6 of [2]). We can extend b̃(x)
to an I-dimensional vector b(x) and normalize it so that
i∈I(x) bi(x) = 1.
Then (A5)′(i) holds with a
in place of a. �
4. Invariance principle. In this section we state and prove an invariance
principle for an SRBM living in the closure of a domain G with piecewise
smooth boundary and having associated reflection fields {γi, i ∈ I}, where
G, {γi, i ∈ I} satisfy assumptions (A1)–(A5) of Section 3. (These assump-
tions hold throughout this section.) We shall first state a preliminary result
called an oscillation inequality (see Theorem 4.1), then we use it to prove
a tightness result (see Theorem 4.2). Finally, we establish the invariance
principle (see Theorem 4.3).
4.1. Oscillation inequality. The following oscillation inequality is the key
to the proof of the tightness result claimed in Theorem 4.2. In this subsec-
tion, for any 0≤ t1 < t2 <∞ and any integer k ≥ 1, D([t1, t2],R
k) denotes
the set of functions w : [t1, t2]→R
k that are right continuous on [t1, t2) and
have finite left limits on (t1, t2]. For w ∈D([t1, t2],R
Osc(w, [t1, t2]) = sup{‖w(t)−w(s)‖ : t1 ≤ s < t≤ t2},(15)
Osc(w, [t1, t2)) = sup{‖w(t)−w(s)‖ : t1 ≤ s < t < t2}.(16)
12 W. KANG AND R. J. WILLIAMS
Note that we do not explicitly indicate the dependence on k in the notation.
Recall the constants a,L from assumptions (A4)–(A5), the functions R(·)
from assumption (A2) and D(·) from (7). Let ρ0 =
Theorem 4.1 (Oscillation inequality). There exists a nondecreasing func-
tion Π: (0,∞)→ (0,∞] satisfying Π(u)→ 0 as u→ 0, such that Π depends
only on the constants I, a and the function D(·), and such that whenever
0 < ρ <min{
R(a/4)
}, 0 < δ <
, 0 ≤ s < t < ∞, w,x ∈ D([s, t],Rd) and
y ∈D([s, t],RI) satisfy:
(i) w(u) ∈Bρ(x0)∩Uδ(G) for all u ∈ [s, t], for some x0 ∈G,
(ii) w(u) = w(s) + x(u) − x(s) +
(s,u] γ
i(w(v))dyi(v) for all u ∈
[s, t],
(iii) for each i ∈ I ,
(a) yi(s)≥ 0,
(b) yi is nondecreasing and ∆yi(u)≤ δ for all u ∈ (s, t],
(c) yi(u) = yi(s) +
(s,u] 1{w(v)∈Uδ(∂Gi∩∂G)} dyi(v) for all u ∈ [s, t],
(iv) D(Π(Osc(x, [s, t]) + δ))<
then we have that the following hold:
Osc(w, [s, t])≤Π(Osc(x, [s, t]) + δ),(17)
Osc(y, [s, t])≤Π(Osc(x, [s, t]) + δ).(18)
Proof. Let
Π0(u) = u for all u > 0.
Define Πm : (0,∞)→ (0,∞], m= 1, . . . , I, inductively such that
Πm(u) = Πm−1(u) + (I+2)u+
(D(Πm−1(u) + (I+2)u) + 2u).
Here the sum of any element of [0,∞) with ∞ is ∞ and D(∞) is defined
to equal ∞. For each m= 0,1, . . . , I, the function Πm is nondecreasing and
depends only on I, a and D(·). For each m= 1, . . . , I and u > 0, Πm−1(u)≤
Πm(u). By assumption (A3), we conclude (using an induction proof) that
Πm(u)→ 0 as u→ 0, for m= 0,1, . . . , I.
Let Π(·) = ΠI(·).
Fix 0 < ρ < min{
R(a/4)
}, 0 < δ <
, 0 ≤ s < t < ∞. Suppose that
w,x ∈D([s, t],Rd) and y ∈D([s, t],RI) satisfy (i)–(iv) in the statement of
Theorem 4.1. For each nonempty interval [t1, t2]⊂ [s, t], let
I[t1,t2] = {i ∈ I :w(u) ∈ Uδ(∂Gi ∩ ∂G) for some u ∈ [t1, t2]},
INVARIANCE PRINCIPLE FOR SRBMS 13
the indices of the boundary surfaces that w(·) comes close to in the time
interval [t1, t2]. For each 0 ≤ m ≤ I, define Tm = {[t1, t2] ⊂ [s, t] : |I[t1,t2]| ≤
m}. Note that under the partial ordering of set inclusion, Tm increases with
m. To prove the theorem, we will prove by induction that for each 0≤m≤ I
and each interval [t1, t2] ∈ Tm, (17)–(18) hold with [t1, t2] in place of [s, t] and
Πm(·) in place of Π(·). The result for m= I yields the theorem.
Suppose that m= 0. Then T0 = {[t1, t2]⊂ [s, t] : |I[t1,t2]|= 0}. Fix an inter-
val [t1, t2] ∈ T0. Since I[t1,t2] =∅ and (iii)(c) holds, the function y does not
increase on the time interval (t1, t2], that is, yi(t2)− yi(t1) = 0 for all i ∈ I .
Then, for t1 ≤ u < v ≤ t2,
w(v)−w(u) = x(v)− x(u).(19)
So in this case,
Osc(w, [t1, t2]) = Osc(x, [t1, t2])≤Osc(x, [t1, t2]) + δ,(20)
Osc(y, [t1, t2]) = 0≤Osc(x, [t1, t2]) + δ.(21)
Thus, (17)–(18) hold with Π0(·) in place of Π(·) and [t1, t2] in place of [s, t]
for each interval [t1, t2] ∈ T0.
For the induction step, let 1 ≤ m ≤ I and suppose that (17)–(18) hold
with Πm−1(·) in place of Π(·) and [t1, t2] in place of [s, t] for each interval
[t1, t2] ∈ Tm−1.
Now fix [t1, t2] ∈ Tm. If |I[t1,t2]| ≤ m − 1, then [t1, t2] ∈ Tm−1 and so by
the induction assumption we have that (17)–(18) hold with [t1, t2] in place
of [s, t] and Πm−1(·) [and hence Πm(·)] in place of Π(·). Thus, it suffices
to consider [t1, t2] ⊂ [s, t] such that |I[t1,t2]| =m. For i /∈ I[t1,t2], by (iii)(c),
yi(t2)− yi(t1) = 0, and so by (ii), for t1 ≤ u < v ≤ t2, we have
w(v)−w(u) = x(v)− x(u) +
i∈I[t1,t2]
(u,v]
γi(w(r))dyi(r).(22)
Let Πm(u) = Πm−1(u) + (I + 2)u for all u > 0, and η = Osc(x, [t1, t2]) + δ.
For any M ∈ (0,∞] and any nonempty set J ⊂ I , let
FMJ = {z ∈R
d : dist(z, ∂Gi ∩ ∂G)<M for all i ∈ J }.
Note that FMJ = ∅ when there is an i ∈ J such that ∂Gi ∩ ∂G = ∅. Since
Πm(·)≤Πm(·)≤Π(·), D(·) and Π(·) are nondecreasing, and Osc(x, [t1, t2])≤
Osc(x, [s, t]), we have by (iv) that
D(Πm(η))≤D(Πm(η))≤D(Π(η))<
.(23)
Note that this implies Πm(η)<∞ since D(∞) =∞.
We now consider two cases.
14 W. KANG AND R. J. WILLIAMS
Case 1. Suppose that w(r) ∈ F
Πm(η)
I[t1,t2]
for all r ∈ [t1, t2].
Fix u, v such that t1 ≤ u < v ≤ t2. Since we have that
w(v) ∈
j∈I[t1,t2]
Πm(η)
(∂Gj ∩ ∂G),
by the definition of D(·) and (23), there is z ∈
j∈I[t1,t2]
(∂Gj ∩ ∂G) such
‖w(v)− z‖ ≤D(Πm(η))<
.(24)
For each r ∈ [t1, t2], by (i) we have that w(r) ∈ Uδ(G), and so there is z
such that
‖w(r)− zr‖ ≤ 2δ.
Hence by (i) and (24) we have
‖zr − z‖ ≤ ‖zr −w(r)‖+ ‖w(r)− x0‖+ ‖x0 −w(v)‖+ ‖w(v)− z‖
≤ 2δ + ρ+ ρ+ ρ/2< 4ρ < R(a/4)
‖w(r)− z‖ ≤ ‖w(r)− x0‖+ ‖x0 −w(v)‖+ ‖w(v)− z‖
≤ ρ+ ρ+ ρ/2< 4ρ.
By (6) and (25) we have
〈nj(z), z − zr〉 ≤
‖z − zr‖ for each j ∈ I(z) and r ∈ [t1, t2].(27)
Note that I(z) ⊃ I[t1,t2]. Recalling the definition of b(·) from assumption
(A5), on dotting the vector
j∈I(z) bj(z)n
j(z) with both sides of (22) and
rearranging, we obtain
i∈I[t1,t2]
(u,v]
j∈I(z)
bj(z)n
j(z), γi(w(r))
dyi(r)
j∈I(z)
bj(z)〈n
j(z),w(v)−w(u)〉(28)
j∈I(z)
bj(z)〈n
j(z), x(v)− x(u)〉.
So by (11), (22), (24)–(28), and the fact that
j∈I(z) bj(z) = 1, bj(z)≥ 0 for
j ∈ I , we have
i∈I[t1,t2]
(yi(v)− yi(u))
INVARIANCE PRINCIPLE FOR SRBMS 15
i∈I[t1,t2]
(u,v]
j∈I(z)
bj(z)n
j(z), γi(w(r))
dyi(r)
j∈I(z)
bj(z)〈n
j(z),w(v)− z〉+
j∈I(z)
bj(z)〈n
j(z), z − zu〉
j∈I(z)
bj(z)〈n
j(z), zu −w(u)〉 −
j∈I(z)
bj(z)〈n
j(z), x(v)− x(u)〉
≤D(Πm(η)) +
‖z − zu‖+ 2δ + ‖x(v)− x(u)‖
≤D(Πm(η)) + 2δ + ‖x(v)− x(u)‖
(‖z −w(v)‖+ ‖w(v)−w(u)‖+ ‖w(u)− zu‖)
≤D(Πm(η)) + 2δ + ‖x(v)− x(u)‖
D(Πm(η)) + ‖x(v)− x(u)‖+
i∈I[t1,t2]
(yi(v)− yi(u)) + 2δ
{D(Πm(η)) + 2δ + ‖x(v)− x(u)‖}+
i∈I[t1,t2]
(yi(v)− yi(u)).
Hence
i∈I[t1,t2]
(yi(v)− yi(u))≤
{D(Πm(η)) + 2δ + ‖x(v)− x(u)‖}
{D(Πm(η)) + 2η}.
On multiplying through by 4
, we obtain
i∈I[t1,t2]
(yi(v)− yi(u))≤
{D(Πm(η)) + 2η} ≤Πm(η).(29)
Hence, by (29) and the fact that for any x ∈Rd, ‖x‖ ≤
i=1 |xi|, we have
Osc(y, [t1, t2])≤Πm(Osc(x, [t1, t2]) + δ),(30)
and by (22), (29) and the definitions of Πm(·) and Πm(·), we have
Osc(w, [t1, t2])≤Osc(x, [t1, t2]) +
{D(Πm(η)) + 2η}
≤Πm(Osc(x, [t1, t2]) + δ),
as desired.
16 W. KANG AND R. J. WILLIAMS
Case 2. Suppose that there is t3 ∈ [t1, t2] such that w(t3) /∈ F
Πm(η)
I[t1,t2]
Define σ = inf{u ∈ [t1, t2] :w(u) /∈ F
Πm(η)
I[t1,t2]
}. Then σ ≤ t2. For each u ∈
[t1, σ), w(u) ∈ F
Πm(η)
I[t1,t2]
and so by a similar analysis to that for Case 1, we
obtain for each v ∈ [t1, σ),
Osc(w, [t1, v])≤ η+
(D(Πm(η)) + 2η)
Osc(y, [t1, v])≤
(D(Πm(η)) + 2η).
By the right continuity of paths we have w(σ) /∈ F
Πm(η)
I[t1,t2]
. Then there is an
i ∈ I[t1,t2] such that dist(w(σ), ∂Gi ∩ ∂G) ≥ Πm(η), and it follows that w
does not reach Uδ(∂Gi ∩ ∂G) during the interval [σ, t2]. To see this, let
τ = inf{u ∈ [σ, t2] : dist(w(u), ∂Gi ∩ ∂G) ≤ δ} with the convention that the
infimum of an empty set is ∞. If τ ≤ t2, then by the right continuity of
w(·) and since Πm(η)> δ, we have τ > σ and dist(w(τ), ∂Gi∩∂G)≤ δ. Also,
since |I[t1,t2]|=m, we have [σ,u] ∈ Tm−1 for each u ∈ [σ, τ). By the induction
assumption and letting u→ τ , we have ‖w(τ−)−w(σ)‖ ≤Πm−1(η). By (ii),
(iii)(b) and since ‖γi(·)‖= 1, we have
‖∆w(τ)‖ ≤ ‖∆x(τ)‖+
∆yi(τ)≤Osc(x, [t1, t2]) + Iδ ≤ Iη.
Then simple manipulations yield
dist(w(σ), ∂Gi ∩ ∂G)≤ ‖w(σ)−w(τ−)‖+ ‖∆w(τ)‖+ dist(w(τ), ∂Gi ∩ ∂G)
≤Πm−1(η) + Iη+ δ
<Πm(η).
This contradicts the fact that dist(w(σ), ∂Gi∩∂G)≥Πm(η), and so confirms
that w does not reach Uδ(∂Gi ∩ ∂G) in [σ, t2]. Thus we must have [σ, t2] ∈
Tm−1. Hence we have by the induction assumption that
Osc(w, [t1, t2])≤ sup
v∈[t1,σ)
Osc(w, [t1, v]) + ‖∆w(σ)‖+Osc(w, [σ, t2])
≤ η +
(D(Πm(η)) + 2η) + Iη+Πm−1(η)
≤Πm(Osc(x, [t1, t2]) + δ)
INVARIANCE PRINCIPLE FOR SRBMS 17
Osc(y, [t1, t2])≤ sup
v∈[t1,σ)
Osc(y, [t1, v]) + ‖∆y(σ)‖+Osc(y, [σ, t2])
(D(Πm(η)) + 2η) + Iη+Πm−1(η)
≤Πm(Osc(x, [t1, t2]) + δ).
On combining all of the cases above, we have
Osc(w, [t1, t2])≤Πm(Osc(x, [t1, t2]) + δ),(31)
Osc(y, [t1, t2])≤Πm(Osc(x, [t1, t2]) + δ).(32)
This completes the induction step. �
Remark. The proof of the above theorem was inspired by the proof of
Lemma 4.3 of [4]. Because of the condition (i) in Theorem 4.1, the oscilla-
tion inequality given here is localized. Similar, but nonlocalized, oscillation
inequalities were proved in [15] when G = Rd+ and in [3] for a sequence of
convex polyhedrons; in these cases, the direction of reflection was constant
on each boundary face.
4.2. C-tightness result. Throughout this subsection and the next, we
suppose that the following assumption holds in addition to (A1)–(A5).
Assumption 4.1. There is a sequence of strictly positive constants {δn}∞n=1
such that for each positive integer n, there are processes W n, W̃ n,Xn, αn
having paths in D([0,∞),Rd) and processes Y n, Ỹ n, βn having paths in
D([0,∞),RI) defined on some probability space (Ωn,Fn, Pn) such that:
(i) Pn-a.s., W n = W̃ n + αn and W̃ n(t) ∈Uδn(G) for all t≥ 0,
(ii) Pn-a.s., W n(t) =Xn(t)+
(0,t] γ
i,n(W n(s−),W n(s))dY ni (s) for
all t≥ 0, where for each i ∈ I , γi,n :Rd ×Rd → Rd is Borel measurable and
‖γi,n(y,x)‖= 1 for all x, y ∈Rd,
(iii) Y n = Ỹ n+βn, where βn is locally of bounded variation and Pn-a.s.,
for each i ∈ I ,
(a) Ỹ ni (0) = 0,
(b) Ỹ ni is nondecreasing and ∆Ỹ
i (t)≤ δ
n for all t > 0,
(c) Ỹ ni (t) =
(0,t] 1{W̃n(s)∈Uδn (∂Gi∩∂G)}
dỸ ni (s),
(iv) δn → 0 as n → ∞, and for each ε > 0, there is ηε > 0 and nε > 0
such that for each i ∈ I , ‖γi,n(y,x)− γi(x)‖< ε whenever ‖x− y‖< ηε and
n≥ nε,
18 W. KANG AND R. J. WILLIAMS
(v) αn → 0 and V(βn)→ 0 in probability, as n→∞,
(vi) {Xn} is C-tight.
Remark. A simple case in which (iv) above holds is where γi,n(y,x)≡
γi(y). In (v), V(βn) is the total variation process for βn (cf. Section 1.1).
The following theorem will play an important role in the proof of the in-
variance principle. It will be used to show that a sequence of processes sat-
isfying suitably perturbed versions of the defining conditions for an SRBM
[cf. (i)–(vi) above] is C-tight.
Theorem 4.2 (C-tightness). Suppose that Assumption 4.1 holds. Define
Zn = (W n,Xn, Y n) for each n. Then the sequence of processes {Zn}∞n=1 is
C-tight.
Remark. Note that C-tightness of {W n}, {Xn} and {Y n} implies C-
tightness of {Zn} (see Chapter VI, Corollary 3.33 of [12] for details).
Proof of Theorem 4.2. References here to (i)–(vi) are to the condi-
tions in Assumption 4.1.
Simple algebraic manipulations yield Pn-a.s.,
W̃ n(t) = X̃n(t) +
(0,t]
γi,n(W n(s−),W n(s))dỸ ni (s)(33)
= X̃n(t) + Ṽ n(t) +
(0,t]
γi(W̃ n(s))dỸ ni (s),(34)
where
X̃n(t) =Xn(t) +
−αn(t) +
(0,t]
γi,n(W n(s−),W n(s))dβni (s)
Ṽ n(t) =
(0,t]
(γi,n(W n(s−),W n(s))− γi(W n(s)))dỸ ni (s)
(0,t]
(γi(W n(s))− γi(W̃ n(s)))dỸ ni (s).
The hypotheses on αn, the total variation process V(βn) of βn, and the
fact that ‖γi,n(y,x)‖ = 1 for all x, y ∈ Rd and each i ∈ I , imply that the
process
−αn(·) +
(0,·]
γi,n(W n(s−),W n(s))dβni (s)
INVARIANCE PRINCIPLE FOR SRBMS 19
converges to 0 in probability as n→∞. Combining this with the fact that
{Xn}∞n=1 is C-tight, we obtain that {X̃
n}∞n=1 is C-tight.
Recall the positive nondecreasing function Π(·) from Theorem 4.1, and
the constants a, L and functions R(·) and D(·) from assumptions (A1)–(A5)
in Section 3. Recall also that ρ0 =
Fix ρ, ε, η, T such that 0 < ρ <min{
R(a/4)
}, ε > 0, η > 0 and T > 0.
By assumption (A3), there is a constant r1 > 0 such that
D(r)<min
for all r ∈ (0, r1].(37)
Since Π(u) → 0 as u → 0, there are constants 0 < r3 < r2 < min{r1,
such that
Π(r)<
for all r ∈ (0, r3].(38)
By (iv), there are 0< ε̃ <min{
} and n0 > 0 such that for all n≥ n0,
‖y−x‖<2ε̃
‖γi,n(y,x)− γi(x)‖<
.(39)
By (iv)–(vi), and Proposition 1.1, there exist an integer n1 > n0, a con-
stant M̃η,T > 0 and λ̃ ∈ (0, T ), such that for all n≥ n1,
0≤s≤T
‖X̃n(s)‖ ≤ M̃η,T
≥ 1− η/2,(40)
Pn{wT (X̃
n, λ̃)≥ ε̃} ≤ η/4,(41)
0≤s≤T
‖αn(s)‖<
6ILr2
≥ 1− η/4,(42)
δn <min
8(1 + I)
.(43)
To prove C-tightness of {W̃ n} and {Ỹ n} (and hence of {W n}, {Y n}),
by Proposition 1.1, it suffices to show that there exists a constant Nη,T > 0
such that for all n≥ n1,
Pn{wT (W̃
n, λ̃)≥ ε} ≤ η,(44)
Pn{wT (Ỹ
n, λ̃)≥ ε} ≤ η,(45)
0≤s≤T
‖W̃ n(s)‖ ≤Nη,T
≥ 1− η,(46)
0≤s≤T
‖Ỹ n(s)‖ ≤Nη,T
≥ 1− η.(47)
20 W. KANG AND R. J. WILLIAMS
For each n≥ 1, let Fn be a set in Fn such that Pn(Fn) = 1 and on Fn,
properties (iii)(a)–(c) hold, (33)–(36) hold, and W̃ n(t) ∈ Uδn(G) for all t≥ 0.
Fix a t such that 0≤ t < t+ λ̃≤ T . Let
τn = inf{s≥ t :W̃ n(s) ∈ Uδn(∂Gi ∩ ∂G) for some i ∈ I}.(48)
For each n≥ n1, let
wT (X̃
n, λ̃)< ε̃, sup
0≤s≤T
‖αn(s)‖<
6ILr2
0≤s≤T
‖X̃n(s)‖ ≤ M̃η,T
Then by (40)–(42) and the definition of Fn,
P{Hn} ≥ 1− η.(50)
Fix ωn ∈Hn. By the definition of wT (x,λ) in (3), we have that,
r,s∈[t,t+λ̃]
‖X̃n(s,ωn)− X̃n(r,ωn)‖< ε̃.(51)
Now there are two cases to consider for n ≥ n1 and u, v fixed such that
t≤ u < v ≤ t+ λ̃.
Case 1. ωn ∈ {τn > v}. In this case, by (iii)(c), Ỹ n(·, ωn) does not increase
on the interval (u, v], that is, Ỹ ni (v,ω
n)− Ỹ ni (u,ω
n) = 0 for all i ∈ I . Then
by (34) and (36),
W̃ n(v,ωn)− W̃ n(u,ωn) = X̃n(v,ωn)− X̃n(u,ωn).(52)
Hence, by (51),
‖W̃ n(v,ωn)− W̃ n(u,ωn)‖ ≤ sup
r,s∈[t,t+λ̃]
‖X̃n(s,ωn)− X̃n(r,ωn)‖< ε̃ < ε/8,
and we also have
‖Ỹ n(v,ωn)− Ỹ n(u,ωn)‖= 0< ε/2.
Case 2. ωn ∈ {τn ≤ v}. Then there is an i ∈ I such that W̃ n(τn, ωn) ∈
Uδn(∂Gi ∩ ∂G), since the set Uδn(∂Gi ∩ ∂G) is closed and W̃
n(·, ωn) is right
continuous. It follows that there is some x0 ∈ ∂G (which depends on ω
such that W̃ n(τn, ωn) is in the closed ball Bδn(x0)⊂Bρ(x0). To apply the
INVARIANCE PRINCIPLE FOR SRBMS 21
oscillation inequality in Theorem 4.1, we first prove the following:
W̃ n(r,ωn) ∈Bρ(x0) for all r satisfying τ
n ≤ r ≤ v.(53)
For the proof of (53), let
ξn = inf{r ≥ τn :W̃ n(r,ωn) /∈Bρ(x0)} ∧ v.(54)
By the definition of ξn, W̃ n(r,ωn) ∈Bρ(x0) for each r ∈ [τ
n, ξn). In order to
apply the oscillation inequality in Theorem 4.1 on the time interval [τn, ξn),
we show that
D(Π(Osc(X̃n(·, ωn) + Ṽ n(·, ωn), [τn, ξn)) + δn))<
.(55)
For each r ∈ (0, T ], by (i)–(iii) and (33), (49), (43), we have that
‖W n(r−, ωn)−W n(r,ωn)‖
≤ ‖W̃ n(r−, ωn)− W̃ n(r,ωn)‖+ ‖αn(r−, ωn)−αn(r,ωn)‖
≤ ‖∆X̃n(r,ωn)‖+2 sup
0≤s≤T
‖αn(s)‖+ Iδn
≤ ε̃+
< 2ε̃.
Hence by (39), for each r ∈ (0, T ],
‖γi,n(W n(r−, ωn),W n(r,ωn))− γi(W n(r,ωn))‖ ≤
.(56)
By (36), (56), Assumption (A4), (i) and (49), we have that for any s1, s2
such that u≤ s1 < s2 ≤ v,
‖Ṽ n(s2, ω
n)− Ṽ n(s1, ω
(s1,s2]
‖γi,n(W n(r−, ωn),W n(r,ωn))
− γi(W n(r,ωn))‖dỸ ni (r,ω
(s1,s2]
‖γi(W n(r,ωn))− γi(W̃ n(r,ωn))‖dỸ ni (r,ω
(Ỹ ni (s2, ω
n)− Ỹ ni (s1, ω
(s1,s2]
L‖W n(r,ωn)− W̃ n(r,ωn)‖dỸ ni (r,ω
(Ỹ ni (s2, ω
n)− Ỹ ni (s1, ω
6ILr2
(Ỹ ni (s2, ω
n)− Ỹ ni (s1, ω
22 W. KANG AND R. J. WILLIAMS
‖Ỹ n(s2, ω
n)− Ỹ n(s1, ω
σn = inf{s≥ τn :Osc(Ỹ n(·, ωn), [τn, s))> r2}.(58)
Note that Osc(Ỹ n(·, ωn), [τn, s)) as a function of s defined on (τn,∞) is
left continuous with finite right limits and is nondecreasing. By the right
continuity of Ỹ n, we know that
Osc(Ỹ n(·, ωn), [τn, s))→ 0 as s ↓ τn.
Thus, σn > τn, Osc(Ỹ n(·, ωn), [τn, σn))≤ r2 and on {σ
n <∞}, Osc(Ỹ n(·, ωn),
[τn, σn]) ≥ r2. By (57), (51), (43), the choice of ε, and since t≤ τ
n ≤ ξn ≤
v ≤ t+ λ̃, we have
Osc(X̃n(·, ωn) + Ṽ n(·, ωn), [τn, ξn ∧ σn)) + δn
≤Osc(X̃n(·, ωn), [τn, ξn ∧ σn))
+Osc(Ṽ n(·, ωn), [τn, ξn ∧ σn)) + δn
≤Osc(X̃n(·, ωn), [τn, ξn ∧ σn))
Osc(Ỹ n(·, ωn), [τn, ξn ∧ σn)) + δn
≤ ε̃+
r2 + δ
n < r3.
Then by (38) and the monotonicity of D(·), we have
D(Π(Osc(X̃n(·, ωn) + Ṽ n(·, ωn), [τn, ξn ∧ σn)) + δn))
≤D(r2)≤D(r1)<
We claim that
σn ≥ ξn.(61)
To prove (61), we proceed by contradiction and suppose that σn < ξn. Then
by (60), with x= X̃n(·, ωn) + Ṽ n(·, ωn) and δ = δn, condition (iv) of The-
orem 4.1 holds with [s, t] = [τn, σn − 1/m] for all m sufficiently large. By
applying Theorem 4.1 and letting m→∞, we obtain using (34), (38) and
(59) that,
Osc(Ỹ n(·, ωn), [τn, σn))
≤Π(Osc(X̃n(·, ωn) + Ṽ n(·, ωn), [τn, ξn ∧ σn)) + δn)(62)
≤Π(r3)<
INVARIANCE PRINCIPLE FOR SRBMS 23
By (62), (iii)(b) and (43), we obtain that
Osc(Ỹ n(·, ωn), [τn, σn])≤
+ Iδn < r2.
This contradicts the fact that Osc(Ỹ n(·, ωn), [τn, σn]) ≥ r2 on {σ
n < ∞},
and so (61) holds and (55) follows by (60).
By applying Theorem 4.1 on [τn, ξn − 1/m] and then letting m→∞, we
obtain using (61), (59) and (38), that
Osc(W̃ n(·, ωn), [τn, ξn))
≤Π(Osc(X̃n(·, ωn) + Ṽ n(·, ωn), [τn, ξn ∧ σn)) + δn)
and similarly,
Osc(Ỹ n(·, ωn), [τn, ξn))<
.(63)
Then we have
‖W̃ n(ξn−, ωn)− x0‖
≤ ‖W̃ n(ξn−, ωn)− W̃ n(τn, ωn)‖+ ‖W̃ n(τn, ωn)− x0‖
+ δn.
Using hypotheses (ii), (iii)(b), and (33), (51), we obtain
‖W̃ n(ξn, ωn)− W̃ n(ξn−, ωn)‖
≤ ‖X̃n(ξn, ωn)− X̃n(ξn−, ωn)‖
‖γi,n(W n(ξn−, ωn),W n(ξn, ωn))‖
× (Ỹ ni (ξ
n, ωn)− Ỹ ni (ξ
n−, ωn))
≤ ε̃+ Iδn.
Hence
‖W̃ n(ξn, ωn)− x0‖ ≤ ‖W̃
n(ξn−, ωn)− x0‖
+ ‖W̃ n(ξn, ωn)− W̃ n(ξn−, ωn)‖
+ δn + ε̃+ Iδn
≤ ε̃+ (I+1)δn +
< ρ/8 + ρ/8 + ρ/8< ρ/2.
24 W. KANG AND R. J. WILLIAMS
It follows from this that ξn = v and (53) holds, as desired.
Then, by (33), (51), (iii)(b), (iii)(c), (63) and (43), we have
‖W̃ n(v,ωn)− W̃ n(u,ωn)‖
≤ sup
r,s∈[u,v]
‖X̃n(s,ωn)− X̃n(r,ωn)‖+
(Ỹ ni (v,ω
n)− Ỹ ni (u,ω
≤ ε̃+
(Ỹ ni (v,ω
n)− Ỹ ni (u∨ τ
n, ωn))
(Ỹ ni (u ∨ τ
n, ωn)− Ỹ ni (u,ω
n))(64)
≤ ε̃+ IOsc(Ỹ n(·, ωn), [u∨ τn, v)) +
∆Ỹ ni (v,ω
n) + Iδn
≤ ε̃+ I
+ Iδn + Iδn <
‖Ỹ n(v,ωn)− Ỹ n(u,ωn)‖ ≤
(Ỹ ni (v,ω
n)− Ỹ ni (u,ω
(Ỹ ni (v,ω
n)− Ỹ ni (u∨ τ
n, ωn))
(Ỹ ni (u ∨ τ
n, ωn)− Ỹ ni (u,ω
Here we have used the fact that Ỹi does not increase on (u, τ
n ∨ u) and can
jump at most by δn at τn, by the definition of τn and (iii)(c).
On combining the results from Case 1 and Case 2, we obtain that for each
n≥ n1,
u,v∈[t,t+λ̃]
‖W̃ n(v,ωn)− W̃ n(u,ωn)‖ : 0≤ t≤ t+ λ̃≤ T
< ε(66)
u,v∈[t,t+λ̃]
‖Ỹ n(v,ωn)− Ỹ n(u,ωn)‖ : 0≤ t≤ t+ λ̃≤ T
< ε.(67)
Hence since ωn ∈Hn was arbitrary, by (50), we have that (44) and (45) hold
for all n≥ n1.
Next we show that there is a constant Nη,T > 0 such that (46) and (47)
hold for all n≥ n1. By (66)–(67) above, we have that for each n≥ n1, ω
INVARIANCE PRINCIPLE FOR SRBMS 25
Hn, t such that 0≤ t < t+ λ̃≤ T and t≤ u < v ≤ t+ λ̃,
‖W̃ n(v,ωn)− W̃ n(u,ωn)‖< ε(68)
‖Ỹ n(v,ωn)− Ỹ n(u,ωn)‖< ε.(69)
Then, for each 0≤ s≤ T , by (68), (69), (49) and (33), we have
‖W̃ n(s,ωn)‖ ≤ ‖W̃ n(s,ωn)− W̃ n(0, ωn)‖+ ‖W̃ n(0, ωn)‖
[T/λ̃]+1∑
‖W̃ n(iλ̃∧ s,ωn)− W̃ n((i− 1)λ̃ ∧ s,ωn)‖+ ‖X̃n(0, ωn)‖
≤ ([T/λ̃] + 1)ε+ M̃η,T
‖Ỹ n(s,ωn)‖ ≤ ‖Ỹ n(s,ωn)− Ỹ n(0, ωn)‖
[T/λ̃]+1∑
‖Ỹ n(iλ̃ ∧ s,ωn)− Ỹ n((i− 1)λ̃∧ s,ωn)‖
≤ ([T/λ̃] + 1)ε.
Here [T/λ̃] is the greatest integer less than or equal to T/λ̃. Let Nη,T =
([T/λ̃] + 1)ε+ M̃η,T . Then we obtain that for n≥ n1 and ω
n ∈Hn,
0≤s≤T
‖W̃ n(s,ωn)‖ ≤Nη,T(70)
0≤s≤T
‖Ỹ n(s,ωn)‖ ≤Nη,T .(71)
Then by (50), we have that (46) and (47) hold for all n≥ n1.
Finally by applying Proposition 1.1, we have the C-tightness of {W̃ n}
and {Ỹ n}. It then follows that {(W̃ n,Xn, Ỹ n)}∞n=1 is C-tight. Since Z
(W̃ n,Xn, Ỹ n) + (αn,0, βn) where αn,V(βn) → 0 in probability as n→∞,
then {Zn}∞n=1 is also C-tight. �
4.3. Invariance principle for SRBMs. The main theorem of the paper is
the following.
Theorem 4.3 (Invariance principle for SRBMs). Suppose that Assump-
tion 4.1 holds. Define Zn = (W n,Xn, Y n) for each n. Then the sequence of
26 W. KANG AND R. J. WILLIAMS
processes {Zn}∞n=1 is C-tight and any (weak) limit point of this sequence is
of the form Z = (W,X,Y ) where continuous d-dimensional processes W,X
and a continuous I-dimensional process Y are defined on some probability
space (Ω,F , P ) such that conditions (i), (ii) and (iv) of Definition 2.1 hold
with Ft = σ{Z(s) : 0≤ s≤ t}, t≥ 0.
If, in addition, the following conditions (vi)′ and (vii) hold, then any
weak limit point of the sequence {Zn}∞n=1 is an extended SRBM associated
with the data (G,µ,Γ,{γi, i ∈ I}, ν). If furthermore the following condition
(viii) holds, then W n ⇒W as n→∞ where W is an SRBM associated with
(G,µ,Γ,{γi, i ∈ I}, ν).
(vi)′ {Xn} converges in distribution to a d-dimensional Brownian mo-
tion with drift µ, covariance matrix Γ and initial distribution ν.
(vii) For each (weak) limit point Z = (W,X,Y ) of {Zn}∞n=1, {X(t) −
X(0)− µt, Ft, t≥ 0} is a martingale.
(viii) If a process W satisfies the properties in Definition 2.1, the law
of W is unique, that is, the law of an SRBM associated with the data
(G,µ,Γ,{γi, i ∈ I}, ν) is unique.
Remark. We note that (vi)′ implies that (vi) of Assumption 4.1 holds.
Proof of Theorem 4.3. By Theorem 4.2, we have that the sequence
{Zn}∞n=1 is C-tight. Let Z = (W,X,Y ) be a (weak) limit point of {Z
n}∞n=1,
that is, there is a subsequence {nk} of {n} such that Z
nk ⇒Z as k→∞. It
also follows that Z̃nk ≡ (W̃ nk ,Xnk , Ỹ nk)⇒Z as k→∞. By the C-tightness
of {Zn}, we obtain that Z has continuous paths a.s. For the purpose of
verifying that Z satisfies the listed properties in Definition 2.1, one may
invoke the Skorokhod representation theorem to assume, without loss of
generality, that Znk and Z̃nk converge u.o.c. to Z a.s. as k→∞ and V(βnk)
converges u.o.c. to 0 a.s. as k → ∞. With this simplification, it is easily
verified that the properties of {Znk} and {Z̃nk} imply that Z has properties
(ii) and (iv)(a)–(b) of Definition 2.1. For the verification of property (i) of
Definition 2.1, note that for each k, a.s. for each t≥ 0,
W nk(t) =Xnk(t) +
(0,t]
γi,nk(W nk(s−),W nk(s))dβ
i (s)
(0,t]
(γi,nk(W nk(s−),W nk(s))− γi(W nk(s)))dỸ
i (s)
(0,t]
γi(W nk(s))dỸ
i (s).
The sum of the first two terms on the right-hand side of the above equality
converges a.s. to X(t) as k → ∞. The third term on the right-hand side
INVARIANCE PRINCIPLE FOR SRBMS 27
converges a.s. to 0 as k→∞, by property (iv) and the fact that a.s.,
s∈(0,t]
‖W nk(s)−W nk(s−)‖
≤ sup
s∈(0,t]
‖∆Xnk(s)‖+ I sup
s∈(0,t]
‖∆Y nk(s)‖→ 0 as k→∞.
It remains to show that for each i ∈ I and t≥ 0, a.s.,
(0,t]
γi(W nk(s))dỸ
i (s)→
(0,t]
γi(W (s))dYi(s) as k→∞.
This follows directly from Lemma A.4.
For the verification of property (iv)(c) of Definition 2.1, it suffices to show
that for each i ∈ I , m= 1,2, . . . , a.s. for each t≥ 0,
Yi(t) =
(0,t]
fm(W (s))dYi(s),(72)
where {fm}
m=1 is a sequence of real valued continuous functions defined
on Rd such that for each m, the range of fm is [0,1], fm(x) = 1 for x ∈
U1/m(∂Gi ∩ ∂G) and fm(x) = 0 for x /∈ U2/m(∂Gi ∩ ∂G). The existence
of such a sequence of continuous functions {fm}
m=1 can be shown using
Urysohn’s lemma (cf. [8], page 122). Then (72) is a consequence of Lemma
A.4, property (iii) of Ỹ
i and the fact that δ
nk → 0 as k→∞. Indeed, a.s.,
for each t≥ 0,
Yi(t) = lim
i (t) = lim
(0,t]
{W̃nk (s)∈U
nk (∂Gi∩∂G)}
i (s)
= lim
(0,t]
fm(W̃
nk(s))dỸ
i (s)
(0,t]
fm(W (s))dYi(s).
Thus, Z satisfies properties (i), (ii) and (iv) of Definition 2.1 with Ft =
σ{Z(s) : 0≤ s≤ t}, t≥ 0.
Assuming properties (vi)′ and (vii) hold, Z satisfies (iii) of Definition
2.1. Then Z is an extended SRBM associated with the data (G,µ,Γ,{γi, i ∈
I}, ν). If in addition, property (viii) holds, then the law of W is unique. Since
each weak limit W is an SRBM associated with the data (G,µ,Γ,{γi, i ∈
I}, ν) and the law of such an SRBM is unique, then by a standard argument,
W n ⇒W as n→∞ where W is an SRBM associated with (G,µ,Γ,{γi, i ∈
I}, ν). �
Some sufficient conditions for (vii) to hold are given in Proposition 4.2 of
[15] for a simpler setting where G=Rd+. Two of those conditions generalize
to our setting here and can be proved in the same manner as in [15]. For
completeness, we state the ensuing result here.
28 W. KANG AND R. J. WILLIAMS
Proposition 4.1. Suppose that Assumption 4.1 and (vi)′ of Theorem
4.3 hold. If, in addition, one of the following conditions (I)–(II) holds, then
condition (vii) of Theorem 4.3 is satisfied, and any weak limit point of
{Zn}∞n=1 is an extended SRBM associated with (G,µ,Γ,{γ
i, i ∈ I}, ν).
(I) For any triple of d-dimensional {Ft}-adapted processes (W,X,Y )
defined on some filtered probability space (Ω,F ,{Ft}, P ) and satisfying con-
ditions (i), (ii) and (iv) of Definition 2.1 together with the condition that
X, under P , is a d-dimensional Brownian motion with drift vector µ, co-
variance matrix Γ and initial distribution ν, the pair (W,Y ) is adapted to
the filtration generated by X and the P -null sets.
(II) Xn = X̌n + εn1 , Y
n = Y̌ n + εn2 , W
n = W̌ n + εn3 , where ε
1 , ε
2 , ε
3 are
processes converging to 0 in probability as n→∞, and:
(a) {X̌n(t)− X̌n(0)}∞n=1 is uniformly integrable for each t≥ 0,
(b) there is a sequence of constants {µn}∞n=1 in R
d such that
limn→∞µ
n = µ,
(c) for each n, {X̌n(t)− X̌n(0)−µnt, t≥ 0} is a Pn-martingale with
respect to the filtration generated by (W̌ n, X̌n, Y̌ n).
In the rest of this work, we focus on applications of the invariance prin-
ciple and in particular on giving sufficient conditions for property (viii) of
Theorem 4.3 to hold.
5. Applications of the invariance principle. In Section 5.1, we prove weak
existence of SRBMs associated with data (G,µ,Γ,{γi, i ∈ I}, ν) satisfying
(A1)–(A5) of Section 3. This is accomplished by constructing a sequence of
approximations whose weak limit points are SRBMs. The invariance prin-
ciple is used to prove the C-tightness of the approximations and that any
weak limit point is an SRBM. In Sections 5.2 and 5.3, using known results
on uniqueness in law for SRBMs, we illustrate the invariance principle for
certain domains and directions of reflection.
5.1. Weak existence of SRBMs.
Theorem 5.1. Suppose that assumptions (A1)–(A5) of Section 3 hold.
Then there exists an SRBM associated with the data (G,µ,Γ,{γi, i ∈ I}, ν).
Proof. We construct a sequence of approximations to an SRBM and
use the invariance principle to establish weak convergence along a subse-
quence to an SRBM.
In the following we will use R(·) from assumption (A2), L > 0 from as-
sumption (A4), a > 0 from assumption (A5), and ρ0 =
. Fix ε > 0 and
INVARIANCE PRINCIPLE FOR SRBMS 29
0 < ρ < min{
R(a/4)
}. By assumption (A3), there is a constant r1 > 0
such that
D(r)<min
for all r ∈ (0, r1].
Recall the properties of Π(·) from Theorem 4.1. Since Π(u)→ 0 as u→ 0,
there are constants 0< r3 < r2 <min{r1,
} such that
Π(r)<
for all r ∈ (0, r3].
Fix ε̃ and δ such that 0< ε̃ <min{
24ILr2
} and 0< 2δ <min{r3
8(1+I)
We will construct a d-dimensional stochastic processW δ and an I-dimensional
“pushing” process Y δ , such that W δ approximately satisfies the conditions
defining an SRBM for the data (G,µ,Γ,{γi, i ∈ I}, ν) (cf. Assumption 4.1).
The idea for this construction is to use a Brownian motion X with drift
vector µ, covariance matrix Γ and initial distribution ν. Away from ∂G, the
increments of W δ are determined by those of X . For any time t ≥ 0 such
that W δ(t−) ∈ ∂G, we add an instantaneous jump to W δ(t−) to obtain
W δ(t) ∈G. Here W δ(0−) =X(0). The size of the jump is such that W δ(t)
is a strictly positive distance (depending on δ) from the boundary of G.
The jump vector is obtained as a measurable function of W δ(t−). To ensure
the measurability, each point x on ∂G is associated with a nearby point x̄,
chosen in a measurable way from a fixed countable set of points in ∂G. The
jump vector for x is one associated with x̄. We now specify the mapping
x→ x̄ and the associated jump vector more precisely.
By assumption (A5)(ii), for each x ∈ ∂G, there is c(x) ∈RI+ such that
i∈I(x)
ci(x) = 1 and min
j∈I(x)
i∈I(x)
ci(x)γ
i(x), nj(x)
≥ a.(73)
By (73), Lemma 2.1 and the fact that ni(·) is continuous on ∂Gi for each
i ∈ I , we have that for each x ∈ ∂G there is rx ∈ (0, δ) such that for each
y ∈Brx(x)∩ ∂G,
I(y)⊂ I(x)(74)
j∈I(x)
i∈I(x)
ci(x)γ
i(x), nj(y)
.(75)
It follows, using the C1 nature of ∂Gi and the fact that n
i(y) is the inward
unit normal to ∂Gi at y ∈ ∂G for each i ∈ I(y), that (by choosing rx even
30 W. KANG AND R. J. WILLIAMS
smaller if necessary) for each x ∈ ∂G there is m(x)> 0 and rx ∈ (0, δ) such
that for each y ∈Brx(x)∩ ∂G, (74)–(75) hold and
y + λ
i∈I(x)
ci(x)γ
i(x) ∈G for all λ ∈ (0,m(x)).(76)
Let Borx(x) denote the interior of the closed ball Brx(x) for each x ∈ ∂G. The
collection {Borx(x) :x ∈ ∂G} is an open cover of ∂G and it follows that there
is a countable set {xk} such that ∂G⊂
kBrxk (xk) and {xk} ∩BN (0) is a
finite set for each integer N ≥ 1. We can further choose the set {xk} to be
minimal in the sense that for each strict subset C of {xk}, {Brx(x) :x ∈C}
does not cover ∂G. Let Dk = (Brxk (xk) \ (
i=1 Brxi (xi)) ∩ ∂G for each k.
Then Dk 6= ∅ for each k, {Dk} is a partition of ∂G, and for each x ∈ ∂G
there is a unique index i(x) such that x ∈Di(x). For each x ∈R
d, let
x, if x /∈ ∂G,
xi(x), if x ∈ ∂G.
Note that for all x ∈Rd,
‖x− x̄‖< δ.(77)
For each i ∈ I and x∈Rd, let
γi,δ(x) = γi(x̄).(78)
The mapping x → x̄ is Borel measurable on Rd and hence γi,δ is a Borel
measurable function from Rd into Rd.
We construct (W δ, Y δ) as follows. Let X defined on some filtered proba-
bility space (Ω,F ,{Ft}, P ) be a d-dimensional {Ft}-Brownian motion with
drift µ and covariance matrix Γ such that X is continuous surely and X(0)
has distribution ν. Let
τ1 = inf{t≥ 0 :X(t) ∈ ∂G}
W δ(t) =X(t), Y δ(t) = 0 for 0≤ t < τ1.
Note that W δ(τ1−) exists on {τ1 <∞} since X has continuous paths and
in the case that τ1 = 0, W
δ(0−)≡X(0). On {τ1 <∞}, define
Y δi (τ1) =
0, i /∈ I(W δ(τ1−)),
ci(W δ(τ1−))
m(W δ(τ1−))
, i ∈ I(W δ(τ1−)),
W δ(τ1) =X(τ1)
m(W δ(τ1−))
i∈I(W δ(τ1−))
ci(W δ(τ1−))γ
i,δ(W δ(τ1−))
INVARIANCE PRINCIPLE FOR SRBMS 31
So W δ, Y δ have been defined on [0, τ1) and at τ1 on {τ1 <∞}, such that:
(i) W δ(t) = X(t) +
i∈I γ
i,δ(W δ(0−))Y δi (0) +
(0,t]
γi,δ(W δ(s−))dY δi (s) for all t ∈ [0, τ1]∩ [0,∞), where W
δ(0−) =X(0),
(ii) W δ(t) ∈G for t ∈ [0, τ1]∩ [0,∞),
(iii) for i ∈ I ,
(a) Y δi (0)≥ 0,
(b) Y δi is nondecreasing on [0, τ1]∩ [0,∞),
(c) Y δi (t) = Y
i (0) +
(0,t] 1{W δ(s)∈U2δ(∂Gi∩∂G)} dY
i (s) for t ∈ [0, τ1] ∩
[0,∞),
(iv) ‖∆Y δ(t)‖ ≡ ‖Y δ(t)−Y δ(t−)‖ ≤ δ for t ∈ [0, τ1]∩ [0,∞), where Y
δ(0−)≡
Note that (iii)(c) above contains the expression W δ(s) ∈U2δ(∂Gi∩∂G). The
reader may wonder why 2δ appears instead of δ. The reason is that at a jump
time s of Y δi , W
δ(s−) ∈ ∂Gi ∩ ∂G and so
dist(W δ(s), ∂Gi ∩ ∂G)≤ ‖W
δ(s)−W δ(s−)‖+ ‖W δ(s−)−W δ(s−)‖
≤ δ + δ = 2δ.
Proceeding by induction, we assume that for some n ≥ 2, τ1 ≤ · · · ≤ τn−1
have been defined, and W δ, Y δ have been defined on [0, τn−1) and at τn−1
on {τn−1 <∞}, such that (i)–(iv) above hold with τn−1 in place of τ1. Then
we define τn =∞ on {τn−1 =∞}, and on {τn−1 <∞} we define
τn = inf{t≥ τn−1 :W
δ(τn−1) +X(t)−X(τn−1) ∈ ∂G}.
For τn−1 ≤ t < τn, let
Y δ(t) = Y δ(τn−1),
W δ(t) =W δ(τn−1) +X(t)−X(τn−1),
and on {τn <∞}, let
Y δi (τn) =
Y δi (τn−1), i /∈ I(W
δ(τn−)),
Y δi (τn−1) + ci(W
δ(τn−))
m(W δ(τn−))
, i ∈ I(W δ(τn−)),
W δ(τn) =W
δ(τn−)
m(W δ(τn−))
i∈I(W δ(τn−))
ci(W δ(τn−))γ
i,δ(W δ(τn−))
32 W. KANG AND R. J. WILLIAMS
In this way, W δ, Y δ have been defined on [0, τn) and at τn on {τn <∞} such
that (i)–(iv) hold with τn in place of τ1.
By construction {τn}
n=1 is a nondecreasing sequence of stopping times.
Let τ = limn→∞ τn. On {τ =∞}, the construction of (W
δ, Y δ) is complete.
We now show that {τ < ∞} = ∅. In fact, if {τ < ∞} 6= ∅, let ω ∈ {τ <
∞}. The above construction gives (W δ(·, ω), Y δ(·, ω)) on the time interval
[0, τ(ω)). For each t ∈ [0, τ(ω)), we have
W δ(t,ω) =X(t,ω) +
γi,δ(W δ(0−, ω))Y δi (0, ω)
(0,t]
γi,δ(W δ(s−, ω))dY δi (s,ω).
SinceX is continuous on [0,∞), ‖γi,δ(x)‖= 1 for each x ∈Rd and
i∈I Y
i (0,
ω)≤ δ, there are constants λ̃ ∈ (0, τ(ω)) and M̃ > 0 (depending on ω) such
wτ(ω)(X(·, ω) + γ
i,δ(W δ(0−, ω))Y δi (0, ω), λ̃)< ε̃(80)
0≤t≤τ(ω)
∥∥∥∥∥X(·, ω) +
γi,δ(W δ(0−, ω))Y δi (0, ω)
∥∥∥∥∥≤ M̃,(81)
where w
(·, ·) is defined in (3). By the choice of ε̃, δ made at the beginning of
this proof, (77)–(78) and the uniform Lipschitz property of the γi(·), i ∈ I , it
follows that (39) and (43) hold with γi,δ(y) and 2δ in place of γi,n(y,x) and
δn, respectively. Then by similar pathwise analysis to that used in Case 1 and
2 of the proof of Theorem 4.2, with W̃ n =W n =W δ, αn = 0, γi,n(y,x) =
γi,δ(y) for each i ∈ I and x, y ∈ Rd, Xn = X +
i∈I γ
i,δ(W δ(0−))Y δi (0),
Y n = Y δ , Ỹ n = Y δ − Y δ(0), βn = Y δ(0) and δn = 2δ, we obtain that (71)
holds for any T < τ(ω) with ωn = ω, Nη,T = ([τ(ω)/λ̃] + 1)ε+ M̃ . It follows
that supi∈I sups∈[0,τ(ω)) Y
i (s,ω) is finite. By the nondecreasing property of
Y δi (·, ω) on [0, τ(ω)) for each i ∈ I , Y
i (τ(ω)−, ω) exists and is finite for each
i ∈ I . Then by (79) and the continuity of X , we see that W δ(τ(ω)−, ω) exists
and is finite. By the construction of Y δ and the fact that
i∈I(x) ci(x) = 1
for all x∈ ∂G, we have that
Y δi (τ(ω)−, ω) =
m(W δ(τn(ω)−, ω))
∧ δ.(82)
Since τn(ω) ↑ τ(ω) as n → ∞ and W
δ(τ(ω)−, ω) exists, it follows that
{W δ(τn(ω)−, ω)}
n=1 converges to W
δ(τ(ω)−, ω) ∈ ∂G as n → ∞. Conse-
quently, {W δ(τn(ω)−, ω)}
n=1 is a bounded sequence in ∂G and so by the
INVARIANCE PRINCIPLE FOR SRBMS 33
definition of the sets {Dk} which form a partition of ∂G, there is a finite set
C such that
{W δ(τn(ω)−, ω)}
n=1 ⊂
Hence,
m(W δ(τn(ω)−, ω))≤ inf
m(xk)> 0,(83)
and so the right-hand side of (82) is infinite. On the other hand, since
supi∈I sups∈[0,τ(ω)) Y
i (s,ω) is finite, the left-hand side of (82) is finite. This
yields the desired contradiction and so {τ < ∞} = ∅ and we have con-
structed (W δ, Y δ) on [0,∞).
From the construction above, we can see that W δ and Y δ are well-defined
stochastic processes with sample paths in D([0,∞),Rd) and D([0,∞),RI).
They are adapted to the filtration generated by X and satisfy (i)–(iv) above
with [0,∞) in place of [0, τ1].
Consider a sequence of sufficiently small δ’s, denoted by {δn}, such that
δn ↓ 0 as n → ∞. For each δn, let (W δ
, Y δ
) be the pair constructed as
above for the same process X . By the above properties and the fact that for
each i ∈ I and x, y ∈Rd,
‖γi,δ
(y)− γi(x)‖ ≤ ‖γi(ȳ)− γi(x)‖ ≤ L‖ȳ− x‖ ≤ L(δn + ‖y − x‖),
we obtain that Assumption 4.1 holds with W̃ n =W n =W δ
, αn = 0, γi,n(y,x) =
(y) for each i ∈ I and x, y ∈Rd, Xn =X+
i∈I γ
i,δn(W δ
(0−))Y δ
i (0),
Y n = Y δ
, Ỹ n = Y δ
(0), βn = Y δ
(0) and 2δn in place of δn. By invok-
ing the first part of Theorem 4.3, we obtain that {Zδ
}∞n=1 = {(W
δn ,Xδ
)}∞n=1 is C-tight and any weak limit point Z of this sequence satisfies
conditions (i), (ii) and (iv) of Definition 2.1 with Ft = σ{Z(s) : 0 ≤ s ≤ t},
t≥ 0. Note that condition (vi)′ of Theorem 4.3 holds trivially. Furthermore,
= {Xδ
(t)−Xδ
(0)− µt, t≥ 0}= {X(t)−X(0)− µt, t≥ 0} is a mar-
tingale with respect to the filtration generated by X . Since W δ
, Y δ
adapted to this filtration, it follows that M δ
is a martingale with respect
to the filtration generated by W δ
, Y δ
(which in fact is the same as
that generated by X). For each t≥ 0, Xδ
(t)−Xδ
(0) =X(t)−X(0) and
so trivially this forms a uniformly integrable sequence as n varies. It fol-
lows from Proposition 4.1 that condition (vii) of Theorem 4.3 holds. Hence,
any weak limit point of {Zδ
}∞n=1 is an extended SRBM with the data
(G,µ,Γ,{γi, i ∈ I}, ν). �
5.2. SRBMs in convex polyhedrons with constant reflection fields. Exis-
tence and uniqueness in law for SRBMs living in convex polyhedrons with a
constant reflection field on each boundary face has been studied by Dai and
34 W. KANG AND R. J. WILLIAMS
Williams [4]. In this subsection, we state a consequence of our invariance
principle using the results in [4] to establish uniqueness in law. In this case,
G is defined in terms of I (I≥ 1) d-dimensional unit vectors {ni, i ∈ I} and
an I-dimensional vector β = (β1, . . . , βI)
′ such that
G≡ {x ∈Rd : 〈ni, x〉 ≥ βi for all i ∈ I}.(84)
It is assumed that G is nonempty and that the set {(n1, β1), . . . , (n
I, βI)} is
minimal in the sense that no proper subset defines G. For each i ∈ I , let Fi
denote the boundary face: {x ∈G : 〈ni, x〉= βi}. Then, n
i is the inward unit
normal to Fi. A constant vector field γ
i of unit length specifies the direction
of reflection associated with Fi.
Definition 5.2. For each ∅ 6=K⊂ I , define FK =
i∈KFi. Let F∅ =G.
A set K⊂ I is maximal if K 6=∅, FK 6=∅ and FK 6= FK̄ for any K̄ ⊃ K such
that K̄ 6=K.
In [4], Dai and Williams introduced the following assumption.
Assumption 5.1. For each maximal K⊂ I ,
(S.a) there is a positive linear combination n=
i∈K bin
i (bi > 0 ∀i ∈K)
of the {ni, i ∈K} such that 〈n,γi〉> 0 for all i ∈K,
(S.b) there is a positive linear combination γ =
i∈K ciγ
i (ci > 0 ∀i ∈K)
of the {γi, i ∈K} such that 〈ni, γ〉> 0 for all i ∈K.
Remark. For the given G and constant vector fields {γi, i ∈ I}, As-
sumption 5.1 is equivalent to assumption (A5).
Definition 5.3. The convex polyhedron G is simple if for each K⊂ I
such that K 6=∅ and FK 6=∅, exactly |K| distinct faces contain FK.
Remark. The polyhedron G is simple if and only if K is maximal for
every K such that ∅ 6= K ⊂ I and FK 6= ∅. It is shown in [4] that when G
is simple, (S.a) holds for all maximal K⊂ I if and only if (S.b) holds for all
maximal K⊂ I.
Dai and Williams [4] showed that Assumption 5.1 is sufficient for exis-
tence and uniqueness in law of SRBMs living in G with the reflection fields
{γi, i ∈ I} and fixed starting point. [They also showed that condition (S.b)
holding for all maximal K⊂ I is necessary for existence of an SRBM starting
from each point in G. Consequently, when G is simple, Assumption 5.1 is
necessary and sufficient for existence of an SRBM starting from each point
in G.] This yields the following consequence of our invariance principle.
INVARIANCE PRINCIPLE FOR SRBMS 35
Theorem 5.4. Let G be a nonempty domain such that G is a convex
polyhedron of the form (84) (with minimal description), and let {γi, i ∈ I} be
a family of constant vector fields of unit length satisfying Assumption 5.1.
Suppose that Assumption 4.1 and (vi)′, (vii) of Theorem 4.3 hold. Then
W n ⇒W as n→∞ where W is an SRBM associated with (G,µ,Γ,{γi, i ∈
I}, ν).
Proof. Clearly (A1) holds. Assumptions (A2)–(A3) hold by Lemma
A.3. Since for each i ∈ I , γi(·) is a constant vector field of unit length, as-
sumption (A4) holds trivially. Assumption (A5) is implied by Assumption
5.1. Hence by Theorem 4.3, the only thing that we have to check is condition
(viii) of Theorem 4.3, that is, uniqueness in law for SRBMs in convex poly-
hedrons with constant reflection fields of unit length. But this is proved in
Theorem 1.3 of [4] for a fixed starting point in G and follows by a standard
conditioning argument for the initial distribution ν. �
5.3. SRBMs in bounded domains with piecewise smooth boundaries. Dupuis
and Ishii [6] have established sufficient conditions for the existence and path-
wise uniqueness of reflecting diffusions living in the closures of bounded
domains with piecewise smooth boundaries. In this subsection, we state a
consequence of our invariance principle using the results in [6] to establish
uniqueness in law.
Theorem 5.5. Let G be a bounded domain and {γi, i ∈ I} be a family
of reflection fields that satisfy assumptions (A1)–(A4) and (A5)′ in Section
3. We further assume that for each i ∈ I , γi(·) is once continuously differ-
entiable with locally Lipschitz continuous first partial derivatives. Suppose
that Assumption 4.1 and (vi)′, (vii) of Theorem 4.3 hold. Then W n ⇒W
as n→∞ where W is an SRBM associated with (G,µ,Γ,{γi, i ∈ I}, ν).
Remark. We remind the reader that in view of Lemma 3.1, to verify
condition (A5)′, one only needs to show that (i) or (ii) holds for all x ∈ ∂G.
However, as can be seen from the proof below, both forms of the condition
can be useful.
Proof of Theorem 5.5. This theorem follows from Theorem 4.3 and
uniqueness in law for the associated SRBMs. The latter follows by a standard
argument from the pathwise uniqueness established in Corollary 5.2 of [6] for
their Case 2. The conditions required for that case are satisfied in particular
because (A5)′(ii) implies condition (3.8) of [6]. That condition (3.8) readily
implies condition (3.6) of [6]; and, by [5], under the additional smoothness
assumptions imposed on the γi in the statement of our theorem, condition
(3.8) also implies condition (3.7) in [6]. In addition, (A5)′(i) implies that
36 W. KANG AND R. J. WILLIAMS
for each x ∈ ∂G, 〈γi(x), ni(x)〉> 0 for each i ∈ I(x), and furthermore, since
(A5)′ implies (A5), we have by (A5)(i) that the origin does not belong to
the convex hull of the {γi(x) : i ∈ I(x)}. �
APPENDIX: AUXILIARY LEMMAS
Lemma A.1. Suppose that G is bounded. If assumption (A1) holds, then
assumption (A2) holds.
Proof. To see this, suppose G is bounded and assumption (A1) holds.
Fix ε ∈ (0,1). For each i ∈ I and z ∈ ∂Gi ∩ ∂G, by the C
1 property of ∂Gi,
there is a neighborhood Vz of z and a constant R(ε, i, z) > 0 such that for
all x ∈ Vz ∩ ∂Gi ∩ ∂G and y ∈Gi such that ‖x− y‖<R(ε, i, z),
〈ni(x), y − x〉 ≥−ε‖y − x‖.(85)
Assumption (A2) then follows by a standard compactness argument. �
Lemma A.2. Suppose that G is a nonempty bounded domain satisfying
(5), where for each i ∈ I , Gi is a nonempty domain. Then assumption (A3)
holds.
Proof. We prove the lemma by contradiction. Suppose that assumption
(A3) does not hold. Then, since there are only finite many J ⊂ I , J 6=∅,
there is an ε > 0, a nonempty set J ⊂ I , a sequence {rn} ⊂ (0,∞) with rn →
0 as n→∞, a sequence {xn} ⊂R
d such that for each n, xn ∈
j∈J Urn(∂Gj∩
∂G) and dist(xn,
j∈J (∂Gj ∩ ∂G)) > ε. But since G is bounded, {xn} is
bounded and without loss of generality we may assume that xn → x as
n→∞ for some x ∈Rd. It follows that x ∈
j∈J (∂Gj ∩ ∂G), since for each
j ∈ J ,
dist(x,∂Gj ∩ ∂G)≤ ‖xn − x‖+dist(xn, ∂Gj ∩ ∂G)≤ ‖xn − x‖+ rn → 0
as n→∞. This is inconsistent with xn → x and dist(xn,
j∈J (∂Gj ∩∂G))>
Lemma A.3. Suppose (A1) holds where
Gi = {x ∈R
d : 〈ni, x〉> βi} for i ∈ I,(86)
{ni, i ∈ I} is a finite collection of d-dimensional vectors of unit length, and
for I= |I|, β = (β1, . . . , βI)
′ is an I-dimensional vector. (Thus, G is a convex
polyhedron.) Assume that for each i ∈ I , ∂Gi ∩ ∂G 6=∅. Then assumptions
(A2) and (A3) hold.
INVARIANCE PRINCIPLE FOR SRBMS 37
Proof. Assumption (A2) holds automatically since G is convex. In or-
der to show that assumption (A3) holds, we just need to show that for each
J ⊂ I with J 6=∅,
(∂Gj ∩ ∂G)
Ur(∂Gj ∩ ∂G)
→ 0(87)
as r→ 0. Fix J ⊂ I such that J 6=∅. Then
j∈J (∂Gj ∩∂G) is the collection
of all solutions x ∈Rd to the following system of linear inequalities:
〈ni, x〉 ≥ βi for all i ∈ I,
〈−ni, x〉 ≥ −βi for all i ∈ J .
Suppose that
j∈J (∂Gj ∩∂G) 6=∅, that is, (LS) has at least one solution.
By a theorem of Hoffman [11], with supporting lemmas proved by Agmon
[1], there is a constant C > 0 (depending only on {ni, i ∈ I} and not on β)
such that for any x ∈Rd there exists a solution x0 ∈R
d of (LS) with
‖x− x0‖ ≤C
(βi − 〈n
i, x〉)+ +
(−βi − 〈−n
i, x〉)+
.(88)
For r > 0, any x ∈
j∈J Ur(∂Gj ∩ ∂G) satisfies the following:
〈ni, x〉 ≥ βi − r for all i ∈ I,
(r-LS)
〈−ni, x〉 ≥ −βi − r for all i ∈ J .
Then by (88), there is x0 ∈
j∈J (∂Gj ∩ ∂G) such that
(∂Gj ∩ ∂G)
≤ ‖x− x0‖ ≤ 2C|I|r.
It follows that (87) holds when
j∈J (∂Gj ∩ ∂G) 6=∅.
Now suppose that
j∈J (∂Gj ∩∂G) =∅, that is, (LS) has no solution. We
shall use an argument by contradiction to show that
j∈J Ur(∂Gj ∩∂G) =∅
for all r sufficiently small. Suppose that this is not true. Then we have that⋂
j∈J Ur(∂Gj ∩ ∂G) 6= ∅ for all r ∈ (0,∞). As we have seen before, any
j∈J Ur(∂Gj ∩ ∂G) is a solution to (r-LS). We now construct a Cauchy
sequence. Let x1 ∈
j∈J U1/2(∂Gj ∩ ∂G). Then x1 is a solution to (
-LS).
Since ( 1
-LS) has at least one solution, by the theorem of Hoffman [11] (using
the fact that the constant C depends only on {ni, i ∈ I}), we conclude that
there is a solution x2 to (
-LS) such that ‖x1−x2‖ ≤
, where C ′ = 2C|I|.
Continuing in this manner, we can obtain a sequence {xn}
n=1 such that for
each n ≥ 1, ‖xn − xn+1‖ ≤
and xn+1 is a solution of (
-LS). The
38 W. KANG AND R. J. WILLIAMS
sequence {xn}
n=1 is Cauchy. Hence, there is an x
∗ ∈Rd such that xn → x
as n→ ∞, and x∗ is a solution to (LS). This contradicts the supposition
j∈J (∂Gj ∩ ∂G) =∅. Thus we have that
j∈J Ur(∂Gj ∩ ∂G) =∅ for
all r sufficiently small, and for such r,
(∂Gj ∩ ∂G)
Ur(∂Gj ∩ ∂G)
by convention.
Combining the above we see that for each J ⊂ I with J 6=∅, (87) holds
and hence assumption (A3) holds. �
Remark. In fact, under the assumptions of Lemma A.3, there is a con-
stant C > 0 such that D(u) ≤ Cu for each u ≥ 0 and D(·) defined as in
assumption (A3).
Lemma A.4. Given T > 0, functions φ,{φn}∞n=1 in D([0,∞),R
d), and
χ,{χn}∞n=1 in D([0,∞),R), suppose that sup0≤s≤T ‖φ
n(s)− φ(s)‖ → 0 and
sup0≤s≤T |χ
n(s)−χ(s)| → 0 as n→∞. Assume that χn is nondecreasing for
each n. Then for any sequence of real valued continuous functions {fn}∞n=1
defined on Rd such that fn converges uniformly on each compact set to a
continuous function f :Rd →R, we have
(0,t]
fn(φn(s))dχn(s)→
(0,t]
f(φ(s))dχ(s) as n→∞,(89)
uniformly for t ∈ [0, T ].
Proof. By replacing χn(·) and χ(·) by χn(·)− χn(0) and χ(·)− χ(0),
respectively, we may assume that χn(0) = χ(0) = 0. It is straightforward to
see by the uniform convergence of {χn} to χ on [0, T ] that χ inherits the
nondecreasing property of the {χn}.
By the triangle inequality,
0≤t≤T
(0,t]
fn(φn(s))dχn(s)−
(0,t]
f(φ(s))dχ(s)
≤ sup
0≤t≤T
(0,t]
(fn(φn(s))− f(φ(s)))dχn(s)
∣∣∣∣(90)
+ sup
0≤t≤T
(0,t]
f(φ(s))d(χn(s)− χ(s))
∣∣∣∣.
For the first term on the right-hand side of the above inequality, we have
0≤t≤T
(0,t]
(fn(φn(s))− f(φ(s)))dχn(s)
≤ sup
0≤s≤T
|fn(φn(s))− f(φ(s))|χn(T ),
INVARIANCE PRINCIPLE FOR SRBMS 39
where the right-hand side member above tends to zero as n → ∞ by the
uniform convergence of φn to φ on [0, T ] (which implies uniform boundedness
of {φn} on [0, T ]), the uniform convergence of fn to f on compact sets, the
continuity of f , and the convergence of χn(T ) to χ(T ). For the second term,
note that since f(φ(·)) ∈D([0,∞),R), by Theorem 3.5.6, Proposition 3.5.3
and Remark 3.5.4 of [7], there is a sequence of step functions {zk}∞k=1 of the
zk(·) =
zk(tki )1[tk
)(·),(91)
where 1 ≤ lk < ∞, 0 = t
1 < t
2 < · · · < t
< ∞ and sup0≤s≤T |f(φ(s)) −
zk(s)| → 0 as k→∞. Then
0≤t≤T
(0,t]
f(φ(s))d(χn(s)− χ(s))
≤ sup
0≤t≤T
(0,t]
(f(φ(s))− zk(s))d(χn(s)− χ(s))
+ sup
0≤t≤T
(0,t]
zk(s)d(χn(s)− χ(s))
≤ sup
0≤s≤T
|f(φ(s))− zk(s)|(χn(T ) + χ(T ))
+ sup
0≤t≤T
|zk(tki )||(χ
n − χ)((tki+1 ∧ t)−)− (χ
n − χ)((tki ∧ t)−)|.
For fixed k, the last term above can be made as small as we like for all n
sufficiently large since χn → χ uniformly on [0, T ]. The desired result follows.
Remark. The proof of Lemma A.4 is a modification of the proof of the
related Lemma 2.4 in [4]. The difference in assumptions is that in [4] it is
assumed that φn → φ in the J1-topology rather than uniformly on [0, T ],
χn, χ ∈C([0,∞),R+) rather than χ
n, χ ∈D([0,∞),R), and there is a single
function f rather than a sequence {fn}.
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Department of Mathematical Sciences
Carnegie Mellon University
Pittsburgh, Pennsylvania 15213
E-mail: [email protected]
Department of Mathematics
University of California at San Diego
9500 Gilman Drive
La Jolla, California 92093
E-mail: [email protected]
http://www.ams.org/mathscinet-getitem?mr=1346729
http://www.ams.org/mathscinet-getitem?mr=2222685
http://www.ams.org/mathscinet-getitem?mr=1096165
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http://www.ams.org/mathscinet-getitem?mr=0959133
http://www.ams.org/mathscinet-getitem?mr=1231926
http://www.ams.org/mathscinet-getitem?mr=1663755
http://www.ams.org/mathscinet-getitem?mr=1663759
mailto:[email protected]
mailto:[email protected]
Introduction
Notation, terminology and preliminaries
Definition of an SRBM
Assumptions on the domain G and the reflection vector fields {i}
Assumptions on the domain G
Assumptions on the reflection vector fields {i}
Invariance principle
Oscillation inequality
C-tightness result
Invariance principle for SRBMs
Applications of the invariance principle
Weak existence of SRBMs
SRBMs in convex polyhedrons with constant reflection fields
SRBMs in bounded domains with piecewise smooth boundaries
Appendix: Auxiliary lemmas
References
Author's addresses
|
0704.0406 | Finite Drude weight for 1D low temperature conductors | Finite Drude weight for 1D low temperature conductors
Dariush Heidarian and Sandro Sorella
Istituto Nazionale di Fisica della Materia (INFM)-Democritos, National Simulation Centre,
and Scuola Internazionale Superiore di Studi Avanzati (SISSA), I-34014 Trieste, Italy
We apply well established finite temperature QuantumMonte Carlo techniques to one dimensional
Bose systems with soft and hardcore constraint, as well as to spinless fermion systems. We give
clear and robust numerical evidence that, as expected, no superfluid density for Bosons or Meissner
fraction for fermions. is possible at any non zero temperature in one dimensional interacting Bose
or fermi lattice models, whereas a finite Drude weight is generally observed in gapless systems, in
partial disagreement to previous expectations.
PACS numbers: 74.25.Fy,71.27.+a,71.10.Fd
I. INTRODUCTION
In the last decades there have been a lot of numerical
and theoretical works to understand the role of strong
correlation in lattice model Hamiltonians.1,2,3,4,5,6,7 Re-
cently this issue has acquired an increasing attention and
remarkable importance, due to the recent advances in the
realization of optical lattices. In these experiments ultra-
cold atoms behave as boson particles trapped on particu-
lar lattice sites, whereas the interaction and the hopping
parameters can be tuned continuously. This important
achievement has opened the possibility to verify directly
the crucial role played by the electron correlation in very
important model Hamiltonians defined on a lattice. An
important example is the realization of a Mott insulating
state in a system with strong on site repulsion8,9. More-
over quite recently the possibility to include the Fermi
statistics in optical lattices appears very promising and
interesting.10
In 1D spinless fermion systems are equivalent to inter-
acting Bose systems with hard-core constraint and are
described by the same low energy theory -the Luttinger
liquid theory-. Indeed this theory holds also for soft-
core bosons, as shown in Ref.(7). Therefore, as far as
the transport properties are concerned one should expect
the same behavior both for fermions and bosons. On the
other hand for lattice models, even in absence of disor-
der, the current does not commute with the Hamiltonian,
implying its possible decay at finite temperature due to
the backscattering processes11. In this case the dynam-
ical current-current correlation function also decays in
time, leading to a current Fourier transform without δ
function at zero energy, namely without a finite Drude
weight within the linear response theory.
Until few decades ago the absence of the Drude weight
was the expected behavior of all interacting metals in lat-
tice models or in real solids at finite temperature. How-
ever a quite clear numerical evidence has been reported
in Ref.12 that current should not decay in integrable 1D
models, namely for Hamiltonians that can be solved by
Bethe ansatz techniques in 1D. These models essentially
possess some hidden conservation law, that was conjec-
tured to forbid the current decay process.12,13 Later sev-
eral groups have reproduced this surprising effect14,15,
with a noticeable exception that a finite Drude weight
at finite temperature was found also for non-integrable
models.15 On the other hand, from purely theoretical
grounds this issue is not settled yet: in Ref.11 it was ar-
gued that backscattering processes can be effective also
at finite temperature and in 1D non integrable models,
whereas in Ref.16, it was proposed that also some par-
ticular non integrable model could provide a conserved
current.
In this work we propose that the general behavior of
1D gapless systems is eventually characterized by a finite
Drude weight at finite temperature, and we have found
no exception in the models that we have studied. This
conclusion is based on a careful and systematic numerical
work on fairly generic one dimensional Bose and Fermi
systems, that all show the same behavior, even though
strong finite size effects are observed in the non integrable
cases.
In the following we investigate the behavior of the
Drude weight in 1D systems in the thermodynamic limit
and finite temperature.
Model and Method : We have studied hardcore and
softcore bosons in a 1D lattice with periodic boundary
conditions. The Hamiltonian studied reads,
iai+1 + h.c.) +
ni(ni − 1)
+V nini+1 +Wnini+2 − µni
The sum is over all lattice sites i, a
i/ai is the boson
creation/annihilation operator at site i, henceforth ni is
the particle number at site i andµ is the chemical poten-
tial. t is the hopping amplitude which is set to one, U is
the on-site repulsion, whereas V and W are the nearest
and the next-nearest neighbor interactions, respectively.
For hardcore bosons in the U → ∞ limit the Hamilto-
nian can be mapped onto an S = 1/2 spin system with
Szi = ni − 1/2 and S
i = a
i . In this work we present
our results for the half filled case of hardcore and soft-
core models. Most of our results have been obtained by
Quantum Monte Carlo (QMC), using the stochastic se-
ries expansion (SSE)5,17 with the directed loop update18.
http://arxiv.org/abs/0704.0406v1
Superfluid density ρs (or spin stiffness in the equiva-
lent spin model), is defined as the second derivative of
the free energy with respect to a twist in the boundary
conditions. In order to compute this quantity by QMC,
it is convenient to apply linear response theory, relating
this quantity to the current current response function
Λ(q, iωn) =
dτ exp(iωnτ)〈J(q, τ)J(−q, 0)〉/N , where
J is the current operator and ωn is Matsubara frequency.
Then the following expression for the superfluid density
is obtained:
ρs = 〈−K〉 − Λ(q = 0; iωn = 0) =
〈W 2〉
where 〈K〉 is the average kinetic energy per site, ωn =
2πn/β are the Matsubara frequencies andW is the wind-
ing number. Similarly the Drude weight is obtained with
the same expression but with a different order in the limit
ω → 0 and q → 0, namely15,19,20
D = 〈−K〉 − ReΛ(q = 0, ω → 0). (3)
In SSE one can obtain Λ very accurately in terms of
Matsubara frequencies. Therefore analytic continuation
of the data is required. In order to avoid difficulties of
extrapolation to iωn → 0 at large temperatures, we have
worked at relatively low temperatures (β ≥ 10).
In principle, due to the different order of lim-
its, the Drude weight and the superfluid density may
be different when the following quantity remains fi-
nite in the thermodynamic limit15: D − ρs =
En=Em
β exp(−βEn)|〈ψn|J |ψm〉|
2/L, where, J is the
current operator, while En and |ψn〉 are the n
th eigen-
value and eigenstate of the many body system, respec-
tively.
The current operator can be written as J(q = 0) =
) where H+
al+1 and b is the bond
index, corresponding to the site index l. The ensemble
average of product of two local operatorsHσ1
andHσ2
(τ)Hσ1
(0)〉 =
n,m=0
(τ − β)n(−τm)
〈ψk|H
HmHσ1
|ψk〉 (4)
where τ is the imaginary time, Z is partition function
and the summation over n and m comes from Taylor-
expansion of e(−β+τ)H and e−τH . Following Ref.17 the
relation (4) can be simplified to
(β − τ)ns−m−2τm
(ns − 1)!
(ns −m− 2)!m!
N b1b2,σ1σ2m
where ns is the length of sequence of the local operators
and it changes in each QMC sampling. N b1b2,σ1σ2m is the
number of times that two operatorsHσ1
and Hσ2
appear
in this sequence with distance of m local operators, and
〈...〉W indicates an arithmetic average using configura-
tions with relative weight W . In this work we introduce
an efficient way to sample by SSE the current-current
response function. To this end, we multiply expression
(5) by eiωnτ and integrate over the imaginary time τ , we
obtain:
1F1(m+ 1, ns; 2iπn)N
b1b2σ1σ2
where
1F1(m+ 1, n; z) =
(n− 1)!
(n−m− 2)!m!
dx exp(zx)xm(1− x)n−m−2 (7)
is the confluent hypergeometric function.
Therefore, the current-current correlation acquires
contributions determined by length of operator string
ns. All these contributions are stochastically sampled
in an efficient way, and in each statistical measurement
the correlation function Λ(q = 0, iωn) has the following
estimator:
σ1,σ2=±
1F 1(m+ 1, ns; 2iπn)N
m (8)
where Nσ1σ2m =
b1,b2
N b1b2,σ1σ2m .
Discussion: At zero temperature, for non degenerate
ground state, the Drude weight and the superfluidity are
the same. In a 1D system at any finite temperature ρs is
expected to be zero in the thermodynamic limit, whereas
the Drude weight can be non-zero. For hardcore and soft-
core bosons in a 1D lattice, a systematic size scaling of
the superfluid density ρs clearly shows that this quantity
vanishes in the thermodynamic limit and for any finite
temperature (see figures 1 and 2). Further, we find that,
for a fixed set of parameters and at half filling, all super-
fluidity data versus 1/L collapse to one curve whenever
the x-axis is appropriately scaled with the temperature
T (see figures 1 and 2). This analysis suggests the scal-
ing form ρs(β, L) ≡ ρs(β/L). If one takes the order of
limit T → 0 after L → ∞, superfluidity remains zero
even at zero temperature. Notice that by taking first the
limit T → 0 and then L → ∞ superfluidity has a finite
value for the gapless phase, but this is not a signature of
superfluidity, rather the occurrence of a finite zero tem-
perature Drude weight. Though in 1D is not possible to
have a finite superfluid density at any non zero tempera-
ture, several authors have identified the finite zero tem-
perature Drude weight with the superfluid density for a
superfluid with vanishing critical temperature. We be-
lieve that this identification is a bit confusing and there-
fore we prefer to think about absence of superfluidity and
superconductivity in 1D systems, as commonly reported
in the textbooks.
Fig. 3 shows the current-current correlation versus ωn
in the metallic and insulating phases of an integrable
FIG. 1: (color online) Superfluid stiffness for an integrable (a)
and a non-integrable (b) model versus β/L. The system size
L is ranging from 50 to 1200.
FIG. 2: (color online) Superfluidity of the soft-core bosons
versus scaled system size at half filling, the on-site interaction
is U = 4
model (W = 0, U = ∞). The zero-frequency value is
the superfluid density ρs and the limit ωn → 0 gives the
Drude weight D. For W = 0 at zero temperature, there
exists a critical value Vc/t = 2 below which the Drude
weight is finite. In the first case (a) shown in Fig.(3)
with V/t = 2 the Drude weight has a finite value at any
finite temperature, which is consistent with the previous
works12. In the insulating phase (case b) with V/t = 3,
the superfluid density coincides with the Drude weight
and they both tend to zero as the system size increases.
In a non-integrable model such as hard-core bosons
with nearest and next nearest neighborer interactions
earlier works have suggested zero Drude weight as system
size increases. With SSE we can go to very large system
sizes and low temperatures and check the scaling depen-
dence of the Drude weight. In Fig. 4 we have plotted
current-current correlation versus Matsubara frequency
for different L, and a fixed temperature T = 1/100. As
shown in the same Figure (4) we have also found a finite
Drude weight at finite T in the celebrated Bose-Hubbard
model with softcore constraint and in several other mod-
FIG. 3: (color online) (a) Current-current correlation for an
integrable model in the metallic phase. The zero frequency
data shows superfluidity while the extrapolation to n → 0 is
the Drude weight. D remains finite with increasing L while
ρs vanishes. (b) In the insulating phase D and ρs have the
same value and both tend to zero by increasing L.
FIG. 4: (color online) Response function vs. n for (a) hard-
core bosons with V/t = 1.5, W/t = 1, T/t = 1/100 and (b)
Bose-Hubbard model with softcore constraint and U/t = 2,
µ/t = −0.4, T/t = 1/25. The system sizes ranges from
L = 100 to L = 800.
els (not shown). Although some evidence that few par-
ticular non integrable models could have a finite Drude
weight at finite temperature have been reported before,
here we have found a very convincing evidence that this
behavior should be generic for 1D gapless system regard-
less from their integrability. We have supported this
statement by state of the art numerical calculations ob-
tained for very large system sizes and low temperature so
that all possible extrapolations are perfectly under con-
trol.
In conclusion it turns out that, at low energy, all gap-
less lattice models studied scale to the Luttinger liquid
fixed point where the backscattering is a marginally irrel-
evant coupling and the current is therefore conserved at
the fixed point. This is therefore a peculiar and generic
feature of 1D. Indeed in 2D systems, such as hardcore
bosons with n.n. repulsion in a square and triangular
lattice, we found no difference between ρs and D.
Acknowledgments
We thank M. Troyer for useful discussions. This work
is partially supported by COFIN-2005 and CNR.
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|
0704.0407 | Density dependent hadronic models and the relation between neutron stars
and neutron skin thickness | arXiv:0704.0407v1 [nucl-th] 3 Apr 2007
Density dependent hadronic models and the relation between neutron stars and
neutron skin thickness
S.S. Avancini,1 J.R. Marinelli,1 D.P. Menezes,1 M.M.W. Moraes,1 and C. Providência2
Depto de F́ısica - CFM - Universidade Federal de Santa Catarina Florianópolis - SC - CP. 476 - CEP 88.040 - 900 - Brazil
Centro de F́ısica Teórica - Dep. de F́ısica - Universidade de Coimbra - P-3004 - 516 - Coimbra - Portugal
In the present work we investigate the main differences in the lead neutron skin thickness, binding
energy, surface energy and density profiles obtained with two different density dependent hadron
models. Our results are calculated within the Thomas-Fermi approximation with two different
numerical prescriptions and compared with results obtained with a common parametrization of the
non-linear Walecka model. The neutron skin thickness is a reflex of the equation of state properties.
Hence, a direct correlation between the neutron skin thickness and the slope of the symmetry
energy is found. We show that within the present approximations the asymmetry parameter for low
momentum transfer polarized electron scattering is not sensitive to the model differences.
PACS number(s): 21.65.+f,24.10.Jv,95.30.Tg,26.60.+c
I. INTRODUCTION
The relation between neutron star properties which are
obtained from adequate equations of state (EoS) and the
neutron skin thickness has long been a topic of investi-
gation in the literature. The details of this relation and
the important quantities to be discussed have been well
established in [1], where it was shown that the difference
between the neutron and the proton radii, the neutron
skin thickness, is linearly correlated with the pressure of
neutron matter at sub-nuclear densities. This is so be-
cause the properties of neutron stars are obtained from
appropriate EoS whose symmetry energy depends on the
density and also controls the size of the neutron skin
thickness in heavy and asymmetric nuclei, as 208 Pb, for
instance. It is important to remember that the EoS in
neutron stars is also very isospin asymmetric due to the
β- equilibrium constraint.
Hence, isospin asymmetry plays a major role in the un-
derstanding of the density dependence of the symmetry
energy and the consequences it may arise [2]. In [3, 4]
it was shown that the models that yield smaller neutron
skins in heavy nuclei tend to yield smaller neutron star
radii due to a softer EoS.
Neutron stars are believed to have a solid crust formed
by nonuniform neutron rich matter in β-equilibrium
above a liquid mantle. In the inner crust nuclei coex-
ist with a gas of neutrons which have dripped out. The
properties of this crust as, for instance, its thickness and
pressure at the crust-core interface depend a lot on the
density dependence of the EoS used to describe it [4, 5].
On the other hand, it is well known [6, 7] that the ex-
istence of phase transitions from liquid to gas phases in
asymmetric nuclear matter (ANM) is intrinsically related
with the instability regions which are limited by the spin-
odals. Instabilities in ANM described within relativistic
mean field hadron models, both with constant and den-
sity dependent couplings at zero and finite temperatures
have already been investigated [7] and it was shown that
the main differences occur at finite temperature and large
isospin asymmetry close to the boundary of the insta-
bility regions. In neutral neutron-proton-electron (npe)
matter the electrons are also included. In a thermody-
namical calculation the instabilities almost completely
disappear due to the high electron Fermi energy [8].
However, in a dynamical calculation which includes the
Coulomb interaction and allows for independent neutron,
proton and electron fluctuations [9, 10], it is seen that
the electron dynamics tends to restore the short wave-
length instabilities although moderated by the high elec-
tron Fermi energy.
Moreover, it is also known that the liquid-gas phase
transition in ANM can lead to an isospin distillation phe-
nomenon, characterized by a larger proton fraction in
the liquid phase than in the gas phase. This is due to
the repulsive isovector channel of the nuclear interaction
[11–13].
In a recent work the spinodal section and related quan-
tities, as the neutron to proton density fluctuations re-
sponsible for the distillation effect, has been studied
within different relativistic models [8]. It was shown that
the distillation effect within density dependent relativis-
tic models decreases with density above a nuclear density
of ∼ 0.02−0.03 fm−3, a result similar to the one obtained
with the SLy230a parametrization of Skyrme interaction
[14] and contrary to the results found with the more com-
mon relativistic parametrizations with no density depen-
dent coupling parameters. In the last case the distillation
effect becomes always larger as the density increases.
Also, the behavior of the symmetry energy obtained
with density dependent models is closer to what one ob-
tains with non-relativistic models than with other rel-
ativistic models with constant couplings [7]. In an at-
tempt to understand this behavior, a comparison be-
tween the non-relativistic Skyrme effective force and rel-
ativistic mean field models at subsaturation densities was
performed [15]. It was shown that the relativistic mod-
els could also be reduced to an energy density functional
similar to the one describing the Skyrme interaction.
http://arxiv.org/abs/0704.0407v1
There have already been some efforts in order to com-
pare nuclear matter and finite nuclei properties obtained
both with relativistic and non-relativistic models [16, 17]
but there is no clear or obvious explanations for the differ-
ences. At very low densities both, the relativistic and the
non-relativistic approaches predict a non-homogeneous
phase commonly named pasta phase, formed by a com-
petition between the long-range Coulomb repulsion and
the short-range nuclear attraction [18].
Based on the above arguments, it is very important
that an accurate experimental measurement of the neu-
tron skin thickness is achieved. This depends on a precise
measurement of both the charge and the neutron radius.
The charge radius is already known within a precision of
one percent for most stable nuclei, using the well-known
single-arm and non-polarized elastic electron scattering
technique as well as the spectroscopy of muonic atoms
[19] . For the neutron radius, our present knowledge has
an uncertainty of about 0.2 fm [20]. However, using po-
larized electron beams it is possible to obtain the neutron
distribution in nuclei in a fairly model independent way,
as first discussed in [21] and, as a consequence, to obtain
the desired neutron radius. In fact, the Parity Radius
Experiment (PREX) at the Jefferson Laboratory [22] is
currently running to measure the 208Pb neutron radius
with an accuracy of less than 0.05 fm, using polarized
electron scattering.
In the present work, we use two different hadronic
models that incorporate density dependence in differ-
ent ways. The first one, to which we refer next as the
TW model is a density dependent hadronic model with
the meson-to-nucleon couplings explicitly dependent of
the density [23, 24]. In the following it is used to cal-
culate the neutron skin thickness of 208Pb, which is a
neutron-rich heavy nucleus. This model was chosen be-
cause it is based on a microscopic calculation, fits well
many nuclei properties and, as stated above, has shown
to provide results which are different from the usual NL3
[25] and TM1 [26] parametrizations for the non-linear
Walecka model (NLWM), having a richer density depen-
dence of the symmetry energy than most of the rela-
tivistic nuclear models. The original motivation for the
development of this density dependent hadronic model
[27, 28] was to reproduce results obtained with the rel-
ativistic Dirac-Brueckner Hartree-Fock (DBHF) theory
[29]. Later the DBHF calculations for nuclear matter
were taken only as a guide for a suitable parametrization
of the density dependence of the meson-nucleon coupling
operators [24, 30]. Moreover, density dependent hadronic
models can also be a useful tool in obtaining EoS for
neutron stars even if hyperons are to be considered [32],
which is not the case if NL3 or TM1 are used. Both, NL3
and TM1, can only be used if the EoS is restricted to ac-
commodate neutrons, protons and the leptons necessary
to enforce β-stability. Once hyperons are included, the
nucleons acquire a negative effective mass above∼ 3−4ρ0
densities [33, 34], where ρ0 is the nuclear saturation den-
sity.
The second model, that we refer to as NLωρ model,
includes non-linear σ − ρ and ω − ρ couplings [3, 35–
37] which allow to change the density dependence of the
symmetry energy of the most common parametrizations
of the NLWM that show essentially a linear behavior of
the symmetry energy with density. However, the symme-
try energy determines the behavior of isospin asymmetric
matter and therefore is intrinsically related to the char-
acteristics of the EoS that can describe neutron stars.
Within this model the authors of [3] have shown that
the neutron skin thickness of 208Pb was sensitive to the
isovector channel of the nuclear interaction and there was
a correlation between neutron skin thickness of nuclei and
properties of neutron stars.
For the sake of completeness, the results of the present
work, whenever possible are compared with the results
obtained with the NL3 parametrization of the NLWM,
known to describe finite nuclei properties well.
We perform two different numerical calculations to ob-
tain the 208Pb properties: a Thomas-Fermi approxima-
tion based on the liquid-gas phase transition developed
in [38] and a Thomas-Fermi approximation based on a
method proposed in [39], where a harmonic oscillator ba-
sis is used. We restrict ourselves to the Thomas-Fermi
approximation because, as we show in the Results sec-
tion at the end of the paper, for the purpose of obtaining
correct surface energy and neutron-skin thickness, it is
almost as good as the solution of the Dirac equation.
At this point it is worth mentioning that the scalar-
isovector δ mesons, which play an important role in the
isospin channel, could also be incorporated in our work
as done in [7, 9, 40] but in order to make the compar-
isons among different approximations as simple as possi-
ble, they will be included in a future work. Finally, as
we are interested in nuclei ground state properties, all
calculations are performed at zero temperature.
II. THE TW DENSITY DEPENDENT
HADRONIC MODEL
Next we describe the main quantities of the TWmodel,
which has density dependent coupling parameters. The
Lagrangian density reads:
L = ψ̄
i∂µ − ΓvV µ −
τ · bµ
−e (1 + τi3)
− (M − Γsφ)
(∂µφ∂
µφ−m2sφ2)−
m2vVµV
Bµν ·Bµν +
m2ρbµ ·bµ−
µν (1)
where φ, V µ, bµ and Aµ are the scalar-isoscalar, vector-
isoscalar and vector-isovector meson fields and the pho-
ton field respectively, Ωµν = ∂µVν − ∂νVµ , Bµν =
∂µbν − ∂νbµ − Γρ(bµ × bν), Fµν = ∂µAν − ∂νAµ and
τp3 = 1, and τn3 = −1. The parameters of the model
are: the nucleon mass M = 939 MeV, the masses of the
mesons ms, mv, mρ, the electromagnetic coupling con-
stant e =
4π/137 and the density dependent coupling
constants Γs, Γv and Γρ, which are adjusted in order
to reproduce some of the nuclear matter bulk properties
shown in Table I, using the following parametrization:
Γi(ρ) = Γi(ρsat)hi(x), x = ρ/ρsat, (2)
hi(x) = ai
1 + bi(x+ di)
1 + ci(x + di)2
, i = s, v (3)
hρ(x) = exp[−aρ(x− 1)], (4)
with the values of the parameters mi, Γi(ρsat), ai, bi, ci
and di, i = s, v, ρ given in [24]. This model does not
include self-interaction terms for the meson fields (i.e.
κ = 0, λ = 0 and ξ = 0 ) as in NL3 or TM1 parametriza-
tions for the NLWM.
The field equations of motion follow from the Euler-
Lagrange equations. When they are obtained, some care
has to be taken since the coupling operators depend on
the baryon fields ψ̄ and ψ through the density. When the
partial derivatives of L are performed relatively to the
fields ψ̄ and ψ, they yield extra terms due to the func-
tional dependence of the coupling operators. The new
terms are absent in the usual Quantum Hadrodynamic
(QHD, NLWM) models [25, 26, 31]. The equations of
motion for the fields read:
µ +m2φ)φ = Γsψ̄ψ, (5)
µν +m2vV
µ = Γvψ̄γ
µψ, (6)
µν +m2ρb
ψ̄τγµψ, (7)
ψ̄(1 + τ3)γ
µψ, (8)
[γµ(i∂
µ − Σµ)−M∗]ψ = 0 , (9)
whereM∗ =M−Γsφ. Notice that in the equation of mo-
tion for the baryon field ψ the vector self-energy consists
of two terms, Σµ = Σ
µ + Σ
µ , where:
Σ(0)µ = ΓωVµ +
τ · bµ +
(1 + τ3)Aµ, (10)
ΣRµ =
V νjν +
bν · jν3 −
where Σ
µ is the usual vector self-energy, ρ̂uµ = jµ with
u2 = 1 jν = ψ̄γνψ, j
3 = ψ̄τγ
νψ and, as a result of
the derivative of the Lagrangian with respect to ρ a new
term appears, ΣRµ , which is called rearrangement self-
energy and has been shown to play an essential rôle in
the applications of the theory. This term guarantees the
thermodynamical consistency and the energy-momentum
conservation. For more detailed calculations, at zero and
finite temperatures, please refer to [41].
In the static case there are no currents in the nucleus
and the spatial vector components are zero. Therefore,
the mesonic equations of motion become:
∇2φ = m2sφ− Γsρs, (12)
∇2V0 = m2vV0 − Γvρ, (13)
∇2b0 = m2ρb0 −
ρ3, (14)
∇2A0 = −eρp, (15)
where ρs =< ψ̄ψ > is the scalar density, ρ = ρp + ρn,
ρ3 = ρp − ρn and ρp and ρn are the proton and neutron
densities.
A. Thomas-Fermi approximation
We first define the functional
Ω = E − µpBp − µnBn, (16)
where E is the energy, µp (µn) is the proton (neutron)
chemical potential and Bp (Bn) is the proton (neutron)
number. Within the semi-classical Thomas-Fermi ap-
proximation, the energy of the nuclear system with par-
ticles described by the one-body phase-space distribution
function f(r,p, t) at position r, instant t with momentum
p is given by
(2π)3
fi(r,p, t)
p2 +M∗2 + Vi
(∇φ)2 +m2sφ2 − (∇V0)2 −m2vV 20
−(∇b0)2 −m2ρb20 − (∇A0)2
where
Vp = ΓvV0 +
b0 + eA0 , Vn = ΓvV0 −
γ = 2 refers to the spin multiplicity and the distribu-
tion functions for protons and neutrons are
fi = θ(k
Fi(r) − p2), i = p, n .
In this approach, the scalar, proton and neutron densities
become:
ρs(r) =
i=p,n
∫ kFi(r)
with ǫ =
p2 +M∗2 and
d3rρi, ρi(r) =
k3Fi(r).
From the above expressions we get for (16)
(∇φ)2 − (∇V0)2 − (∇b0)2 − (∇A0)2
+ Vef
Vef =
2 −m2vV 20 −m2ρb20
− µpρp − µnρn
i=p,n
∫ kFi
dpp2ǫ+ΓvV0ρ+Γρ
ρ3 + eA0ρp (18)
Minimization of Ω with respect to kFi(r), i = p, n,
gives rise to the following conditions
k2Fp +M
∗2 − ΓvV0 −
b0 − eA0 − ΣR0
k2Fn +M
∗2 − ΓvV0 +
b0 − ΣR0
where the rearrangement term is
ΣR0 =
ρ V0 +
− ∂ Γs
ρs φ.
From the above equations we obtain kFp = 0 and kFn = 0
or, for kFp or kFn different from zero,
k2Fp +M
∗2 + ΓvV0 +
b0 + eA0 +Σ
0 , (19)
k2Fn +M
∗2 + ΓvV0 −
b0 +Σ
0 . (20)
The values of kFp and kFn are obtained inverting these
two last equations.
Such density dependences in the coupling parameters
do not affect the energy functional but of course affect
its derivative such as the pressure density and the chem-
ical potentials. As already discussed in the literature
[7–9, 32], the rearrangement term is crucial in obtaining
different behaviors in physical properties related to the
chemical potentials or to their derivatives with respect to
the density, such as spinodal regions, as compared with
the more common NL3 or TM1 parametrizations.
III. NLωρ MODEL
The Lagrangian density that incorporates the extra
non-linear σ − ρ and ω − ρ couplings [3, 35–37] reads
L = ψ̄
i∂µ − gvV µ −
τ · bµ
−e (1 + τi3)
− (M − gsφ)
(∂µφ∂
µφ−m2sφ2)−
κφ3 − 1
λφ4 − 1
m2vVµV
µ − 1
Bµν ·Bµν +
m2ρbµ · bµ −
+g2ρbµ · bµ[Λsg2sφ2 + Λvg2vVµV µ], (21)
where Ωµν , Bµν and Fµν are defined after eq.(1). The
parameters of the model are again the masses and the
couplings, which are now constants, i.e., gs replaces Γs,
gv replaces Γv and gρ replaces Γρ. Non-linear σ terms are
also included. We have followed the prescription of [3],
where the starting point was the NL3 parametrization
and the gρ coupling was adjusted for each value of the
coupling Λi studied in such a way that for kF = 1.15
fm−1 (not the saturation point) the symmetry energy is
25.68 MeV. In the present work we set Λs = 0 as in [37].
Notice that other possibilities for this model with σ − ρ
and ω − ρ couplings have already been discussed in the
literature as in [4], for instance.
The mesonic equations of motion in the Thomas-Fermi
approximation become
∇2φ = m2sφ− gsρs +
φ3 (22)
∇2V0 = m2vV0 − gvρ+ 2Λvg2v V0 g2ρb20, (23)
∇2b0 = m2ρb0 −
ρ3 + 2Λvg
0 , (24)
∇2A0 = −eρp, (25)
and the expression for the energy reads
∫ kFi(r)
(2π)3
p2 +M∗2
(∇φ)2 +m2sφ2 − (∇V0)2 −m2vV 20
−(∇b0)2 −m2ρb20 − (∇A0)2
+gvV0ρ+
ρ3b0 + eA0ρp
φ4 − Λvg2vV 20 g2ρb20
. (26)
All other expressions are very similar to the ones ob-
tained from the TW model and can be read off from
them bearing in mind that the density dependent cou-
plings have to be replaced by the constant couplings. In
particular the chemical potentials do not contain the re-
arrangement term ΣR0 .
IV. NUMERICAL RESULT VIA A
NUCLEATION PROCESS
At this point, eqs. (12-15) for the TW model and eqs.
(22-25) for the NLωρmodel have to be solved numerically
in a self-consistent way and hence, initial and boundary
conditions for each equation are necessary. One of the
methods we use here is based on a prescription given in
[38], where these conditions are obtained from a situation
of phase coexistence in a mean field approximation with
classical meson fields and no electromagnetic interaction.
The method is well explained in [38] and, as we are using
different models here, just the main equations are written
next.
For the TW model, the equilibrium equations for ho-
mogeneous matter for the fields are:
m2sφ− Γs ρs = 0, (27)
m2vV0 − Γv ρ = 0, (28)
m2ρb0 −
ρ3 = 0, (29)
and for the energy and pressure density:
E = 1
∫ kFi
p2 +M∗2
V 20 +
b20, (30)
∫ kFi
V 20 +
+ρΣR0 . (31)
For the NLωρ model, the equilibrium equations for ho-
mogenous matter, energy density and pressure become:
m2sφ− gsρs +
φ3 = 0, (32)
m2vV0 − gvρ+ 2Λvg2v V0 g2ρb20 = 0, (33)
m2ρb0 −
ρ3 + 2Λvg
0 = 0, (34)
E = 1
∫ kFi
p2 +M∗2
2 −m2vV 20 −m2ρb20
+ gvV0ρ+
φ4 − Λvg2vV 20 g2ρb20. (35)
∫ kFi
V 20 +
φ3 − λ
φ4 + Λvg
0. (36)
Based on the geometrical construction and Gibbs con-
ditions for phase coexistence, i.e., the pressure and both
chemical potentials are equal in both phases, we build
the binodal section given in Fig. 1. Notice that we have
defined the proton fraction of the system as
. (37)
The binodal section yields the boundary conditions which
we need. For the same pressure, two points, with differ-
ent proton fractions are found. For each of these points,
the meson fields and the densities are well defined and
used as the initial and boundary conditions in eqs. (12-
15), which are then solved. Once the meson fields are
obtained, all the quantities that depend on them, as the
energy, pressure densities, chemical potentials, baryonic
densities, etc are also computed. The solution is a droplet
with a certain proton fraction surrounded by a gas of
neutrons. If stable nuclei are calculated, the gas vanishes
because the energy of the system lies below the neutron
drip line and the finite nuclei properties are easily calcu-
lated. This is the general method, but the results depend
strongly on the model used because of the reasons dis-
cussed in Section VI.
V. NUMERICAL RESULT WITHIN A
HARMONIC OSCILLATOR BASIS
Here a different prescription for solving the equations
of motion and the thermodynamical quantities within the
Thomas-Fermi approximation is used. According to [39],
meson field equations of motion of the Klein-Gordon type
with sources can be carried out by an expansion in a com-
plete set of basis states. The harmonic oscillator func-
tions with orbital angular momentum equal to zero are
then chosen. The oscillator length is given by
, b0 =
, (38)
where M is the nucleon mass and ω0 is the oscillator
frequency. The meson fields and their corresponding in-
homogeneous part can be expanded as
Λ(r) =
ΛnRn0(r), SΛ(r) =
SΛnRn0(r), (39)
where Λ(r) = φ(r), V0(r), b0(r) and
Rnl(r) =
l+1/2
n−1 (x
2)exp(−x2/2), (40)
where x = r/b0 is the radius measured in units of the
oscillator length,
Nnl =
2(n− 1)!/(l + n− 1/2)! (41)
is the normalization constant and Lmn (x
2) are the asso-
ciated Laguerre polynomials. For the calculation of the
meson fields l = 0 in the expressions given below. Once
the ansatz given by eqs.(39) are substituted into eqs.(12-
14), a set of inhomogeneous equations is obtained:
Hnn′Λn′ = SΛn (42)
where
Hnn′ = δnn′
b−2B (2(n− 1) + 3/2) +m
+δnn′+1b
n(n+ 1/2) + δn+1n′b
n′(n′ + 1/2).
vΛ =0.01
vΛ =0.025
0 0.1 0.2 0.3 0.4 0.5
FIG. 1: Binodal section for the NL3, TW and NLωρ
parametrizations.
Only the massive fields can be calculated with this
method because the convergence of the Coulomb field,
which has a long range, is very slow. The Green’s func-
tion method is then chosen to describe the electromag-
netic interaction:
A0(r) = e
r′2dr′ρp(r
′)Gc(r, r
′), (44)
Gc(r, r
1/r for r > r′
1/r′ for r′ > r.
VI. RESULTS
A. Parity Violating Electron Scattering and the
Neutron Radius
We start this section by defining the asymmetry for
polarized electron scattering of a hadronic target as
A = dσ+/dΩ− dσ−/dΩ
dσ+/dΩ+ dσ−/dΩ
, (46)
where dσ±/dΩ is the differential cross section for initially
polarized electrons with positive(+) and negative (−) he-
licities. As the electromagnetic interaction is not sensi-
tive to the above difference, the asymmetry becomes de-
pendent of the weak interaction between the electron and
the target. Moreover, we know from the Standard Model
that the neutral Z-boson couples more strongly to the
neutron than to the proton. Those reasonings were then
used in [21] to first propose a clean way to determine
the neutron distribution in nuclei. If we consider elastic
scattering on an even-even target nucleus, the asymmetry
can be written in the form:
V + β
ρn(q)
ρp(q)
]. (47)
In the above expression, G, α, a and β
V are Standard
Model coupling constants as defined in [21], q is the trans-
ferred momentum by the electron to the nucleus and,
ρn(p)(q) =
d3r j0(qr)ρn(p)(r), (48)
ρn(p)(r) being the neutron (proton) distribution in con-
figuration space and j0 the spherical Bessel function of
order zero. It is then clear that a small q measurement
of the asymmetry gives the neutron radius of the distri-
bution once the proton radius is well known. The proton
and neutron mean-square radius are defined as
R2i =
d3rr2ρi(r)
d3rρi(r)
, i = p, n. (49)
The neutron skin thickness is defined as
θ = Rn −Rp. (50)
In the PREX experiment mentioned in the Introduc-
tion, the asymmetry is expected to be measured at
q ≈ 0.4 fm−1 [22]. Also, because the target is a heavy nu-
cleus (208Pb), the above results for the asymmetry should
be reconsidered for a detailed comparison with the exper-
iment, since they were obtained using a Plane Wave Born
Approximation for the electron [43]. For our present pur-
poses, eq. (47) is sufficient to illustrate the sensitivity to
the different model parametrizations and is used next in
the presentation of our numerical results.
The surface energy per unit area of the droplets in
the small surface thickness approximation, excluding the
electromagnetic field, reads [38]
. (51)
However, as the electromagnetic interaction does not con-
tribute to surface properties directly, we have kept the
same definition for the surface energy.
In Table II we show the neutron and proton radius, the
neutron skin thickness, the binding energy and the sur-
face energy obtained within the Thomas-Fermi approxi-
mation and the two different numerical prescriptions de-
scribed in the previous sections. All the results are sensi-
tive to the numerical calculation although the analytical
approximation is the same. When the nucleation method
is performed, the neutron radius is systematically larger,
what results in a thicker neutron skin. This is correlated
with the fact that the surface energy is lower within the
nucleation calculation than within the harmonic oscilla-
tor method. Within the same numerical prescription,
the neutron skin thickness is smaller with the TW model
than with the NL3. As the coupling strength Λv increases
in the NLωρ model, the results move from the original
NL3 to the TW results for all quantities, except the pro-
ton radius, which oscillates a little. We have also in-
cluded the results obtained with the HS parametrization
[44] because we have used this parametrization in order
to compare the TF and the Dirac results for the cross sec-
tions, as discussed in the following. As this parametriza-
tion is known not to give as good results as the other
parametrizations of the NLWM for finite nuclei, we do
not comment on the results it provides. Notice that the
experimental radius for the protons is obtained from the
charge radius Rc and it is given by Rp =
R2c − 0.64 in
fm [39]. Our results can be compared with experimental
and other theoretical results found in the literature. The
proton radius, which is known to better than 0.001 fm is
better described within the TW model. This quantity is
practically independent of the ω− ρ interaction strength
in the NLωρ model as far as the HO numerical prescrip-
tion is used. The neutron radius, on the other had, is
strongly model dependent with drastic consequences in
the neutron skin thickness calculation. The experimen-
tal values for θ are still very uncertain and all our re-
sults fall inside the experimental confidence interval. We
shall comment on possible restrictions to the neutron skin
thickness in the next section. NL3 provides the best re-
sults for the binding energy.
In [25], the results shown for the proton and neutron
radius are respectively 5.52 and 5.85 fm, yielding a skin
of 0.33 fm, larger than ours. Notice, however, that in
[25] the Dirac equation was explicity solved. In [4], the
authors obtained a value of 0.21 fm for the neutron skin
thickness and a binding energy of -7.89 MeV within a
different parametrization of the NLωρ model. Again in
this case the Dirac equation was solved.
In Fig. 2 we show the difference between neutron and
proton densities at the Pb surface for the models dis-
cussed in the present work with the Thomas-Fermi ap-
proximation solved in a harmonic oscillator basis. While
the curves deviate a little in between 6.0 and 8.0 fm, at
the very surface they are similar, but a small discrepancy,
reflecting the differences in the neutron skin can be seen.
In Fig. 3 we display again the difference between neu-
tron and proton densities within both numerical calcula-
tions of the TW and NL3 models. These two Thomas-
Fermi calculations should have given more similar results.
However the nucleation method predicts a very small sur-
face energy for the NL3 parametrization, and therefore,
a large radius. This may be related to the choice of the
boundary conditions and a deeper comparison between
the two methods will be pursued.
Next we present our results for the asymmetry given
by eq. (47) as a function of the transfered momentum.
We begin with Fig.4 which displays the results for the HS
parametrization of the Walecka model. The curve labeled
no structure means the case where Zρn(r) = Nρp(r) and
the other two curves are obtained within the TF approx-
imation and the full solution of the Dirac equation in
6 7 8 9
0.000
0.005
0.010
0.015
0.020
0.025
NL ( v )
NL ( v )
r(fm)
FIG. 2: Difference between neutron and proton densities ob-
tained with the Thomas-Fermi approach solved in a harmonic
oscillator basis for the models discussed in the present work.
the Hartree approximation. At the momentum trans-
fer values of recent experimental interest (around ≃ 0.4
fm−1), the curves are almost identical. A careful analy-
sis of the same results in a different scale shows us that
the asymmetry changes 12 and 11 percent respectively
within the Dirac and TF approximations in comparison
with the no structure case. Since it is the measurement
of the asymmetry in this low momentum transfer region
that will provide the accurate result for the neutron skin
thickness, we have restricted our calculations to the TF
approximation, as stated in the Introduction.
In Fig. 5a we show the asymmetry obtained with the
NL3 model for both numerical calculations in the TF ap-
proximation, i.e, nucleation and HO expansion methods.
In this case, the agreement is very satisfactory even for
larger q-values, although the small numerical discrepan-
cies is reflected in a ∼ 10 percent difference in the pre-
dicted neutron skin thickness, as can be seen from Table
II. Finally, in Fig. 5b our results for the NLωρ (using
two different values for the ω− ρ coupling constant) and
the TWmodels within the HO numerical prescription are
shown. Again, at low momentum transfers, all curves co-
incide. However, it should be noticed that even for two
different model parametrizations which lead us to identi-
cal neutron skin thicknesses, a measurement of the asym-
metry in a higher q-region with a modest experimental
precision, can distinguish between them. Also, we should
expect that the asymmetry presents more structure in
this high momentum transfer region if we solve the Dirac
equation instead of using the TF approach, once the high
q value region is much more sensitive to the central part
of the neutron distribution, which is known to be flat in
the TF approximation. These differences can be seen in
Fig.4.
NL3 nucl
TW nucl
NL3 HO
TW HO
r(fm)
FIG. 3: Difference between neutron and proton densities ob-
tained with the Thomas-Fermi approach solved with both nu-
merical prescriptions for the TW model.
VII. DIFFERENT EOS, DIFFERENT NEUTRON
SKINS
For the sake of completeness, at this point, we discuss
some of the differences between the TW, the NLωρ mod-
els and the NL3 parametrization of the NLWM.
From Fig. 1 one can see that the largest possible pres-
sure for a phase coexistence in the TW model is much
lower, and appears at a lower proton fraction than the
NL3 model. This gives rise to a thinner crust within
the TW model, which may imply that the more exotic
pasta shapes will not form [5]. The NLωρ model goes
on a different direction, i.e., the pressure becomes higher
than the one obtained with the NL3 as the Λv coupling
is turned on.
Although the nuclear matter properties fitted to
parametrize the models are quite similar (see Table I),
the way the EoS behaves when extrapolated to higher
or lower densities can vary a lot from a density depen-
dent hadron model to one of the parametrizations of the
NLWM. Moreover, as seen from Table I, although the ef-
fective mass at saturation density is lower with the TW
than with the NL3, it can accommodate hyperons if an
EoS for stellar matter is necessary, contrary to the usual
0.0 0.5 1.0 1.5 2.0
no structure
HS-Dirac
HS-TF A
q(fm-1)
0.30 0.32 0.34 0.36 0.38 0.40 0.42 0.44 0.46 0.48 0.50
2.5x10-7
5.0x10-7
7.5x10-7
1.0x10-6
1.3x10-6
1.5x10-6
no structure
HS-Dirac
HS-TF
q(fm-1)
FIG. 4: Parametrization HS, comparison Thomas-Fermi-HO
versus Dirac-HO
NL3 parametrization [32–34].
TABLE I: Nuclear matter properties.
NL3 NLωρ TW
[25] [36] [24]
Λv = 0.01 Λv = 0.02 Λv = 0.025
B/A (MeV) 16.3 16.3 16.3 16.3 16.3
ρ0 (fm
−3) 0.148 0.148 0.148 0.148 0.153
K (MeV) 271 271 271 271 240
Esym. (MeV) 37.4 34.9 33.1 32.3 32.0
M∗/M 0.60 0.60 0.60 0.60 0.56
L (MeV) 118 88 68 61 55
Ksym (MeV) 100 -46 -53 -34 -124
Another quantity of interest in asymmetric nuclear
matter is the nuclear bulk symmetry energy, shown in
Table I for the saturation point. The differences in the
symmetry energy at densities larger than the nuclear sat-
uration density is still not well established, but has al-
ready been extensively discussed in the literature even
0.0 0.5 1.0 1.5 2.0
q(fm-1)
NL3 nucl
NL3 HO
no structure
0.0 0.5 1.0 1.5 2.0
q(fm-1)
NL ( v=0.01)
NL ( v=0.025)
FIG. 5: Asymmetry obtained with a) NL3 with both numer-
ical prescriptions and b)parametrizations NLωρ and TW
for the TW model [7, 8, 16, 32]. Again, for the sake
of completeness we reproduce these results here because
the neutron skin thickness and the neutron star EoS are
related by this quantity [1–4], which is usually defined
as Esym = 12
∂2E/ρ
, with δ = −ρ3/ρ = 1 − 2yp. The
symmetry energy can be analytically rewritten as
Esym =
ρ, (52)
for the TW model and as
Esym =
ρ (53)
with the effective ρ-meson mass defined as [3]
= m2ρ + 2g
for the NLωρ model. In both cases
kFp = kF (1 + δ)
1/3, kFn = kF (1− δ)1/3,
with kF = (1.5π
2ρ)1/3 and ǫF =
k2F +M
∗2. In equa-
tions (52) and (53) the second term dominates at large
densities. It is seen that the non-linear ρ − ω terms
introduce a non-linear density behavior in the symme-
try energy of the NLWM parametrizations such as NL3
and TM1. In TW the non-linear density behavior en-
ters through the density dependent coupling parameters.
These non-linear density behavior is important because
the linear behavior of NL3 and TM1 parametrizations
predicts too high symmetry energy at densities of impor-
tance for neutron star matter which has direct influence
on the proton fraction dependence with density. From
Fig. 6, it is easily seen that the symmetry energy ob-
tained with the TW model behaves in a very different
way, as compared with NL3. In [4] a relation between
the symmetry energy and the nuclear binding energy is
discussed : the harder the EoS, the more the symmetry
(fm )−3ρ
vΛ =0.01
vΛ =0.025
0 0.05 0.1 0.15 0.2 0.25 0.3 0.35 0.4
FIG. 6: Symmetry energy for the NL3, TW and NLωρmodels.
energy rises with density. The density dependence dis-
cussed in [4] is of the type introduced in [3, 36] through
the inclusion of a σ− ρ and/or ω− ρ couplings and then,
similar with the NLωρ model discussed here. One can
observe that as the strength of the coupling increases,
the symmetry energy gets closer to the TW curve. In
fact, in [8] it was shown that once this kind of coupling
is introduced with a reasonable strength, the symmetry
energy at low densities tends to behave as the TW model.
The symmetry energy can be expanded around the nu-
clear saturation density and reads
Esym(ρ) = Esym(ρ0) +
ρ− ρ0
ρ− ρ0
where L and Ksym are respectively the slope and the
curvature of the nuclear symmetry energy at ρ0 and they
are calculated from
L = 3ρ0
∂Esym(ρ)
|ρ=ρ0 Ksym = 9ρ20
∂2Esym(ρ)
|ρ=ρ0 .
These two quantities can provide important information
on the symmetry energy at both high and low densi-
ties because they characterize the density dependence of
the energy symmetry. In a recent work [49], the authors
found a correlation between the slope of the symmetry
energy and the neutron skin thickness. In their work 21
sets of the non-relativistic Skyrme potential were inves-
tigated and only 4 of them were shown to have L values
consistent with the values extracted from experimental
isospin diffusion data from heavy ion collisions. In fact,
the extracted value was L = 88 ± 25 MeV [50], which
gives a very strong constraint on the density dependence
of the nuclear symmetry energy and consequently on the
EoS as well. A detailed analysis of Table I shows that, if
this constraint is to be taken seriously, neither the NL3
nor the TW model satisfy it. Nevertheless, the NLωρ
slope interpolates beautifully between the NL3 and TW
slope values. Once again it is seen that the increase in
Λv approximates the NL3 model values for the slope and
energy symmetry to the TW values. Moreover, we have
also tried to find a correlation between the θ values shown
in Table II and L values displayed in Table I. We found
that, as far as some numerical imprecision are consid-
ered, larger values of L correspond to larger values of the
neutron skin, as seen in Fig. 7.
Let’s now go back to the problem of solving the dif-
ferential equations within the nucleation numerical pre-
scription. As we need boundary conditions arising from
the liquid-gas phase coexistence in order to solve eqs.
(12-15) for the TW model and eqs. (22-25) for the NLωρ
model, the binodal sections are essential and the spinodal
sections, which separate the regions of stable to unsta-
ble matter are also of interest. If we had displayed the
binodals in a ρp versus ρn plot, as it is done with the
spinodals in Fig 8, we could see that the spinodals sur-
faces lie inside the binodal sections and share the critical
point corresponding to the highest pressure.
In Fig. 8 the spinodals for the three different mod-
els discussed in this work are shown. Once again, some
of these results can also be found in the recent litera-
ture [7, 8], but we include them here to make a direct
link with the binodals. The instability of the ANM sys-
tem is essentially determined by density fluctuations in
the isoscalar channel. Although the spinodals are, by
themselves, not relevant in calculations performed at the
thermodynamical equilibrium, the isospin channel is very
sensitive to the instabilities occurring below the nuclear
saturation density. The spinodal is determined by the
values of pressure, proton fraction and density for which
the determinant of
Fij =
∂ρi∂ρj
, (56)
where F is the free energy density, goes to zero. A de-
tailed analysis of this quantity can be found in [8, 42].
From Fig. 8, it is seen that the instability region in the
ρp/ρn plane, defined by the inner section of the spinodal
curve is larger for the TW than for the NL3 model. The
size of the instability region depends on the derivative
of the chemical potentials with respect to the neutron
and proton densities. At low densities different models
exhibit different behaviors.
The presence of the rearrangement term in the TW
model also plays a decisive role. Even though a rela-
tively large compensation exists between scalar and vec-
tor mesons in the isoscalar channels within the rearrange-
ment term at low densities, the spinodal region is defined
by the derivative of the chemical potential and therefore
of the rearrangement term.
Next we examine the spinodals obtained with differ-
ent coupling strengths for the NLωρ model. As seen in
Fig. 8, there is almost no difference between the different
curves. They all fall around the original NL3 curve but
once again, they tend to the TW curve as the coupling
strength increases. However, contrary to the TW model,
it was shown in [9] that the direction of the instability in
Λ =0.01
Λ =0.02
Λ =0.025v 0.16
0.18
0.22
0.24
50 60 70 80 90 100 110 120
L (MeV)
FIG. 7: Correlation between the neutron skin θ and the slope
of the symmetry energy L.
NLωρ increases distillation as the density increases, and
the larger the coupling Λv the larger the effect.
Finally, to end this section, let’s make our points clear:
we have used a simple mean field theory approach to ob-
tain the boundary conditions for the equations of motion
of the meson fields in the nucleation prescription. These
boundary conditions depend on the model used and are
intrinsically related with the liquid-gas phase transition
which, in turn, can be well understood by studying the
coexistence surfaces of the corresponding models. On
the other hand, the neutron skin thickness shows a lin-
ear correlation with the slope of the symmetry energy,
as already pointed out in [49] for non-relativistic mod-
els. Based on the different behaviors found with density
dependent hadronic models and the NLWM, an obvious
consequence is the fact that the neutron skin thickness
depends on the choice of the model.
VIII. CONCLUSIONS
We have calculated the 208Pb neutron skin thickness
with two different density dependent hadronic models,
the TW and the NLωρ model, and one of the most used
parametrizations of the NLWM, the NL3. The calcu-
lations were done within the Thomas-Fermi approxima-
tion, which gives quite accurate results for the asymme-
try in the momentum transfer range of interest for the
calculation of neutron skins. In implementing the nu-
merical results two different prescriptions were used: the
first one based on the nucleation process and the second
one based on the harmonic oscillator basis method. We
have seen that when the nucleation method is performed,
the neutron radius is systematically larger, what results
vΛ =0.01
vΛ =0.025
ρ (fm )−3
0.02
0.04
0.06
0.08
0 0.02 0.04 0.06 0.08
FIG. 8: Spinodal section in terms of ρp versus ρn for the NL3,
TW and NLωρ models.
in a thicker neutron skin. This is a consequence of the
fact that the surface energy is lower within the nucleation
calculation than within the harmonic oscillator method.
Within the same numerical prescription, the neutron skin
thickness is smaller with the TW model than with the
NL3. As the coupling strength Λv increases in the NLωρ
model, the neutron skin thickness moves from the orig-
inal NL3 towards the TW results. We have also found
that although the neutron skin thickness is model depen-
dent, the asymmetry at low momentum transfers (below
0.5 fm−1) is very similar for all models and all numerical
prescriptions. As q increases, the asymmetry also be-
comes model dependent. The density profiles obtained
from the solution of the Dirac equation exhibits oscil-
lations near the center of the nucleus, behavior which is
not reproduced within the Thomas-Fermi approximation.
This fact shows up in the asymmetry at large momentum
transfers and therefore all the calculations should be re-
produced by solving the Dirac equation. This calculation
is already under investigation.
It is worth mentioning that the neutron skin thickness
has shown to give hints on the equations of state that
are suitable to describe neutron stars. Moreover, in [49]
a correlation between the slope of the symmetry energy
and the neutron skin thickness was found for Skyrme-
type models. We have observed that this correlation was
also present in the density dependent models we have
studied in the present work.
ACKNOWLEDGMENTS
This work was partially supported by CNPq(Brazil),
CAPES(Brazil)/GRICES (Portugal) under project
100/03 and FEDER/FCT (Portugal) under the projects
POCTI/FP/63419/2005 and POCTI/FP/63918/2005.
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TABLE II: 208 Pb properties
model approximation Rn Rp θ B/A σ
(fm) (fm) (fm) MeV Mev/fm2
NL3 TF+nucleation 5.88 5.65 0.24 -7.77 0.76
NL3 TF+HO 5.79 5.57 0.22 -7.79 0.96
NLωρ, Λv = 0.01 TF+HO 5.77 5.57 0.20 -7.73 0.98
NLωρ, Λv = 0.02 TF+HO 5.75 5.57 0.17 -7.65 0.99
NLωρ, Λv = 0.025 TF+HO 5.74 5.58 0.16 -7.63 1.00
TW TF+nucleation 5.71 5.50 0.22 -6.42 1.08
TW TF+HO 5.68 5.52 0.16 -7.46 1.10
HS TF+HO 5.70 5.47 0.24 -6.10 1.37
exp.[45] 5.44
exp. [46] -7.87
exp. [47] 0.12± 0.07
exp. [48] 0.20± 0.04
|
0704.0409 | On the over-barrier reflection in quantum mechanics with multiple
degrees of freedom | CERN-PH-TH/2007-065
On the over-barrier reflection in quantum mechanics
with multiple degrees of freedom
D.G. Levkova1, A.G. Panina,b2, S.M. Sibiryakovc,a3
aInstitute for Nuclear Research of the Russian Academy of Sciences,
60th October Anniversary prospect 7a, Moscow 117312, Russia.
bMoscow Institute of Physics and Technology,
Institutskii per. 9, Dolgoprudny 141700, Moscow Region, Russia.
cTheory Group, Physics Department, CERN, CH-1211 Geneva 23, Switzerland.
Abstract
We present an analytic example of two dimensional quantum mechanical system,
where the exponential suppression of the probability of over–barrier reflection changes
non-monotonically with energy. The suppression is minimal at certain “optimal” ener-
gies where reflection occurs with exponentially larger probability than at other energies.
1 Introduction
Tunneling and over–barrier reflection are the characteristic non–perturbative phenomena in
quantum mechanics. They typically occur with exponentially small probabilities,
P ∝ e−F/~ , (1)
where F is the suppression exponent; still, the above phenomena are indispensable in under-
standing a wide variety of physical situations, from the generation of baryon number asym-
metry in the early Universe [1] to chemical reactions [2] and atom ionization processes [3].
During the last decades extensive investigations of tunneling processes in systems with
many degrees of freedom have been performed [2, 4, 5, 6, 7, 8, 9, 10, 11, 12, 13]. These studies
[email protected]
[email protected]
[email protected], [email protected]
http://arxiv.org/abs/0704.0409v2
revealed a rich variety of features of multidimensional tunneling which are in striking contrast
to the properties of one–dimensional tunneling and over–barrier reflection. In particular, the
following phenomenon has been observed: the probability of tunneling may depend non-
monotonically on the total energy of the system and exhibit resonance-like peaks. One can
envisage three physically different mechanisms of this phenomenon. The first mechanism,
present already in one-dimensional case, is tunneling via creation of a metastable state.
In this case the tunneling probability at the maximum of the resonance is exponentially
higher than at other energies. On the other hand, the resonance width ∆E is exponentially
suppressed; so, after averaging with an energy distribution of a finite width the effect of the
resonance is washed out in the semiclassical limit ~ → 0. The second possible mechanism of
non-monotonic behavior of P(E) is quantum interference [7, 13] (see also [14]). In this case
the peak value of the tunneling probability is only by a factor of order one higher than the
average value, while the width of the resonance scales as ∆E ∝ ~. Again, the resonances
become indiscernible in the semiclassical limit. In both these cases the resonances can
be attributed to the subleading semiclassical corrections, i.e. non-monotonic behavior of
the pre-exponential factor omitted in Eq. (1). The third possibility is that the suppression
exponent F (E) is non-monotonic. In this case the existence of the “resonances” is the leading
semiclassical effect: the optimal tunneling probability at the maximum of the resonance is
exponentially higher than the probability at other energies. At the same time the resonance
width scales as4 ∆E ∝
~. This last possibility of “optimal tunneling” is definitely of
interest; yet, it did not receive much attention in literature. We are aware of only a few works
mentioning non–monotonic dependence of the suppression exponent on energy [15, 16, 14].
It is worthwhile studying this phenomenon in detail; this can provide a new insight into the
dynamics of multidimensional tunneling.
In this paper we consider the process of over–barrier reflection in a simple model with two
degrees of freedom. Our setup is interesting in two respects. First, the model under study
is essentially non–linear and the variables cannot be separated; still, over–barrier reflections
in this model can be described analytically within the semiclassical framework. Thus, this
model can serve as an analytic laboratory for the study of multidimensional tunneling. Sec-
ond, the suppression exponent F of the reflection process behaves non–monotonically as the
4This follows from the representation
P(E) ∝ exp
F (Eo)
F ′′(Eo)(E − Eo)2
of the tunneling probability in the vicinity of the maximum.
total energy E changes. We demonstrate that the function F (E) possesses a number of
local minima E = Eo, where reflection is optimal. We stress that the process we study is
exponentially preferable at “optimal” energies as compared to other energies.
Our model describes the motion of a quantum particle in the two dimensional harmonic
waveguide (see Refs. [8, 10, 14] for similar models). The Hamiltonian is
w2(x, y) ,
where x, y are the Cartesian coordinates and m is the mass of the particle. The function
U = mω2w2/2 represents the waveguide potential in two dimensions: a particle with small
energy is bound to move along the line w(x, y) ≈ 0. We do not introduce a potential barrier
across the waveguide and consider the case when the line w = 0 stretches all the way from
x → −∞ to x → +∞. We also assume that the function w(x, y) is linear in the initial
asymptotic region,
w(x, y) → y as x→ −∞ .
In the present paper we consider two particular cases of the function w(x, y) describing
waveguides with one and two sharp turns5, see Fig. 1.
The motion of the particle at x → −∞ is a superposition of free translatory motion in
x direction and oscillations of frequency ω along y coordinate; the state of such a particle is
fully characterized by two quantum numbers, the total energy E and y–oscillator excitation
number N . The particle sent into the waveguide from the asymptotic region x→ −∞ with
given E, N may either continue to move towards x → +∞, or reflect back into the region
x→ −∞. We are interested in the probability P(E,N) of reflection.
Let us discuss reflections at the classical level. [Note that the classical counterpart of N
is the energy of transverse oscillations.] Consider first the waveguide with one sharp turn
(Fig. 1a). One observes that the outcome of the classical evolution, i.e. whether or not
the particle reflects from the turn, depends not only on the total energy E, but on other
dynamical quantities as well. In particular, the direction of the momentum of the particle
in the vicinity of the turn (point C on the graph) is important. This means that the entire
dynamics in the waveguide should be taken into account in order to determine the possibility
of classical reflection. This is in sharp contrast with the situation in one–dimensional case,
where reflection from the potential barrier (or transition through it) is ensured by the value
of the conserved energy of the particle.
5The explicit expressions for the waveguide functions w(x, y) will be presented in the subsequent sections.
Figure 1: The equipotential contour U = E for the waveguides with (a) one and (b) two
sharp turns. An example of classical trajectory is shown in the case (b).
Now, consider the waveguide with two turns. The model is characterized by the angles
of the turns and the distance L between them (see Fig. 1b). Suppose the particle starts
moving classically from x → −∞ with N = 0 along the valley w = 0. Then, the transverse
oscillations get excited only after the particle crosses the first turn, point C ′ on the plot, so
that at the time of arrival to the second turn (point C) approximately ωτ/2π oscillations
are made, where τ ∼ L
m/2E is the time of motion between the two turns. The state of
the particle (coordinates and momenta) at which it comes across the second turn depends
periodically on the phase of transverse oscillations ωτ . Hence, one expects that the regime
of motion of the classical particle can change from transmission to reflection and back as the
energy grows (τ decreases); the energies where it happens can be roughly estimated as
mω2L2
2(2πn)2
. (2)
We will see that this is indeed the case for the waveguides with certain angles of the turns.
At some values of E, N the reflection process cannot proceed classically. Then, at the
quantum mechanical level its probability is exponentially suppressed, F (E,N) > 0. It is
natural to call such a process “over–barrier reflection”6. The central quantity to be studied
below is the suppression exponent F (E,N) of this process. The above discussion suggests
that F (E,N), being determined by the entire dynamics in the waveguide, may be a highly
non–trivial function. For the particular case of the waveguide with alternating regimes of
6By this term we want to emphasize that the process is classically forbidden. Recall, however, that there
is no actual potential barrier across the waveguide in our setup.
classical reflections and transmissions F should oscillate: F = 0 at the energies where the
classical reflections are allowed, and F > 0 at the energies where the reflections are classically
forbidden. One can expect that the similar oscillatory behavior of the suppression exponent
persists for other two–turn models as well. Now, instead of reaching zero, F may possess
a number of local positive minima implying that the reflection at the “optimal” energies is
still a tunneling process.
Let us emphasize the difference of the “optimal tunneling” from quantum interference
and resonance phenomena in our two–turn model. The interference of the de Broglie waves
reflected from the two turns can, in principle, lead to oscillations in the reflection proba-
bility P(E). One can estimate the positions of the interference peaks by equating the De
Broglie wavelength of the particle to an integer fraction of the distance between the turns,
2mE ∼ L/n. This yields the energies of the interference peaks,
Eintn ∼
(2πn)2~2
This formula is completely different from Eq. (2) for the peaks due to “optimal tunneling”.
In particular, the distance between the adjacent inteference peaks,
∆Eint ∼
scales proportional to ~. Thus, these peaks should be averaged over in the semiclassical
limit. Besides, the amplitude of the interference peaks is at most of order one and does
not affect the suppression exponent. Indeed, the exponential increase of the scattering
amplitude can arise due to quantum interference only in the presence of a resonant state
with exponentially long life–time. This state should be supported somewhere in between the
turns and should be classically stable. In Sec. 4.2 we show that such states are absent in
our system. One concludes that the peak–like structure of the probability P(E) of “optimal
tunneling” is caused by completely different physical reasons as compared to the case of
resonance scattering in quantum theory.
It is worth noting that the phenomenon of “optimal tunneling” has an important imple-
mentation in field theory. Recently it was argued [17] (see also Ref. [16]) that the probability
of tunneling induced by particle collisions [18, 19] reaches its maximum at a certain “op-
timal” energy and stays constant7 at higher energies. This result, if generic, provides the
7As opposed to the quantum mechanical case, the tunneling probability does not decrease at energies
higher than the “optimal” one. This is due to the possibility, specific to the field theoretical setup, to emit
the excess of energy into a few hard particles, so that tunneling effectively occurs at the “optimal” energy.
answer to the long–standing question [20] about the high–energy behavior of the probability
of collision–induced nonperturbative transitions in field theory. The quantum mechanical
model presented here supports the generic nature of the phenomenon of “optimal tunnel-
ing”; the simplicity of our model enables one to get an intuitive insight into the nature of
this phenomenon.
The paper is organized as follows. In Sec. 2 we review the semiclassical method of complex
trajectories, which is exploited in the rest of the paper. Reflections in the waveguides with
one and two turns are considered in Secs. 3 and 4 respectively. We discuss our results in
Sec. 5. In appendix we analyze the validity of some assumptions made in the main body of
the paper.
2 The semiclassical method
We start by describing the semiclassical method8 of complex trajectories which will be used
in the study of over–barrier reflections. We concentrate on the derivation of the formula
for the suppression exponent F (E,N) (see Refs. [2, 8, 9] for the details of the method and
Ref. [19] for the field theory formulation). In what follows we use the system of units
~ = m = ω = 1 ,
where the Hamiltonian takes the form,
p2x + p
y + w
2(x, y)
. (3)
One starts with the amplitude of reflection into the state with definite coordinates
xf < 0 , yf ,
A = 〈xf , yf |e−iĤ(tf−ti)|E, N〉 . (4)
Here |E, N〉 is the initial state of the particle moving in the asymptotic region xi → −∞
with fixed translatory momentum p0 =
2(E −N) and the oscillator excitation number N .
Semiclassically,
〈xi, yi|E, N〉 = eip0xi cos
′)dy′ + π/4
, (5)
8Note that the method has been confirmed by the explicit comparison with the exact quantum mechanical
results in Refs. [8, 9, 14]; specifically, the recent check [14] deals with the case when the dependence of the
suppression exponent on energy is not monotonic.
where xi, yi denote initial coordinates,
2N − y′2 , (6)
and we omitted the pre-exponential factor which is irrelevant for our purposes. Using Eq. (5),
one rewrites the amplitude (4) as a path integral,
dxidyi
[dx][dy]
xf , yf
xi, yi
eiS+ip0xi cos
′)dy′ + π/4
, (7)
where S is the classical action of the model (3).
In the semiclassical case the integral (7) is dominated by the (generically complex) saddle
point. Note that, as we continue the integrand in Eq. (7) into the plane of complex coor-
dinates, one of the exponents constituting the initial oscillator wave function grows, while
the other becomes negligibly small. Within the validity of our approximation, we omit the
decaying exponent by writing
′)dy′ + π/4
→ exp
′)dy′
, (8)
with the standard choice9 of the branch of the square root in Eq. (6).
One proceeds by finding the saddle point for the integral (7) with the substitution (8).
Extremization with respect to x(t), y(t) leads to the classical equations of motion,
ẍ = −wwx , ÿ = −wwy . (9)
Differentiating with respect to xi ≡ x(ti), yi ≡ y(ti), one obtains,
ẋi = p0 =
2(E −N) , ẏi = py(yi) =
2N − y2i .
The latter equations are equivalent to fixing the total energy E and initial oscillator energy
N of the complex trajectory,
ẋ2i +N , (10a)
ẏ2i + y
. (10b)
9 The correct branch is fixed by drawing a cut between the oscillator turning points y = ±
2N , and
choosing Im py > 0 at y ∈ R, y >
2N , see, e.g., Refs. [21].
Substituting the saddle–point configuration10 into Eq. (7), one obtains the amplitude of the
process with exponential accuracy,
A ∝ eiS+iB(xi, yi) ,
where the term
B(xi, yi) = p0xi +
′)dy′ (11)
is the initial–state contribution. For the inclusive reflection probability one writes,
dxfdyf |A|2 ∝
dxfdyf e
iS−iS∗+iB−iB∗ .
The integral over the final states can also be evaluated by the saddle point technique; ex-
tremization with respect to xf ≡ x(tf ), yf ≡ y(tf) fixes the boundary conditions in the
asymptotic future,
Im ẋf = Im xf = 0 , Im ẏf = Im yf = 0 . (12)
In this way one obtains the expression (1) for the reflection probability, where the suppression
exponent F is given by the value of the functional
F (E, N) = 2 ImS + 2 ImB(xi, yi)
evaluated on the saddle–point configuration — a complex trajectory satisfying the boundary
value problem (9), (10), (12).
The contribution B(xi, yi) of the initial state is simplified after one uses the asymptotic
form of the solution at t→ −∞ (xi → −∞),
x = p0t + x0 , y = ae
−it + āeit . (13)
Equations (10) guarantee that the quantities p0 =
2(E −N) and 2aā = N are real, since
E, N ∈ R. Therefore, one may introduce two real parameters T , θ as follows,
2 Im x0 = −p0T , ā = a∗eT+θ . (14)
One finds for the initial term (11),
2 Im B(xi, yi) = Im
2p0xi − 2Narccos(yi/
2N) + yi
2N − y2i
= −p20T −N(T + θ) + Im(yiẏi) ,
10For simplicity we assume that the saddle–point configuration is unique. Otherwise, one should take the
saddle point corresponding to the weakest exponential suppression.
and thus
F = 2 Im S̃ −ET −Nθ , (15)
where S̃ is the classical action of the system (3) integrated by parts,
S̃ = −1
xẍ+ yÿ + w2(x, y)
. (16)
Let us comment on the physical meaning of the parameters T , θ. Consider two trajectories
which are solutions to the boundary value problem (9), (10), (12) at neighbouring values of
E, N . The differential of the quantity 2 Im S̃ as one deforms one trajectory into the other is
d (2 Im S̃) = d Im(2S + xiẋi + yiẏi) = Im(xidẋi − ẋidxi + yidẏi − ẏidyi) = EdT +Ndθ ,
where in the last equality we used the asymptotic form (13), (14) of the solution. Then,
from Eq. (15) one finds,
dF (E,N) = −TdE − θdN . (17)
Thus, the parameters T and θ are (up to sign) the derivatives of the suppression exponent
with respect to energy E and initial oscillator excitation number N respectively.
Our final remark is that the boundary value problem (9), (10), (12) is invariant with
respect to the trivial time translation symmetry,
t→ t+ δt , δt ∈ R , (18)
which can be fixed in any convenient way.
3 The model with one turn
To warm up, we consider the simplest model, where the waveguide has one sharp turn,
w = y θ(−x+ y tg β) + cos β (x sin β + y cos β) θ(x− y tg β) . (19)
Here θ(x) is the step function. It is convenient to use the rotated coordinate system,
cos β − sin β
sin β cos β
The waveguide function takes the form,
w = η cos β − ξ sin β θ(−ξ) . (20)
Figure 2: The equipotential contour w2(x, y) = 2N for the waveguide (20) and the trajectory
of the critical solution with energy N/ cos2 β.
The equipotential contour w2(ξ, η) = const is shown in Fig. 2. One observes that the motion
of the particle in two regions, ξ < 0 and ξ > 0, decomposes into the translatory motion and
oscillations in the coordinates x, y and ξ, η respectively (see. Eqs. (19) and (20)); the
frequency of η–oscillations in the latter case is cos β.
Due to the presence of the step function, the first derivatives of the potential (20) are
discontinuous11 at ξ = 0. Strictly speaking, the semiclassical method is not applicable in
this situation [21]. Thus, the formula (20) should be regarded as an approximation to some
waveguide function with smooth turn. Generically the width of the smoothened turn is
characterized by a parameter b; the sharp–turn approximation (20) corresponds to b → 0.
An example of smoothening is provided by the following substitution in Eq. (20),
θ(ξ) → θb(ξ) =
1 + e−ξ/b
. (21)
The semiclassical description can be used as long as the de Broglie wavelength of the particle
is small compared to the linear size of the potential12, 1/
E ≪ b. We conclude that the
sharp–turn and semiclassical approximations are valid simultaneously for smooth waveguides
1 ≫ b≫ 1/
E . (22)
11Note that the potential itself is continuous.
12Another semiclassical condition is that the energy is sufficient to excite a lot of oscillator levels, E ≫ 1.
It is satisfied provided Eq. (22) holds.
An important property of the model (20) is invariance of the classical equations of motion
(9) under the rescaling of the coordinates,
x→ Λx , y → Λy . (23)
Using the transformation (23), one may express a solution x(t), y(t) with energy E in terms
of the “normalized” one,
x = x̃
E , y = ỹ
where the solution x̃(t), ỹ(t) has unit energy; its initial oscillator excitation number is
ν = N/E .
The suppression exponent (15) takes the form,
F (E, N) = Efβ(ν) , (24)
where fβ(ν) is the exponent for the “normalized” solution. Substituting the expression (24)
into Eq. (17), one obtains,
fβ(ν) = −T − θν . (25)
We will exploit Eq. (25) in the end of this section. Now, we proceed to finding the “normal-
ized” trajectories.
At certain initial data ν > νcr the particle can reflect from the turn classically, so that
fβ(ν > νcr) = 0 .
Let us find the value of νcr. In the region ξ < 0 the classical solution takes the form,
x(t) = p0t+ x0 , (26a)
y(t) = A0 sin(t + ϕ) . (26b)
Having crossed the line ξ = 0 (line AB in Fig. 2), the classical particle can never return
back into the region ξ < 0. Indeed, in this case it moves at ξ > 0 with constant momentum
pξ > 0. Thus, the particle can reflect classically only if its trajectory touches the line
ξ = 0. The potential of our model has ill–defined derivatives at ξ = 0, and the fate of the
particle moving along the line AB depends on the particular choice of the smoothening of
the potential. In appendix we consider the motion of the classical particle in the case when
nonzero smoothening of width b is switched on. For a class of smoothenings we show that
in the small vicinity (δξ ∼ b) of any trajectory touching the line ξ = 0 there exists some
“smoothened” trajectory, which reflects classically from the turn. Consequently, below we
associate the trajectories touching the line ξ = 0 with the classical reflected solutions.
One notices that the inclination of the trajectory (26) is bounded from above
therefore, the classical trajectory of the particle can touch the line ξ = 0, that is, y/x = ctg β
only at
A0/p0 ≥ ctg β . (27)
From Eqs. (27), (26), (10) one extracts the condition for the particle to reflect classically
from the turn,
ν ≥ νcr = cos2 β . (28)
The critical classical solution at ν = νcr touches the line ξ = 0 at η = 0 (point C in Fig. 2),
where its trajectory
xcr(t) =
2t sin β , (29)
ycr(t) =
2 sin t cos β .
has the largest inclination.
We now turn to the classically forbidden reflections at ν < νcr, which are described by
the boundary value problem (9), (10), (12). One makes the following important observation.
The waveguide function (20) has the form of two analytic functions glued together at ξ = 0.
Hence, the equations of motion (9) can be continued analytically to the complex values of
coordinates in two different ways, starting from the regions ξ < 0 and ξ > 0 respectively. In
this way one obtains two complex solutions, ξ−(t), η−(t) and ξ+(t), η+(t). These solutions
and their first derivatives should be matched at some moment of time t1, ξ(t1) = 0. [Note
that the matching time t1 does not need to be real.] Below we conventionally refer to these
solutions as the ones belonging to the regions ξ < 0 and ξ > 0.
By the same reasoning as above we find that once the particle arrives into the region
ξ > 0, it never reflects back to ξ < 0, unless pξ = 0. So, in the region ξ > 0 one writes,
ξ+(t) = 0 , (30a)
η+(t) =
cos β
sin(t cos β + ϕη) , (30b)
where the “normalization” condition E = 1 has been used explicitly. Due to the conditions
in the asymptotic future, Eqs. (12), the parameter ϕη is real. We use the translational
invariance (18) to set ϕη = 0. Note that we again associate the trajectory going along the
line ξ = 0 with the reflected one.
The physical picture of over–barrier reflection that comes to mind matches with the new
mechanism of multidimensional tunneling proposed recently in Refs. [9, 11]. The process
proceeds in two steps. The first step, which is exponentially suppressed, is formation of the
periodic classical orbit (30) oscillating along the line ξ = 0. This orbit is unstable. At the
second step of the process the unstable orbit decays classically forming a trajectory going
back to x → −∞ at t → +∞. Clearly, the second step does not affect the suppression
exponent of the whole process, and we do not consider it explicitly. In what follows we
concentrate on the determination of the tunneling trajectory describing the first step of the
process.
One should find the solution at ξ < 0 and impose the boundary conditions (10). Note,
however, that the energy of our solution is fixed already. As for the initial oscillator excitation
number ν, it does not change during the evolution in the region ξ < 0. Thus, one may fix it
at the matching time t = t1. One writes,
(ẏ2 + y2)
= cos2 β + sin2 β sin2(t1 cos β) .
This complex equation allows one to express t1 as
sin(t1 cos β) = −i
νcr − ν
sin β
, (31)
where the choice of the sign is dictated by the condition in footnote 9. It is convenient to
introduce notation t1 = iT1, T1 ∈ R.
In order to find the suppression exponent fβ(ν), one needs to evaluate the parameters
T (ν), θ(ν). At ξ < 0 the solution has the form,
x−(t) = p0(t− iT/2) + x′0 , (32a)
y−(t) = ae
−it + a∗eT+θ+it , (32b)
where the definitions (13), (14) have been taken into account explicitly, so that p0, x
0 ∈ R.
One evaluates p0, x
0, a, T , θ by matching the coordinates x±, y± and their first derivatives
νcr0.20.150.10.050
Figure 3: The suppression exponent fβ(ν) for the waveguide (20); β = π/3.
ẋ±, ẏ± at t = iT1; this yields
x′0 = 0 , p0 =
2(1− ν) , a = i
1− ν/ cos2 β
T + θ
cos2 β − ν
sin β
The last two equations, together with Eq.(25), define the function fβ(ν),
fβ(ν) =
cos β
arcsh
νcr − ν
sin β
− ν cos β arcsh
νcr − ν
sin β
(νcr − ν)(1− ν)
this finction is plotted in Fig. 3. One observes that at ν → νcr the quantities T1, T, θ, fβ
tend to zero, and the complex trajectory tends to the classically allowed critical solution, cf.
Eqs. (29),
2 sin β , a→ i√
cos β .
At ν = 0 one has,
fβ(0) = −2 +
cos β
arcth (cos β) . (33)
To summarize, we obtained the suppression exponent for the reflection of a particle in the
simplest waveguide with one sharp turn.
Figure 4: The equipotential contour w2(x, y) = 2N ′ for the waveguide (35) and the trajectory
of the critical solution with energy N ′/ cos2 β > EB. The matching points C, C
′ are shown
by the thick black dots.
4 The model with two turns
4.1 Introducing the system
In the model of the previous section the suppression exponent was proportional to energy
because of the coordinate rescaling symmetry (23). Now, we are going to demonstrate that
small violation of this symmetry results in highly non–trivial graph for F (E).
One introduces a second turn into the waveguide, see Fig. 4. We want to consider this
turn as a small perturbation, so, we assume its angle α to be smaller than β. It is convenient
to introduce two additional coordinate systems, x′, y′ and ξ, η, bound to the central and
rightmost parts of the waveguide respectively. They are related to the original coordinate
system x, y as follows,
cosα sinα
− sinα cosα
cos β − sin β
sin β cos β
x′ − L
Note that the origin of the coordinate system ξ, η is shifted by the distance L. The waveguide
function is
w = θ(−x′)θ(−ξ)y + θ(−ξ)θ(x′)y′ cosα + θ(ξ)η cosα cos β ; (35)
it consists of three pieces glued together continuously at x′ = 0 and ξ = 0 (lines A′B′ and
AB in Fig. 4 respectively). At t→ −∞ the particle comes flying from the asymptotic region
x′ < 0, where w = y. In the intermediate region x′ > 0, ξ < 0 the particle moves in the x′
direction oscillating along the y′ coordinate with the frequency cosα. Finally, in the region
ξ > 0 its motion is free in the coordinates ξ, η; the frequency of η–oscillations is cosα cos β.
The model (35) no longer possesses the symmetry (23): rescaling of coordinates changes
the length L of the central part of the waveguide. In what follows it is convenient to work
in terms of the rescaled dynamical variables,
x̃ = x/L , ỹ = y/L .
In new terms the parameter L disappears from the classical equations of motion, entering
the theory through the overall coefficient L2 in front of the action. The initial–state quantum
numbers are also proportional to L2,
E = L2Ẽ , N = L2Ñ . (36)
Thus, the conditions (22) for the validity of the semiclassical approximation are satisfied in
the limit
L→ ∞ , Ẽ, Ñ = fixed .
The suppression exponent takes the form
F (E,N) = L2F̃ (Ẽ, Ñ) . (37)
To simplify notations, we omit tildes over the rescaled quantities in the rest of this sec-
tion. Rescaling back to the physical units can be easily performed in the final formulae by
implementing Eqs. (36), (37).
4.2 Classical evolution
Let us begin this subsection by demonstrating that there are no stable classical solutions
localized in the region between the turns. This is important for the determination of the
tunneling probability, since such stable solutions could lead to exponential resonances in the
tunneling amplitude. The argument proceeds as follows. Any trajectory which is localized in
the intermediate region should reflect from the line AB infinitely many times. Each reflection
involves touching the unstable orbit living at the line AB. This implies that the trajectory
itself is unstable.
We proceed by determining the region of initial data E, N , which correspond to the
classical reflections. [For brevity we will refer to this region as the “classically allowed
region”, as opposed to the “classically forbidden region” where reflections occur only at the
quantum mechanical level. We stress that these are the regions in the plane of quantum
numbers E, N .] Let us search for the critical classical solutions which correspond to the
smallest initial oscillator number N = Ncr(E) at given energy E. As in the previous section,
one finds that the particle must get stuck at the line13 AB for some time in order to reflect
back. Let us first make an assumption inspired by the study of the one-turn model that the
critical solutions touch the line AB at their maximum inclination point (point C in Fig. 4).
We will see shortly that this is true only at energies above a certain value EB, see Eq. (50).
Still, the analysis based on the above assumption enables one to catch the qualitative features
of the critical line N = Ncr(E). Besides, the analysis is considerably simplified in this case;
we postpone the accurate study until the end of this subsection. Keeping in mind the above
remarks, one writes for the solution in the intermediate region,
x′cr(t) = t
2E sin β + 1 , (38a)
y′cr(t) =
cos β
sin(t cosα) . (38b)
Before entering the intermediate region, the particle crosses the line A′B′ (point C ′ in Fig. 4).
The initial oscillator number N is most conveniently calculated at the moment
t = t0 ≡ −
2E sin β
of crossing. Using the relations (34) one obtains,
ẋcr(t0) =
sin β cosα− cos β sinα cos
cosα√
2E sin β
, (39)
and thus
Ncr(E) = E −
ẋ2cr(t0) = E − E
sin β cosα− cos β sinα cos
cosα√
2E sin β
, E > EB .
0.05
0.15
0.90.80.70.60.50.40.30.10
Figure 5: The boundary N = Ncr(E) of the classically allowed region at E > EB for the
waveguide model (35); β = π/3, α = π/30. The region of the classically allowed initial data
lies above this boundary. The empty circles correspond to the energies E = En, where the
curve N = Ncr(E) touches its lower envelope N = E cos
2(β + α).
As an example, we show in Fig. 5 the region of the classically allowed initial data for β = π/3,
α = π/30. One observes that the function Ncr(E) oscillates between two linear envelopes,
E cos2(β + α) and E cos2(β − α); the period of oscillations decreases as E → 0. Moreover,
the curve Ncr(E) has a number of minima at the points E = E
n . This means that the
energies E = Ecrn are optimal for reflection: in the vicinity of any point E = E
n , N =
Ncr(E
n ) reflections become exponentially suppressed independently of whether the energy
gets increased or decreased. This feature is particularly pronounced in the case α+β = π/2,
when the lower envelope coincides with the line N = 0. Then, the classical reflections (i.e.
reflections with the probability of order 1) at N = 0 are possible only in the vicinities of the
points
8π2(n− 1/2)2
This is the case we used in Introduction to illustrate the effect.
The minima E = Ecrn exist at other values of the parameters as well. For instance, let
13We do not consider reflections from the line A′B′. They disappear at larger values of N than reflections
from the line AB if α is small enough.
us find the positions of these minima in the case α ≪ 1. One differentiates Eq. (40) with
respect to energy and obtains,
Ecrn = En
π(n− 1/2)
arcsin
ctg β
2πα(n− 1/2)
+O(α2)
, (41)
where
8π2(n− 1/2)2 sin2 β
are the points where the curve N = Ncr(E) touches its lower envelope. The argument of
arcsine in Eq. (41) should be smaller than one, so, the minima Ecrn exist only at large enough
n ≥ n0 ≡
ctg β
+ 1 , (43)
where [·] stands for the integer part.
Let us make several comments. First, note that n0 ∼ O(1/α), consequently, all the
optimal points Ecrn lie in the region of small energies E ∼ 1/n20 ∼ O(α2). Second, as we
pointed out before, the formula (40) for the function Ncr(E) holds at E > EB. Comparing
the expressions (42), (43) and (50), one observes that En0 > EB if tg β > 1. So, there
does exist a range of energies where the non-monotonic behavior of the function Ncr(E) can
be inferred from the formula (40). In fact, the conclusion about the existence of the local
minima of Ncr(E), as well as the expressions (41), (42), (43) determining their positions,
remain valid also at E < EB. This follows from the rigorous analysis of the boundary of the
classically allowed region to which we turn now. The reader who is more interested in the
tunneling processes may skip this part and proceed directly to subsection 4.3.
Now, we do not appeal to the Ansatz (38). Instead, we start with the general solution
in the intermediate region,
x′ = p′0(t− t0) , (44a)
y′ = A′0 sin [(t− t0) cosα+ ϕ′] . (44b)
It is convenient to parametrize it by the total energy E = p′20 /2 + cos
2 αA′20 /2 and the
“inclination” γ defined by the relation
p′0/A
0 = tg γ cosα .
Expressions (44) take the following form,
2E (t− t0) sin γ , (45a)
cos γ
sin [(t− t0) cosα + ϕ′] . (45b)
The constants t0 and ϕ
′ are fixed by demanding the trajectory (45) to reflect classically from
the second turn, i.e. touch the line ξ = 0 at t = 0,
(x′ − 1) cos β − y′ sin β
= 0 ,
= ctg β .
These conditions imply,
t0 = −
2E sin γ
tg2 β
tg2 γ
− 1 , (46a)
ϕ′ = −
cosα√
2E sin γ
tg2 β
tg2 γ
− 1− arccos
. (46b)
One sees that the classical reflections are possible only at γ ∈ [0; β]; the boundary value
γ = β reproduces the solution (38).
In order to find Ncr(E), one should minimize the value of the incoming oscillator exci-
tation number with respect to γ at fixed E. At t = t0, when the particle crosses the first
turn,
p0 ≡ ẋ(t0) =
2E(cosα sin γ − sinα cos γ cosϕ′) . (47)
Since N = E − p20/2, one can maximize the value of the translatory momentum p0 instead
of minimizing N(γ). Formula (39) represents the value γ = β lying at the boundary of the
accessible γ–domain; this value should be compared to p0(γ) taken at local maxima.
Let us consider the case α ≪ 1. At large enough energies, E ∼ 1, Eq. (47) is dominated
by the first term, which grows with γ, so that the maximum of p0(γ) is indeed achieved at
γ = β. At small energies, however, the second term in Eq. (47) becomes essential because of
the quickly oscillating cosϕ′ multiplier: the frequency of cosϕ′ oscillations grows as E → 0,
and at E ∼ α2, in spite of the small magnitude proportional to sinα, the second term
produces the sequence of local maxima of the function p0(γ).
One expects the parameters of the trajectory at small α not to be very different from the
ones at α = 0 (the latter case was considered in Sec. 3). So, we write,
γ = β − δγ ,
where 0 < δγ ≪ 1. Expanding the expressions (46), (47) and taking into account that
E ∼ α2 one obtains,
ϕ′ = −
2E sin β
(1 + δγ ctg β) , (48a)
2E(sin β − δγ cos β − α cos β cosϕ′) . (48b)
Now, the local maxima of the initial translatory momentum can be obtained explicitly by
differentiating Eqs. (48) with respect to δγ. One finds the sequence of them,
δγn = −tg β +
sin2 β
cos β
2πn− π − arcsin
2E sin2 β
α cos β
. (49)
Only the maxima with δγn > 0 should be taken into account. The local maxima exist when
E ≤ EB ≡
α2 cos2 β
2 sin4 β
. (50)
Substituting Eq. (49) into the expressions (48), one evaluates the values of p0 at the local
maxima,
p0,n(E) =2
2E sin β − 2E sin2 β
2πn− π − arcsin
2E sin2 β
α cos β
2E cos β
1− 2E sin
α2 cos2 β
The graphs Nn(E) = E − p20,n(E)/2 are shown in Fig. 6 for the case β = π/3, α = π/30.
Each graph is plotted for the energy range E > EAn restricted by the condition δγn > 0.
They are presented together with the curve given by the formula (40). By definition, the
critical solution corresponds to the lowest of these graphs. Clearly, for each “local” curve
representing the n-th local minimum of N(γ) there is a range of energies EAn < E < EBn
where it lies lower than the “global” curve (40). This means that the parameter γ of the
critical solution changes discontinuously across the points E = EBn . Correspondingly, the
curve Ncr(E) has a break at these points. On the other hand, the function Ncr(E) is smooth
at the points An as the “local” graphs end up exactly at δγ = 0, where the parameters of
the n-th “local” solution coincide with the ones of the “global” solution.
To summarize, we have observed that the boundary of the classically allowed region is
given by a collection of many branches of classical solutions, each branch being relevant in
its own energy interval. We will see that a similar branch structure is present in the complex
trajectories describing over–barrier reflections in the classically forbidden region of E, N .
4.3 Classically forbidden reflections
In this subsection we demonstrate that the suppression exponent F (E, N) viewed as a
function of energy at fixed N exhibits oscillations deep inside the classically forbidden region
0.02
0.04
0.06
0 0.05 0.1 0.15 0.2
Figure 6: The graphs Nn(E) corresponding to the local minima of the function N(γ) (dashed
lines) plotted together with the “global” curve, Eq. (40) (solid line); β = π/3, α = π/30.
The critical curve N = Ncr(E) is obtained by taking the minimum among all the graphs.
of initial data. This result comes without surprise if one takes into account the non-monotonic
behavior of the boundary Ncr(E) of the classically allowed region. Indeed, the curve N =
Ncr(E) coincides with the line F (E,N) = 0. One has,
= −∂EF
N=Ncr(E)
so that
(Ecrn , N
n ) = 0 .
We conclude that the points E = Ecrn are the local minima of the function F (E) at fixed
N = N crn . It is natural to expect that such local minima of F (E) exist at other values of
N as well. To illustrate this fact explicitly, we study the complex trajectories, solutions to
Eqs. (9), (10), (12).
Following the tactics of the previous section, we find solutions in three separate regions:
initial region x′ < 0, final region ξ > 0, and the intermediate region x′ > 0, ξ < 0. These
solutions, together with their first derivatives, should be glued at t = t0, when the complex
trajectory crosses the line x′ = 0, and at t = t1, when ξ = 0. Besides, we are looking for the
tunneling solution which ends up oscillating along the line AB, see Fig. 4. As discussed in
Sec. 3 this assumes existence of the second step of the process: classical decay of the unstable
orbit living at ξ = 0; the latter decay is described by a real trajectory14 going to x → −∞
at t→ +∞.
The solution in the final region ξ > 0 is (cf. Eqs. (30)),
ξ+(t) = 0 , (51a)
η+(t) =
cosα cos β
sin(t cosα cos β) , (51b)
where we used the time translation invariance (18) to fix the final oscillator phase ϕη = 0.
In the intermediate region x′ > 0, ξ < 0 one writes,
x′(t) = p′0t + x
0 , (52a)
y′(t) = a′e−it cosα + ā′eit cosα . (52b)
Note that the final solution (51) does not contain free parameters; thus, the matching of x′,
ẋ′, y′, ẏ′ at t = t1 enables one to express all the parameters in Eqs. (52) in terms of one
complex variable t1,
p′0 =
2E sin β cos φ1 , (53a)
x′0 = 1 +
[sin φ1 − φ1 cos φ1] , (53b)
eiφ1/ cos β [sinφ1 + i cos β cosφ1] , (53c)
ā′ =
e−iφ1/ cos β [sinφ1 − i cos β cos φ1] , (53d)
where we introduced φ1 = t1 cosα cos β.
As the energy of the solution has been fixed already, the only remaining initial condition
involves initial oscillator excitation number at x′ < 0, see Eqs. (10). It is convenient
to impose this condition at the matching point t = t0. One recalls the definition of the
matching time t0,
p′0t0 + x
0 = 0 ,
14One wonders why this trajectory does not reflect from the turn A′B′ on its way back. This concern is
removed by the observation that the trajectory produced in the decay of the unstable orbit is not unique:
in appendix we show that the decay can occur at any point of the segment AC giving rise to a whole bunch
of potential decay trajectories. Most of these trajectories pass through the turn A′B′ without reflection.
which, after taking into account the expressions (53a), (53b), leads to the following equation,
cosα√
2E sin β
sinφ1
cos β
− cosφ1∆φ = 0 , (54)
where ∆φ = cosα(t1 − t0). At t = t0 one has,
ẋ(t0) = p
0 cosα− ẏ′(t0) sinα =
2(E −N) ,
and thus √
= ctgα sin β cosφ1 − sin φ1 sin∆φ− cos β cosφ1 cos∆φ . (55)
As before, ν = N/E.
Two complex equations (54), (55) determine the matching times t0, t1, and, consequently,
the complex trajectory. Although these equations cannot be solved explicitly, they can
be simplified in the case α ≪ 1, which we consider from now on. For concreteness, we
study reflections at N = 0. It is important to keep in mind that in the region of interest
E ∼ Ecrn ∼ O(α2); thus, one should regard all the momenta p and oscillator amplitudes a,
ā, as the quantities of order O(α). At the same time, for the distances along the waveguide
one has x ∼ O(1), so that the real parts of time intervals may be parametrically large,
Re t ∼ x/p ∼ O(1/α).
Further on, it will be convenient to work in terms of real variables, so, we represent φ1
and ∆φ as
φ1 = cosα cos β(τ1 + iT1) , ∆φ = cosα(τ + i∆T ) .
Note that τ and ∆T are the real and imaginary parts of the time interval t1 − t0 which the
particle spends in the intermediate region. Now, equation (54) enables one to express
2E sin βch(T1 cos β)
+O(α) , (56)
τ1 = −
τ cos β
cos β
−∆T cth(T1 cos β)
+O(α3) . (57)
Note that τ1 ∼ O(α), τ ∼ O(1/α). Then, the real part of Eq. (55) implies that
ch(T1 cos β) =
sin β
1 + α ctgβ cos τe∆T
+O(α2) . (58)
While deriving this formula we imposed T1 < 0 which follows from the requirement that
in the limit α → 0 equation (31) should be recovered; besides, we assumed e∆T ∼ O(1).
Substituting Eq. (58) into Eq. (56) and the imaginary part of Eq. (55), we obtain the final
set of equations,
2E = α ctgβ cos τe∆T +O(α2) , (59a)
(1 + ∆T )e−∆T = α ctgβτ sin τ +O(α) . (59b)
These two nonlinear equations, still, cannot be solved explicitly. Nevertheless, one can get
a pretty accurate idea about the structure of their solutions.
Before proceeding to the analysis of the above equations, let us derive a convenient
expression for the suppression exponent F0(E) ≡ F (E,N = 0). Note that on general
grounds one expects to obtain an expression of the form,
F0(E) = E(fβ(0) +O(α)) ,
where fβ(0) is given by Eq. (33). We are interested in the O(α) correction in this expression,
so, one must be careful to keep track of the subleading terms during the derivation.
Making use of the equations of motion, one obtains for the incomplete action (16) of the
system,
2 Im S̃ = Im p′0 =
2E sin β Im(cosφ1) .
Substitution of Eqs. (56), (57), (58) into this formula yields
2 Im S̃ = 2E
−1−∆T − α ctg β cos τe∆T
cos2 β
+ 2∆T
+O(α2)
For the parameter T one has (see Eqs. (14)),
T = −
2 Im x0
2 Im(x(t0)− p0t0)
= 2(T1 −∆T ) +
sinα Im y′(t0) , (60)
where in the last equality we used Eqs. (34) and x′(t0) = 0. The quantity Im y
′(t0) is
evaluated by using Eqs. (52b), (53) and (58); one finds,
Im y′(t0) = −
ctg β cos τe∆T +O(α)
Substituting everything into the formula (15), we obtain,
F0(E) = E
fβ(0)− 4α ctg β cos τ ∆T e∆T +O(α2)
. (61)
This expression implies that determination of the O(α) correction to the suppression expo-
nent involves finding τ , ∆T with O(1)–accuracy. This is precisely the level of accuracy of
Eqs. (59). Below we will also need the following formulae, which can be easily obtained by
using T = −F.
and Eq. (60),
= fβ(0) + 2(∆T + 1) +O(α) , (62)
2(∆T + 1 +O(α))
. (63)
Note that, though the suppression exponent differs from that in the one–turn case only by
O(α) correction, its derivative gets modified in the zeroth order in α.
Now, we are ready to analyze Eqs. (59). One begins by solving Eq. (59b) graphically,
see Fig. 7. The important property of this equation is as follows. One notices that the l.h.s.
of Eq. (59b) is always smaller than 1, the maximum being achieved at ∆T = 0. Therefore,
the solutions to this equation are confined to the bands
τ sin τ <
This corresponds to
τ ∈ [0; 2π(n1 − 1) + δτn1 ] or τ ∈ [2πn− π − δτn; 2πn + δτn] , n ≥ n1 (64)
where
δτn = arcsin
2πα(n− 1/2)
+O(α) ,
+ 1 , (65)
with [·] in the last formula standing for the integer part. The forbidden bands, where
τ sin τ > tgβ/α, are marked in Fig. 7 by yellow shading. The property (64) introduces a
topological classification of the solutions τ , ∆T to Eqs. (59). Namely, these solutions fall
into a set of continuous branches: the “local” branches τn(E), ∆Tn(E) living inside the
strips τ ∈ [2πn − π − δτn; 2πn + δτn], n ≥ n1, and the “global” branch τg(E), ∆Tg(E)
inhabiting the very first band τ ∈ [0; 2π(n1 − 1) + δτn1 ]. As follows from the definition of
τ , the topological number n counts the number of y′–oscillations during the evolution in the
intermediate region.
Let us consider the “global” branch. From Eqs. (59) one has,
τg → 2π(n1 − 1) +O(α lnα) , ∆Tg → ln(tg β/α) , E → 0 ,
τg → 0 , ∆Tg → −1 , E → +∞ .
10π 9π 8π 7π 6π 5π4π 3π 2π π 0
g 4 5
10π 9π 8π 7π 6π 5π4π 3π 2π π 0
g 4 5
Figure 7: Curves representing solutions to Eq. (59b); β = π/3, α = π/30.
By inspection of Fig. 7 one can work out the qualitative behavior of the functions τg(E),
∆Tg(E). Alternatively, these functions can be found numerically. They are plotted in Fig. 8
for the case β = π/3, α = π/30 (the curves marked with “g”). One observes that at
high enough energies the function ∆Tg(E) exhibits oscillations around the line ∆T = −1.
According to the formula (63) this means that the function F0(E)/E is non-monotonic, it
attains local minima at the points
E ′n =
8π2(n− 1/2)2
1 + 2αe−1ctgβ +O(α2)
. (66)
Moreover, if
n ≥ n′0 ≡
fβ(0) exp
fβ(0)
+ 1 (67)
there exist Eon = E
n(1 + O(α)), such that ∆T (E
n) = −1 − fβ(0)/2. Then, according to
Eq. (62) the points Eon are the “optimal” energies corresponding to the local minima of the
suppression exponent F0(E).
At low energies the function ∆Tg(E) ceases to oscillate and becomes large and positive.
According to Eq. (62) this means that the suppression exponent F0,g(E) of the “global”
solution becomes negative at low energies15, see Fig. 9. This is a clear signal that the
15It is worth mentioning that Eqs. (59) and the expression (61) for the suppression exponent become
0 0.1 0.2 0.3
0 0.1 0.2 0.3
0 0.1 0.2 0.3
E’4E’5
Figure 8: Several first branches of solutions to Eqs. (59): “global” branch (“g”) and two
“local” branches (“4”, “5”); β = π/3, α = π/30.
0.05
0.15
0.25
0 0.2 0.4 0.6
0.02
0.04
0.120.110.100.09
Figure 9: The suppression exponent F0(E) for the “global” and first “local” (n = 4) branches;
β = π/3, α = π/30. The vicinity of intersection of the graphs is enlarged in the upper right
corner.
“global” solution becomes unphysical at these energies and its contribution to the reflection
probability should be discarded: negative suppression exponent contradicts the unitarity
requirement16, P < 1. One is forced to conclude that at low energies reflection is described
by the “local” solutions. Let us study them in detail.
For the n-th branch one obtains,
τn → 2πn+O(α lnα) , ∆Tn → ln(tg β/α) , E → 0 ,
τn → 2πn− π , ∆Tg → +∞ , E → +∞ .
From Fig. 7 one learns that the n-th solution passes through the points
∆Tn = −1 , τ = 2πn or τ = 2πn− π . (68)
inapplicable at large ∆T : the assumption e∆T ∼ O(1) which was used in the derivation of these equations
gets violated. Nevertheless, by analyzing the full equations (54), (55) one can show that dF0,g/dE = −Tg
is large and positive at E → 0. This is sufficient for concluding that F0,g(E) is negative in the low–energy
domain.
16Another indication that the “global” solution is unphysical at small E is that the function τg(E) is
bounded from above. Indeed, τ is the time interval the particle spends in the intermediate part of the
waveguide, one expects it to tend to infinity as E → 0 for a physically relevant solution.
Thus, each curve ∆Tn(E) has one sharp dip, its minimum is smaller than −1, see Fig. 8. As
in the case with the “global” branch, the points (68) represent the extrema of the functions
F0,n(E)/E; the positions of the local minima are again given by Eq. (66).
Making use of Eq. (61), we find that the suppressions F0,n(E) of the “local” branches are
large and positive at high energies. Hence, these solutions give subdominant contributions
to the reflection probability at such E as compared to the “global” solution. As energy
decreases, F0,n(E) also decreases, then makes one oscillation and drops to negative values at
small E. The latter property means that each “local” branch becomes unphysical at small
enough energies. The suppression exponent of the first “local” branch (corresponding to
n = 4 in the case β = π/3, α = π/30) is presented in Fig. 9.
An alert reader may have already guessed that we have met here the typical Stokes
phenomenon [21]. In fact, the Stokes phenomenon is specific to the situations where some
integral (e.g., the path integral (7) in our case) is evaluated by the saddle–point method.
Essentially, it means the following: as one gradually changes the parameters of the integral
in question, a given saddle point may become non–contributing after the values of these
parameters cross a certain curve drawn in the parameter space, the Stokes line. Since the
result of the computation should be continuous, this phenomenon occurs only for subdomi-
nant saddle points (saddle–point trajectories in our case). Unfortunately, apart from several
heuristic conjectures [21, 12], sometimes rather suggestive [13], there is presently no general
method of dealing with the Stokes phenomenon in the semiclassical calculations. However,
in the situation encountered above it suffices to use the simplest logic lying at the heart of
all other approaches17.
When gathering the final result for the suppression exponent, we follow two guidelines.
First, it is clear that, as energy decreases, each branch becomes unphysical before F0,n(E)
crosses zero. On the other hand, at high energies one should pick up the branch corresponding
to the smallest value of the suppression exponent. Looking at Fig. 9, one notes that the
curves F0,g(E), F0,4(E) have two intersections, A and B. At E > EB one chooses the
“global” branch. In the region EA < E < EB we switch to the first “local” branch, because
in this region F0,4(E) < F0,g(E). Naively, at E = EA one should jump back to the “global”
branch; however, in order to preserve unitarity at small energies, we suppose that somewhere
in between the points B and A the “global” branch becomes non–contributing, so that one
should stay at the “local” branch at E < EA. Similarly, the adjacent “local” branches have
17The simplification in the present case is related to the fact that we concentrate on the dominant semi-
classical contribution, leaving aside the subdominant ones.
two intersections; as the energy decreases, we switch from n-th branch to n + 1-th at the
first intersection, and stay there until the intersection with the n+2-th branch. Overall, one
obtains the graph for the suppression exponent plotted in Fig. 10. The suppression exponent
0.02
0.04
0.06
0.08
0 0.05 0.1 0.15 0.2 0.25
Figure 10: The final result for the suppression exponent F0(E) in the region of small energies;
β = π/3, α = π/30. The points where different branches merge are shown with thick black
dots.
oscillates between two linear envelopes, F = E(fβ(0) ± 4e−1α ctg β); oscillations pile up in
the region of low energies. The reflection process is optimal in the vicinities of the minima
of the function F0(E).
5 Discussion
By considering a class of two–dimensional waveguide models, we have demonstrated explicitly
that the probability of over–barrier reflection can be non–monotonic function of energy. The
origin of the effect lies in the classical dynamics: the parameters of the complex trajectory
describing over–barrier reflection change quasi-periodically as the energy gets decreased.
This results in the oscillatory behavior of the suppression exponent. Reflection occurs with
exponentially larger probability in the vicinities of “optimal” energies (local minima of the
suppression exponent) while being highly suppressed in between.
Our results are obtained for a fairly specific class of waveguides, namely, the ones with
very sharp turns. However, the qualitative features observed in this paper should be valid
for quite general waveguide models: a classical particle with high energy feels any large–scale
turn of the waveguide as a sharp one18; if two turns are separated by a long interval of free
motion, one arrives to the model (35). We remark that the phenomenon of optimal tunneling
has been observed also in numerical investigation of a smooth waveguide, see Ref. [14].
The branch structure of solutions observed in the region of small energies is interesting
from the mathematical point of view. We have shown that there exists an infinite sequence of
complex trajectories marked by the topological number n. Each branch produces physically
consistent result for the suppression exponent in some energy interval; outside of this interval
the n-th branch would correspond either to highly suppressed transitions (high energies) or
to violation of unitarity (low energies). We collected the final graph for the suppression
exponent basing on the empirical considerations, which hardly may be acknowledged as
satisfactory. Our study clearly shows that the method of complex trajectories should be
equipped with a convenient rule to pick up the physical trajectory among the discrete set of
solutions to the boundary value problem (9), (10), (12) (in other words, the method to deal
with the Stokes phenomenon). Presently, such a rule is absent.
We note that the described physical phenomenon of optimal tunneling is present inde-
pendently of the way the branches of solutions are glued together. The result at relatively
high energies is given by the “global” branch, which displays a large number of local minima
if n′0 > n1, see Eqs. (67), (65). This is the case for the illustrative example considered
throughout this paper, see Fig. 9.
As a final remark, we point out some open issues. We have calculated the suppression
exponent of reflection using the sharp–turn approximation. It would be instructive to extend
our analysis by finding corrections due to the finite turn widths. The motivation is twofold.
First, the analysis performed in appendix implies existence of a rich variety of distinctive
semiclassical solutions contributing almost equally into the reflection probability. This fea-
ture might be a manifestation of chaos [7] which is present in our system but hidden by the
sharp–turn approximation. [Note that chaos is inherent in a very similar waveguide model
18More precisely, one should compare the width b of the turn to the quantity 2π
, where p0 is the
translatory momentum of the particle and ω stands for the frequency of transverse oscillations; if b≪ 2πp0
one is in the class of models with sharp turns.
with smooth potential, see Ref. [14].] Clearly, the structure of solutions in the vicinities of
the turns is worth further investigation.
Second, it was proposed recently in Refs. [9, 11] that the process of dynamical tunneling
in quantum systems with multiple degrees of freedom (including field theoretical models,
see Refs. [19]) can proceed differently from the ordinary case of one–dimensional tunnel-
ing. Namely, classically unstable state can be created during the process; this state decays
subsequently into the final asymptotic region. The analysis performed in the present paper
naturally conforms with this tunneling mechanism: all our complex trajectories are matched
with the unstable orbit living at the turn. Still, the sharp–turn approximation does not allow
to distinguish between the truly unstable trajectories staying at the turn forever and those
which reflect from the turn in a finite time. To decide whether the tunneling mechanism of
Refs. [9, 11] is indeed realized in our model one needs to go beyond the sharp–turn approx-
imation. Then, the candidate for the “mediator” unstable state is the “excited sphaleron”,
the solution considered in the appendix. Presumably, in our model one can answer analyti-
cally to the question of whether or not the “excited sphaleron” acts as an intermediate state
of the tunneling process. This study is quite beyond the scope of the present paper and we
leave it for future investigations.
Acknowledgments. We are indebted to F.L. Bezrukov and V.A. Rubakov for the en-
couraging interest and helpful suggestions. This work is supported in part by the Russian
Foundation for Basic Research, grant 05-02-17363-a; Grants of the President of Russian
Federation NS-7293.2006.2 (government contract 02.445.11.7370), MK-2563.2006.2 (D.L.),
MK-2205.2005.2 (S.S.); Grants of the Russian Science Support Foundation (D.L. and S.S.);
the personal fellowship of the “Dynasty” foundation (awarded by the Scientific board of
ICFPM) (A.P.) and INTAS grant YS 03-55-2362 (D.L.). D.L. is grateful to Universite Libre
de Bruxelles and EPFL (Lausanne) for hospitality during his visits.
A Classical motion near the turn
In this appendix we analyze the motion of the particle near the sharp turn of the waveguide
(20) at nonzero smoothening of the turn, see, e.g., Eq. (21). We suppose that in the small
vicinity of the turn the function w(ξ, η) can be represented in the form
w(ξ, η) = cos β (η − bv(ξ/b)) , (69)
where v(ψ) does not depend explicitly on b. Moreover, we consider the case when v(ψ) has
a maximum19,
v′(ψ0) = 0 . (70)
Due to the property (70) one immediately obtains the exact periodic solution to the
equations of motion (9), which we call “excited sphaleron” [9],
ξsp = bψ0 , ηsp = Aη sin(t cos β + ϕη) + bv(ψ0) . (71)
We are going to show that this solution is unstable: a small perturbation above it grows
with time and the particle flies away to either end of the waveguide. In particular, there are
solutions that describe the decay of the sphaleron to ξ → −∞ both at t → ±∞. Clearly,
such solutions correspond to reflections from the turn.
In the vicinity of the sphaleron the trajectory of the particle can be represented in the
form,
ξ = bψ(t) , η = ηsp(t) + bρ(t) , (72)
where ψ, ρ ∼ O(1). Writing down the classical equations of motion (9) in the leading order
in b, one obtains,
Aη sin(2s)v
′(ψ) , (73)
+ 4ρ = 4[v(ψ)− v(ψ0)] , (74)
where s = (t cos β +ϕη)/2. It is worth noting that the right hand side of Eqs. (73), (74) are
of different order in b. We will see that due to this difference ρ = 0 in the leading order in b.
Let us first consider the linear perturbations above the excited sphaleron,
ψ = ψ0 + δψ , δψ ≪ 1 .
Equation (73) can be linearized with respect to δψ leading to the Mathieu equation
δψ + 2q sin(2s)δψ = 0 ,
with canonical parameter q = −2v′′0Aη/b > 0. As q ∼ O(1/b) ≫ 1, one can apply the WKB
formula,
A cosW
dW/ds
, (75)
19For the smoothening (21), the properties (69), (70) hold with v(ψ) = ψtgβ
, ψ0 ≈ 1.28.
where |A| ≪ 1, and
sin(2s′) .
Note that we have chosen the solution symmetric with respect to time reflections,
δψ(π/2− s) = δψ(s) . (76)
At s ∈ [0; π/2] the exponent W is real and the particle gets stuck at ψ ≈ ψ0, oscillating
around this point with high frequency dW/ds ∼ O(b−1/2). At s < 0 the solution (75) grows
exponentially, meaning that the particle flies away from the excited sphaleron,
δψ(s < 0) =
A cos(W (0)− π/4)
|dW/ds|
e|W (s)−W (0)| .
In what follows, we choose A cos(W (0)− π/4) < 0, so that δψ < 0 at s < 0. Let us denote
by s1 < 0 the point where δψ becomes formally equal to −1,
A cos(W (0)− π/4)
|dW/ds|
e|W (s1)−W (0)| = −1 .
In what follows we suppose that s1 ∼ O(1), hence, A is exponentially small. Then, in the
vicinity of this point, |s− s1| ≪ 1, one has,
δψ = − exp
−2q sin(2s1)(s1 − s)
= − exp
4v′′0Aη sin(2s1)
(s1 − s)√
. (77)
We notice that δψ evolves from exponentially small values to δψ ∼ O(1) during the charac-
teristic time |s− s1| ∼ O(
When δψ ∼ O(1) the linear approximation breaks down and one has to solve the nonlinear
equation (73). Using s = s1 +O(
b) one writes
Aη sin(2s1)v
′(ψ) . (78)
This equation permits to draw a useful analogy with one–dimensional particle moving in
the effective potential Veff (ψ) = −4b−1Aη sin(2s1)v(ψ) (see Fig. 11). This auxiliary particle
starts in the region near the maximum of the potential at (s − s1)/
b → +∞ with energy
E ≈ Vmax and rolls down toward ψ → −∞ at (s−s1)/
b→ −∞. In this limit v(ψ) → ψ tg β
and the solution takes the form
ψ = C1 + C2(s− s1) + 2b−1Aη sin(2s1) tg β (s− s1)2 .
Vmax
Figure 11: The effective potential for Eq. (78).
Note that the coefficients C1, C2 here are not independent: they are determined by the
parameter s1 through matching of the solution with Eq. (77) at (s− s1)/
b→ +∞. We do
not need their explicit form, however.
Let us argue that the function ρ remains small during the whole evolution of the particle
in the vicinity of the sphaleron. Indeed, in the linear regime one has δψ ≪ 1 and the r.h.s.
of Eq. (74) is small. So, ρ does not get excited. On the other hand, the nonlinear evolution
of ψ proceeds in a short time interval ∆s = O(
b); so, again, ρ is suppressed by some power
of b.
The trajectory (72) found in the vicinity of the sphaleron should be matched at
1 ≫ |s− s1| ≫
with the free solution in the asymptotic region ξ < 0, see Eqs. (26). It is straightforward to
check that matching can be performed up to the second order in (t− t1), which is consistent
with our approximations. In this way one determines the free asymptotic solution which, up
to corrections of order O(b), coincides with the sinusoid coming from ξ → −∞ at t → −∞
and touching the line ξ = 0 at t = t1.
Now we recall that, by construction, the obtained solution is symmetric with respect to
time reflections,
ξ(s) = ξ(π/2− s) , η(s) = η(π/2− s) .
This means that it satisfies ξ → −∞ at t → ±∞. This solution describes reflection of the
particle from the turn.
The reasoning presented in this appendix puts considerations of the main body of this
paper on the firm ground: we have found the “smoothened” solutions which reflect classically
from the turn, and in the limit b→ 0 coincide with the free solutions of Sec. 3 touching the
line ξ = 0.
It is worth mentioning that, apart from the reflected solution we have found, in the
vicinity of any trajectory touching the line ξ = 0 there exists a rich variety of qualitatively
different motions. First of all, one may successfully search for solutions which are odd with
respect to time reflections (Eq. (76) with minus sign). Such solutions, though close to the
reflected ones at t < 0, describe transmissions of the particle through the sharp turn into the
asymptotic region ξ → +∞. Relaxing the time reflection symmetry, one can find solutions
leaving the vicinity of the turn at any point η < 0, which is different, in general, from the
starting point η = η(s1). Yet another types of solutions are obtained in the case when the
amplitude A of δψ–oscillations at s ∈ [0; π/2] is so small that δψ does not reach the values
of order one during the time period s ∈ [−π/2; 0]. If the particle is still in the vicinity of the
point ψ0 at s = −π/2, it remains for sure in this vicinity at s ∈ [−π; −π/2], because the r.h.s.
of Eq. (73) is positive again. In this way one obtains solutions, which spend two, three, etc.
sphaleron periods at ψ ≈ ψ0 before escaping into the asymptotic regions ψ → ±∞. In the
leading order in b all these solutions correspond to the identical initial state, and (in the case
of classically forbidden transitions) to the same value of the suppression exponent. However,
an accurate study of the dynamics in the vicinity of the the sphaleron is generically required
to obtain the correct value of the suppression exponent in the case b ∼ 1, cf. Ref. [14].
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Introduction
The semiclassical method
The model with one turn
The model with two turns
Introducing the system
Classical evolution
Classically forbidden reflections
Discussion
Classical motion near the turn
|
0704.0411 | Molecular circuits based on graphene nano-ribbon junctions | Microsoft Word - APL-Ribbon_junctions_20070328-Text_with_Figures-ARXIV.doc
Molecular circuits based on graphene nano-ribbon junctions
Zhiping Xu†
Department of Engineering Mechanics, Tsinghua University, Beijing, 100084, China
Graphene nano-ribbons junctions based electronic devices are proposed in this Letter. Non-
equilibrium Green’s function calculations show that nano-ribbon junctions tailored from
single layer graphene with different edge shape and width can act as metal/semiconductor
junctions and quantum dots can be implemented. In virtue of the possibilities of patterning
monolayer graphene down to atomic precision, these structures, quite different from the
previously reported two-dimensional bulk graphene or carbon nanotube devices, are
expected to be used as the building blocks of the future nano-electronics.
Keyword: graphene nano-ribbon, electronic transport, metal/semiconductor junction, quantum dot
† Email: [email protected]
Nano-electronics, or molecular electronics have been proposed as the alternative to silicon in
future technical applications1 and have attracted great interests recently. In virtue of their unique
structures and various functions, these nanostructures possess intriguing electromagnetic,
mechanical and optical features. Especially, carbon based nanostructure, such as fullerene, graphene
and carbon nanotubes are the most interesting structures because of their rich variety of excellent
physical properties. For instance, anomalous quantum hall effects (QHE) and massless Dirac
electronic behavior have been discovered in the graphene systems2, 3, and these discovery has
sparked lots of investigations on this unique two-dimensional material. Tailored from monolayer
graphene, graphene ribbon (GNR) with finite width has been shown to hold unusual electronic
properties4, depending on their edge shape and width. In more details, ribbons with zigzag edges
(ZGNRs) possess spin-polarized peculiar edges states and spin-polarized electronic state provides
half-metallicity under transverse electric field and has great potential in the application as
spintronics5. In contrast, the armchair edged ribbons (AGNR) can be either metallic or
semiconducting depending on their width6, AGNR with width Na (named as NaAGNR in the
conventional nomenclature) has been shown to be metallic only if Na = 3k + 2 and semiconducting
otherwise, where k is an integer.
From the experimental point of view, the fascinating feature of the ribbons is that the
graphene material can be easily patterned using standard micro- or nano-electronics lithography
methods. Unlike the carbon nanotubes or other low-dimensional nanostructures, the GNRs with
intricate sub-micrometer structures can now be fabricated7, 8, 9, and it is believed that a combination
of standard lithographic and chemical methods will help to pattern the graphene with atomic
precision down to the molecular level. The high mobility μ = 2.7 m2/V.s, large elastic mean free
path le = 600nm and phase coherence lengths lφ= 1.1 μm observed7 in the epitaxial graphene
patterned suggest the use of pure GNR structures as the building blocks of the nanoscale confined
and coherent electronic circuits. To realize the components such as field transistors9 and coulomb
blockade devices, experimentally controllable metal/semiconductor junctions and quantum dots will
be essential. As proposed by Chico et al.10, 11, these can be achieved by jointing different carbon
nanotubes. However, the fabrication and control of the nanostructure of graphene ribbons are much
more convenient than introducing pentagon-heptagon defects in carbon nanotubes as discussed,
therefore it is interesting to investigate the possibilities of ribbon junction based nano-circuits.
To this end, we have proposed several kinds of the GNRs based electronic devices in this
Letter. We show that, by controlling the tailoring process of GNRs with different edge shape and
width, the metal/semiconductor junctions and quantum dots can be easily implemented
experimentally. To validate this, electronic transport calculation using the non-equilibrium Green’s
function method have been carried out following Landauer’s approach12. The electronic structure of
the graphene lattice is described using the nearest-neighbour π-orbital tight-binding model and the
hopping parameter Vppπ = 2.75 eV is used. This simple topological model gives quantitative results
comparing with the LDA results except for the gap opening at small width as the consequence of
the length changing of σ bonds6. By solving the Green’s function, the conductance was finally
calculated as G = G0Tr[ΓLGRΓRGA] and the density of state is expressed as D = –ImTr[GR]/π11,
where G0 = 2e2/h is the unit quanta of conductance including the spin degeneracy, GR(A) is the
retarded (advanced) Green’s function of the conductor and ΓL(R) is the spectral density describing
the coupling between the left (right) lead and the conductor. In our model, the leads are represented
using semi-infinite graphene ribbons attached to the conductor region, with the same shape and
width.
First of all we investigate straight metal/semiconducting junction 11AGNR/10AGNR. The
structure of the junction is considered by simply patching two different straight ribbons together,
leave a width mismatching at the interface. The result shown in Fig. 1 indicates a gap Eg = 0.93 eV
near the Fermi energy and the imperfection at the interface induces a deviation of conductance from
the step-like curve of the perfect ribbon. However, the van Hove singularities which are the
characteristics of 1D system remain.
To examine the detailed electronic structure of the junction, a spatial-resolved localized
density of states (LDOS) analysis is helpful. We have grouped the atoms into slices according to
their distance from the interface. Each 4.26 Å long slice (a unit cell of the perfect AGNR) in the
10AGNR, 11AGNR and interface part contain 20, 22 and 21 atoms respectively. The LDOS
averaged at different slices are plotted in Fig. 1. From the semiconducting 10AGNR side we find
the LDOS is distorted near the interface and gap state appears through the contact with metallic
11AGNR. However at slices far from the interface, at slice 3 for example, the perfect
semiconduting behavior is mostly recovered. At the scattering interface the van Hove singularities
have been smoothed and 1D metallic structure gradually emerges as the distance from interface
increases from the 11AGNR side. The arising of gap state near the interface characterizes the metal-
semiconductor junction and suggests the possibilities of building Schottky devices.
Furthermore, L-shape GNR junctions with different orientations can be constructed. For
instance, the LDOS of 8ZGNR/15AGNR junction with a π/6 joint is analysized in Fig. 2. As
expected, the edge state of the 8ZGNR spreads into the semiconducting 15AGNR side. Because of
the ZGNR possess spin-polarized structure, so this half-metal/semiconducting junction inspires
interests in the spin-transport devices.
Beside of the metal/semiconductor junction, the semiconductor/semiconductor junctions
have also been investigated and defect states in gap appear at the interface. Moreover in the
ZGNR/ZGNR junctions, zero-conductance dips13 near Fermi energy have been observed, caused by
the complete backward scattering.
The metal-semiconductor junction also suggests quantum dot devices through combing two
of them together. We now consider the junction 12AGNR/11AGNR/12AGNR. In this structure a
central metallic ribbon is sandwiched by two semiconductor barriers where quantized states can be
formed. Our calculation results depicted in Fig. 3 show two sharp DOS peaks inside the gap of
semiconducting 12AGNR containing 7 unit cells, with energy E1,2 = 0.2025 and -0.2025 eV. As
seen from the spatial-resolved LDOS at E = 0.2025, the bounded state is localized inside the
11AGNR region. The structure of the quantum levels can be further tuned by changing the length of
11AGNR. From our calculation, as it changes from 1 to 8 unit cells, the energy spacing between the
nearest peaks around Fermi energy, i.e. ΔE = E1-E2, gradually decreases from 0.785 eV to 0.385 eV
and their DOS becomes higher and sharper.
We have also observed quantized edge states within the 10AGNR/7ZGNR/10AGNR
junctions through introducing two π/6 joints. The results are shown in Fig. 3 where we can found 7
LDOS peaks inside the zero-conductance gap. The quantized states with E = -0.3525, -0.1625, -
0.05, 0, 0.05, 0.1625 and 0.3525 correspond to different LDOS patterns (see Fig. 4 for E = 0.3525).
The higher the energy, the more nodes of the bounded standing wave have. The electron wave
quantized pattern depends on the structure of the central region.
In conclusion, we have proposed nano-electronic circuits based on graphene nano-ribbon
junctions. Through tailoring GNRs into junctions of different edge shape and width, we can
implement metal/semiconductor junctions and quantum dots in principle. In virtue of the possibility
of molecular level patterning based on lithography and chemical methods, these devices are
expected to be fabricated easier in comparison with other structures such as the single molecule or
carbon nanotubes junctions, and are expected to find great applications in the large-scale integrated
nano-circuits in future.
The work is supported by the National Science Foundation of China through Grants
10172051, 10252001, and 10332020 and the Hong Kong Research Grant Council (NSFC/RGC N
HKU 764/05 and HKU 7012/04P). ZX also thanks Prof. Wenhui Duan, Dr. Tao Zhou and Dr.
Haiyun Qian from the Department of Physics in Tsinghua University for their help on the
calculation.
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FIG. 1. The metal/semiconducting junction 11AGNR/10AGNR: (Top) Conductance and DOS of
the whole system. (Bottom) LDOS at slices near the interface. Slice n (n = 1, 2 and 3) represents the
n-th nearest slice to the interface and the vertical scale of DOS is 0.2.
FIG. 2. Spatial-resolved LDOS in metal/semiconducting junction 8ZGNR/15AGNR, the vertical
scale of DOS is 0.2.
FIG. 3. Quantum dot structure based on 12AGNR/11AGNR/12AGNR junction: (Top) Conductance
and DOS at low bias, where two isolated sharp peaks appear inside the gap; (Bottom) Spatial-
resolved LDOS at E = 0.2025 eV, the grey dot represents the ionic site and the radius of circle
around it corresponds to the value of LDOS.
FIG. 4. Quantum dot structure based on 10AGNR/7ZGNR/10AGNR junction: (Top) Conductance
and DOS; (Bottom) Spatial-resolved LDOS at E = 0.3525 eV.
Figure 1
Figure 2
Figure 3
Figure 4
|
0704.0412 | Unit groups of integral finite group rings with no noncyclic abelian
finite subgroups | UNIT GROUPS OF INTEGRAL FINITE GROUP RINGS WITH
NO NONCYCLIC ABELIAN FINITE SUBGROUPS
MARTIN HERTWECK
Abstract. It is shown that in the units of augmentation one of an integral
group ring ZG of a finite group G, a noncyclic subgroup of order p2, for some
odd prime p, exists only if such a subgroup exists in G. The corresponding
statement for p = 2 holds by the Brauer–Suzuki theorem, as recently observed
by W. Kimmerle.
1. Introduction
Is a finite subgroup H of units in the integral group ring ZG of a finite group
G necessarily isomorphic to a subgroup of G? Of course, torsion coming from the
coefficient ring should be excluded, that is, only finite subgroups H in V(ZG), the
group of units of augmentation one in ZG, will be considered. The question was
raised by Higman in his thesis (1940), where he gave an affirmative answer when G
is metabelian nilpotent or the affine group over a prime field; cf. Sandling (1981).
In the survey of Sandling (1984) it is included as Problem 5.4, and noted that an
affirmative answer for metabelian G was finally given by Roggenkamp (1981); but
see also Cliff, Sehgal and Weiss (1981), and Marciniak and Sehgal (2003) for a more
recent result, giving a generalization based on a theorem of Weiss (1988). These
results are really about certain ‘large’ torsion-free normal subgroups of V(ZG). For
a more complete discussion, see Chapter 4 in Sehgal’s book (1993).
As a sort of converse, one may fix a finite group H and look for groups G for
which H embeds into V(ZG), again hoping for the best, but little is known in this
respect. What is known is that if a cyclic group H of prime power order embeds
into some unit group V(ZG), then H also embeds into G (due to an observation
of Cohn and Livingstone (1965); see also Zassenhaus (1974)), and only recently in
Hertweck (2007b) it was shown that the restriction on the order can be removed
if in addition G is assumed to be solvable. In this spirit, Marciniak, at a satellite
conference of the ICM 2006, asked whether a group G necessarily has a subgroup
isomorphic to Klein’s four group provided this is the case for V(ZG). Kimmerle
immediately observed that this is implied by the Brauer–Suzuki theorem (rendered
in Kimmerle (2006)), see Section 2. Our complementary result is as follows.
Theorem A. Let G be a finite group. Suppose that V(ZG) has a noncyclic abelian
subgroup of order p2, for some odd prime p. Then the same is true for G (i.e.,
Sylow p-subgroups of G are not cyclic).
Date: October 30, 2018.
2000 Mathematics Subject Classification. Primary 16S34, 16U60; Secondary 20C05.
Key words and phrases. integral group ring, torsion unit, partial augmentation.
http://arxiv.org/abs/0704.0412v1
2 MARTIN HERTWECK
It is easy to verify that a finite p-group with no noncyclic abelian subgroup is
either cyclic or a (generalized) quaternion group, see Theorem 4.10 in Gorenstein
(1968). It comes to mind that the theory of cyclic blocks might be used in the
proof, but it is pretty simple and makes only use of a fact about vanishing of
partial augmentations of torsion units, established in Hertweck (2006, 2007a).
We remark that both results (whether p is even or odd) for a solvable group G
are covered by Theorem 5.1 in Dokuchaev and Juriaans (1996).
Note that a group G whose Sylow 2-subgroups are cyclic has a normal 2-comple-
ment, by Burnside’s well known criterion, see Theorem 4.3 in Gorenstein (1968).
We obtain the following corollary.
Corollary 1. Let G be a finite group having cyclic Sylow p-subgroups for some
prime p. Then any finite p-subgroup of V(ZG) is isomorphic to a subgroup of G.
Finally, we remark that, as with other results in this field, the theorem can be
formulated for more general coefficient rings than Z, notably for the semilocalization
of Z at the prime divisors of the order of G. Unfortunately, it is definitely wrong
for p-adic coefficient rings.
2. Kimmerle’s observation
Coming back to the initial question, we mention that in the hope for further
positive results, it is natural to impose restrictions on the prime divisors of the
finite subgroup H , i.e., to consider only π-groups H for some set π of primes (a
singleton {p}, to begin with), as has been done before in work on the stronger
Zassenhaus conjecture (ZC3), cf. Dokuchaev and Juriaans (1996). It is well known
that then, one can assume that Oπ′(G), the largest normal π
′-subgroup of G, is
trivial, for H has an isomorphic image under the natural map ZG → ZG/Oπ′(G),
see the remark after Theorem 2.2 in Dokuchaev and Juriaans (1996).
This derives from the vanishing of certain partial augmentations of the elements
of H . Recall that for a group ring element u =
g∈G agg (all ag in Z), its partial
augmentation with respect to an element x of G, or rather its conjugacy class xG
in G, is the sum
g∈xG ag; we will denote it by εx(u). The result of Cohn and
Livingstone mentioned in the introduction really says that if an element h of H is
of prime power order, then there exists an element x in G of the same order such
that εx(h) 6= 0. Note that εz(u) = az for an element z in the center of G. An
old yet fundamental result from Berman (1955) and Higman (1940) asserts that if
εz(h) 6= 0 for an element h in H and some z in the center of G, then h = z.
Coming to Marciniak’s question, suppose that G has no subgroups isomorphic to
Klein’s four group. For our purpose, we can assume that O2′(G) = 1 and that Sylow
2-subgroups of G are not cyclic. Thus Sylow 2-subgroups of G are (generalized)
quaternion, and by the Brauer–Suzuki theorem, from Brauer and Suzuki (1959), G
contains a unique involution z. For an involution u in V(ZG), the Cohn–Livingstone
result gives εz(u) 6= 0, and therefore u = z by the Berman–Higman result, answering
Marciniak’s question in the affirmative.
Theorem B (Kimmerle). Let G be a finite group. Suppose that V(ZG) has a
subgroup isomorphic to Klein’s four group. Then the same is true for G.
We do not know of a proof avoiding the use of the Brauer–Suzuki theorem.
Suppose that Sylow 2-subgroups of G are quaternion groups. Then the theorem
implies that finite 2-subgroups of V(ZG) are cyclic or quaternion groups. Taking
UNIT GROUPS WITH NO NONCYCLIC ABELIAN FINITE SUBGROUPS 3
into account the structure of the quaternion groups, and the Cohn–Livingstone
result, one obtains the following corollary.
Corollary 2. Let G be a finite group whose Sylow 2-subgroups are quaternion
groups (ordinary or generalized). Then any finite 2-subgroup of V(ZG) is isomor-
phic to a subgroup of G.
3. Proof of Theorem A
The partial augmentations of a torsion unit in V(ZG) encode its character values
in a way establishing a connection to group elements which respects a divisibility
relation between orders. We will make use of a lemma which is an easy consequence
of this fact.
Lemma 3. Let u be a torsion unit in V(ZG) of, say, order n. Let s be a natural
integer coprime to n, so that st ≡ 1 mod n for another natural integer t. Then for
all x in G whose order divide n, we have εx(u
s) = εxt(u).
Proof. Let ζ be a primitive n-th complex root of unity, and let σ be the Galois auto-
morphism ofQ(ζ) sending ζ to ζs. Let x1, . . . , xk be representatives of the conjugacy
classes of G whose elements have order dividing n. Note that then xt1, . . . , x
k is an-
other system of representatives. By Theorem 2.3 in Hertweck (2007a), εx(u) 6= 0 is
possible only for elements x whose order divide n. Thus for any ordinary irreducible
character χ of G, we have
εxi(u
s)χ(xi) = χ(u
s) = χ(u)σ =
εxi(u)χ(xi)
εxi(u)χ(x
i ) =
(u)χ(xi).
Since the character table of G, stripped off from any additional information, is an
invertible matrix, it follows that εxi(u
s) = εxt
(u) for all indices i, which proves the
lemma. �
Corollary 4. Let u be a torsion unit in V(ZG) of, say, order n. Then for any x
in G whose order divides n,
s∈(Z/nZ)×
s∈(Z/nZ)×
εxs(u).
Corollary 5. Suppose that for a prime divisor p of the order of G, all elements
of order p in G are conjugate to a power of some fixed element x. Let u be a
torsion unit in V(ZG) of order p. Then
i=1 u
i and
i=1 x
i have the same
partial augmentations.
Proof. Let k be the number of conjugacy classes of elements of order p in G. By
Corollary 4 and Theorem 2.3 in Hertweck (2007a),
εxi(u) =
yG : y∈〈x〉
εy(u) =
Applying again Theorem 2.3 from Hertweck (2007a), the corollary follows. �
4 MARTIN HERTWECK
We will apply this by means of the following formula relating ranks of an idem-
potent to arithmetical properties of the group.
Corollary 6. Suppose that for a prime divisor p of the order of G, all elements
of order p in G are conjugate to a power of some fixed element x. Suppose further
that V(ZG) contains an elementary abelian subgroup U of order p2. Then for any
ordinary character χ of G,
(1) χ
χ(1) + (p+ 1)
χ(xi)
We now turn to the proof of Theorem A. Suppose that G has a cyclic Sylow
p-subgroup P (p = 2 is allowed). Let x be an element of order p in P , and
set N = NG(〈x〉). Suppose further that V(ZG) contains an elementary abelian
subgroup U of order p2. Let χ be the character of G which is induced from the
principal irreducible character of P . Then the rank in (1) is
(|G : P |+ |N : P |(p2 − 1)).
If χ is a character of G which is induced from a faithful irreducible character of P ,
the rank in (1) is
(|G : P | − |N : P |(p+ 1)).
The difference of these ranks is |N : P |(p2 + p)/p2, which is not an integer. This
contradiction proves the theorem.
In view of Corollaries 1 and 2, one may be tempted to investigate the analogous
problem for groups with dihedral Sylow 2-subgroups. These groups were classified
by Gorenstein and Walter, and listed, for example, on p. 462 in Gorenstein (1968).
To indicate what can be done by now, we end with an example.
Note that the order of a finite subgroup of V(ZG) divides the order of G, see
Lemma 37.3 in Sehgal (1993); a fact which, surprisingly enough from today’s point
of view, is in this generality not recorded in Higman’s thesis.
Example 7. For the alternating group A7, any finite 2-subgroup of V(ZA7) is
isomorphic to a subgroup of A7.
Proof. Sylow 2-subgroups of A7 are dihedral of order 8. Let x be an element of
order 4 in A7. Then x
G and (x2)G are the only conjugacy classes of elements of
order 4 and 2, respectively. There is an (irreducible) character χ of A7 of degree 6
which is afforded by a deleted permutation representation. We have χ(x) = 0 and
χ(x2) = 2.
Let U be a finite 2-subgroup of V(ZA7). If U is of order 2, then U is rationally
conjugate to a subgroup of A7 by Corollary 3.5 in Hertweck (2006). If U is of
order 4, the Luthar–Passi method as described in Hertweck (2007a) is not sufficient
to guarantee rational conjugacy to a subgroup of A7: for a unit u of order 4 in
V(ZA7) one cannot exclude the possibility of having (εx2(u), εx(u)) = (2,−1) when
χ(u) = 4. In this case, also χ(u−1) = 4. Anyway, U is isomorphic to a subgroup of
A7, and the same is true if U is a Klein’s four group.
Suppose that U is abelian of order 8. By the Cohn–Livingstone result, U is not
cyclic. Set e = 1
u∈U u. Since e is an idempotent, χ(u) is a rational integer. If
U is elementary abelian, then χ(e) = 1
(χ(1) + 7χ(x2)) = 20
, which is impossible.
UNIT GROUPS WITH NO NONCYCLIC ABELIAN FINITE SUBGROUPS 5
Thus U contains 3 elements of order 2 and 4 elements of order 4. Trying out all
possibilities shows that again χ(e) is not a rational integer.
It remains to consider the case when U is the quaternion group. Let u be an
element of order 4 in U . Since χ(u2) = χ(x2), the restriction of the character χ
to U is the sum of four linear characters and the one of degree two. But this is
impossible since χ is afforded by a rational representation, while the character of
degree two of the quaternion group comes from the block of the rational quaternion
algebra (whence the name of the group). �
References
Berman, S. D. (1955). On the equation xm = 1 in an integral group ring. Ukrain. Mat.
Ž. 7:253–261.
Brauer, R., Suzuki, M. (1959). On finite groups of even order whose 2-Sylow group is a
quaternion group. Proc. Nat. Acad. Sci. U.S.A. 45:1757–1759.
Cliff, G. H., Sehgal, S. K., Weiss, A. R. (1981). Units of integral group rings of metabel-
ian groups. J. Algebra 73(1):167–185.
Cohn, J. A., Livingstone, D. (1965). On the structure of group algebras. I. Canad. J.
Math. 17:583–593.
Dokuchaev, M. A., Juriaans, S. O. (1996). Finite subgroups in integral group rings. Can-
ad. J. Math. 48(6):1170–1179.
Gorenstein, D. (1968). Finite groups. New York: Harper & Row.
Hertweck, M. (2006). On the torsion units of some integral group rings. Algebra Colloq.
13(2):329–348.
Hertweck, M. (2007a). Partial augmentations and Brauer character values of torsion units
in group rings. Comm. Algebra, to appear (e-print arXiv:math.RA/0612429v2).
Hertweck, M. (2007b). The orders of torsion units in integral group rings of finite solvable
groups. Comm. Algebra, to appear (e-print arXiv:math.RT/0703541).
Higman, G. (1940). Units in group rings. Ph.D. thesis. University of Oxford (Balliol
College).
Kimmerle, W. (2006). Arithmetical properties of finite groups. Talk delivered at the Math
Colloquium of the Vrije Universiteit Brussel.
Marciniak, Z., Sehgal, S. K. (2003). The unit group of 1 + ∆(G)∆(A) is torsion-free. J.
Group Theory 6(2):223–228.
Roggenkamp, K. W. (1981). Units in integral metabelian grouprings. I. Jackson’s unit
theorem revisited. Quart. J. Math. Oxford Ser. (2) 32:209–224.
Sandling, R. (1981). Graham Higman’s thesis “Units in group rings”. In: Integral rep-
resentations and applications (Oberwolfach, 1980). Lecture Notes in Math. Vol. 882.
Berlin: Springer, pp. 93–116.
Sandling, R. (1984). The isomorphism problem for group rings: a survey. In: Orders
and their applications (Oberwolfach, 1984). Lecture Notes in Math. Vol. 1142. Berlin:
Springer, pp. 256–288.
Sehgal, S. K. (1993). Units in integral group rings. Pitman Monographs and Surveys in
Pure and Applied Mathematics Vol. 69. Harlow: Longman Scientific & Technical.
Weiss, A. (1988). Rigidity of p-adic p-torsion. Ann. of Math. (2) 127(2):317–332.
Zassenhaus, H. (1974). On the torsion units of finite group rings. In: Studies in mathe-
matics (in honor of A. Almeida Costa). Lisbon: Instituto de Alta Cultura, pp. 119–126.
Universität Stuttgart, Fachbereich Mathematik, IGT, Pfaffenwaldring 57, 70550
Stuttgart, Germany
E-mail address: [email protected]
arXiv:math.RA/0612429v2
arXiv:math.RT/0703541
1. Introduction
2. Kimmerle's observation
3. Proof of Theorem A
References
|
0704.0413 | Exotic Hadron in Pole-dominated QCD Sum Rules | Exotic Hadron in Pole-dominated QCD Sum Rules
Toru Kojo 1,∗), Daisuke Jido, 2 and Arata Hayashigaki 3
1 Department of Physics, Kyoto University, Kyoto 606-8502, Japan
2 Yukawa Institute for Theoretical Physics, Kyoto University, Kyoto 606–8502,
Japan
3Institut für Theoretische Physik, J.W. Goethe Universität, D-60438 Frankfurt am
Main, Germany
We study pentaquark Θ+(I = 0, J = 1/2) in the QCD sum rules emphasizing that we
can not extract any properties of the pentaquark outside of the Borel window. To find the
appropriate Borel window, we prepare a favorable set up of the correlators and carry out
the operator product expansion up to dimension 15 within factorization hypothesis. Our
procedures reduce the unwanted high energy contaminations and enhance the low energy
correlation. In the Borel window, independent analyses for the chiral-even and odd sum
rules give the consistent values of the Θ+ mass, 1.68±0.22 GeV, and the residue. The parity
is found to be positive.
§1. Introduction
The experimental announcement for the discovery of the pentaquark Θ+(1540)1)
triggered tremendous amount of theoretical and experimental works on the exotic
states. Although the existence of such exotic states is still not so obvious, the exotics
provide a good opportunity to get the deeper insight of the hadron structures and
their connection to QCD. One of approaches from QCD to exotics is the QCD sum
rule (QSR),2) which relates informations of QCD to the hadronic properties through
the correlator analysis for the interpolating fields of hadrons. The Borel transformed
sum rules with the simplest pole + continuum parametrization are given as (i = 0, 1
correspond to the chiral even and odd part, respectively)
(ope)
i (−Q
2) = λ2i e
−m2/M2 +
ds e−s/M
(ope)
i (s), (1
where the relation ±mλ20 = λ21 holds due to the spinor structure and the relative
sign of the residues λ2i represents with the parity of the resonance state. Using these
sum rules, we can derive the approximated expressions of the mass and residue as a
function of M and sth. To extract properties of the low energy excitations with good
accuracy, we need to treat sum rules in the appropriateM2-region, i.e., Borel window,
where the low energy correlation is large enough compared to the contaminations
from high energy components which have no relations with properties of low-lying
resonances. The setting the Borel window is the most important step in QSR and,
only within this window, we can evaluate the physical quantities.
In the exotic cases, as reported in Ref. 3), it is extremely difficult to find
the appropriate Borel window in contrast to the usual meson and baryon cases.
∗) e-mail address: [email protected]
typeset using PTPTEX.cls 〈Ver.0.9〉
http://arxiv.org/abs/0704.0413v1
2 Toru Kojo, Daisuke Jido and Arata Hayashigaki
This is because the OPE convergence is very slow and the unwanted high energy
components dominate the spectral integral. In addition, we often encounter the
artificial stability of the physical quantities against M2-variation. This is the case
that physical quantities depend too strongly on the threshold parameter sth and not
on the low energy correlations which we want to extract. To attack these serious
problems common to the exotics, we propose a new approach and apply it to the Θ+,
assuming its quantum number as I = 0, J = 1/2, as one example of the exotics.4)
§2. OPE and favorable set up of the correlators
To find the Borel window, it is necessary to increase low energy informations in
the spectral function efficiently and, at the same time, reduce high energy contami-
nations. For these purposes, we take the following treatments.
Inclusion of the higher dimension terms of OPE is particularly important be-
cause they strongly reflect the low energy dynamics. For example, in the case of the
sum rules for ρ and A1 mesons, the dim.0 and 4 terms are the same due to the chiral
symmetry realized in the high energy, and the splitting of the masses is explained
only after the inclusion of dim.6 terms, 〈q̄q〉2, which appear due to the chiral sym-
metry breaking. From these observations, we perform the OPE calculation up to
dim.15 within factorization hypothesis both for taking into account the low energy
correlations and for the confirmation of good OPE convergence.
As later shown, simple inclusion of the low energy correlations through the higher
dimension terms is found to be not sufficient to find the Borel window because
high energy contaminations are too large in the QSR for the exotics. To reduce
the high energy contaminations, we take the difference between correlators for two
interpolating fields with similar structure but different chirality each other, i.e.,
d4x eiq·x
∣T [P (x)P̄ (0) − t S(x)S̄(0)]
ds e−s/M
Im[ΠP0 (s)− tΠS0 (s)] q̂ + Im[ΠP1 (s)− tΠS1 (s)]
, (2.1)
where Π0, Π1 correspond to the chiral even and odd part respectively, and
P = ǫabcǫdef ǫcfg{uTaCdb}{uTd Cγµγ5de}γµCsTg , (2.2)
S = ǫabcǫdef ǫcfg{uTaCγ5db}{uTd Cγµγ5de}γµCsTg . (2.3)
Here the only difference in these interpolating fields is that the first diquark structures
have the opposite chirality.
Let us first explain in the case of the chiral even part. Since they show the same
behavior in high energy due to the chiral symmetry, after the subtraction of two
correlators (t = 1 case), the irrelevant high energy contributions are expected to be
canceled out in the similar way as the Weinberg sum rules.5) In terms of OPE, this
cancellation corresponds to the cancellation of the lower dimension terms. It is not
a priori evident whether the low energy correlations remain enough even after the
subtraction because the low energy contribution could also cancel out. Our Borel
Exotic Hadron in Pole-dominated QCD Sum Rules 3
analysis, however, reveals that, in the case of t = 1, the large low energy correlation
remains enough even after the subtraction. As a result, we can find the Borel window
in the relatively large M2-region.
On the other hand, for the chiral odd part, the subtraction procedure corre-
sponding to t = 1 case does not lead the cancellation of the high energy components
because chiral odd part is constructed of the chiral symmetry breaking terms. How-
ever, in the case of t = 1, the OPE convergence is found to be very good and then we
can find the Borel window in the small M2-region where high energy contaminations
are suppressed due to the Borel factor e−s/M
in the integral of the spectral function.
§3. Borel analysis for mass and residue
Here we explain our criterion for the Borel window. The lower bound of the
Borel window is given so that the highest-dimensional terms in the truncated OPE
are less than 10% of its whole OPE, while the upper bound is determined by the
region where the absolute value of the pole contribution is larger than the absolute
value of the integrated spectral function in the region s ≥ sth. Note that the 50%
pole contribution in our criterion is extremely large in comparison with any prior
pentaquark sum rules, where the pole contributions are not more than 20%.3)
To recognize the problems in the case of QSR for exotics, let us see Fig.1 for
M2-dependence of the mass in the cases of t = −1, 0, 10 corresponding to PP̄ +SS̄,
PP̄ , SS̄ cases respectively. The threshold value is fixed to typical value
sth = 2.2
GeV. In these cases, we fail to find stabilities of the mass in the M2-region lower
than the upper bound of the Borel window. The stabilities above the upper bound
are simply artifacts which often appear in QSR. Fig.1 shows that typical mass of PP̄
case is much smaller than that of SS̄, and then we can expect that the low energy
correlation of PP̄ is much larger than that of SS̄. This observation leads that even
after the subtraction PP̄ − SS̄ (t = 1 case), the low energy correlation can remain
enough.
Now we see the case of around t = 1. We tune the value of t around t = 1 to
get the widest Borel window. As expected, for even part (t = 0.9), the high energy
SS (t = 10)
PP (t = 0)
PP (t = 0)
SS (t = ‐1)
SS (t = ‐1)
SS (t = 10)
Fig. 1. The behavior of the mass as a function of M2 for t = −1, 0, 10. The left arrows represent
the upper bound of the Borel window. In the smaller M2-region than the upper bound, we can
not find stable region of the mass. The stabilities above the upper bound are simply artifacts
which often appear in QSR.
4 Toru Kojo, Daisuke Jido and Arata Hayashigaki
even odd
(t=1.1)(t=0.9)
Fig. 2. The behavior of the mass as a function of M2. The left (right) arrows represent the upper
(lower) bound of the Borel window. We succeed to find the Borel window and stabilities of
mass.
contaminations are canceled out due to chiral symmetry and we find the wide Borel
window in the relatively large M2-region. On the other hand, for odd part (t = 1.1),
thanks to the good OPE convergence, we also find the wide Borel window in the
small M2-region. The threshold values are taken to make the physical quantities
most stable in the Borel window.
The best stability is achieved with
sth = 2.2 GeV (even) and 2.1 GeV (odd),
giving mΘ+ = 1.64 GeV (even) and 1.72 GeV (odd) respectively. The values of the
residue are also obtained from the chiral even and odd sum rules as λ20 = (3.0±0.1)×
10−9 GeV12 and λ21/mΘ+ = (3.4 ± 0.2) × 10−9 GeV12. It is remarkable that these
numbers are quite similar with the close t. This implies that our analysis investigates
consistently the same state in the two independent sum rules. Note that from the
relative sign of the residues, we assign positive parity to the observed Θ+ state.
In conclusion, we perform the Borel analysis for Θ+ with special setup of the
correlators in order to find the Borel window. Within uncertainties of the condensate
value, independent analyses for the chiral-even and odd sum rules give the consistent
values of the Θ+ mass, 1.68± 0.22 GeV, and the residue. The parity is found to be
positive.
Acknowledgements
We thank Profs. M. Oka, A. Hosaka and S.H. Lee for useful discussions about
QSR for the exotics during the YKIS2006 on ”New Frontiers on QCD” held at
the Yukawa Institute for Theoretical Physics. This work is supported in part by
the Grant for Scientific Research (No.18042001) and by Grant-in-Aid for the 21st
Century COE ”Center for Diversity and Universality in Physics” from the Ministry
of Education, Culture, Sports, Science and Technology (MEXT) of Japan.
References
1) T. Nakano et al., Phys. Rev. Lett. 91 (2003), 012002.
2) M.A. Shifman, A.I. Vainshtein, and V.I. Zakharow, Nucl. Phys. B 147 (1979), 385.
3) R.D. Matheus and S. Narison, Nucl. Phys. Proc. Suppl. 152 (2006), 236.
4) T. Kojo, A. Hayashigaki, and D. Jido, Phys. Rev. C 74 (2006), 045206
5) S. Weinberg, Phys. Rev. Lett. 18 (1967), 507.
Introduction
OPE and favorable set up of the correlators
Borel analysis for mass and residue
|
0704.0414 | Leaky modes of a left-handed slab | Leaky modes of a left-handed slab
A. Moreau
LASMEA, UMR CNRS 6602, Université Blaise Pascal, 24 avenue des Landais, 63177
Aubière, France.
D. Felbacq
GES UMR CNRS 5650, Université de Montpellier II, Bat. 21, CC074, Place E.
Bataillon, 34095 Montpellier Cedex 05, France.
Using complex plane analysis we show that left-handed slab may support either leaky
slab waves, which are backward because of negative refraction, or leaky surface waves,
which are backward or forward depending on the propagation direction of the surface
wave itself. Moreover, there is a general connection between the reflection coefficient of
the left-handed slab and the one of the corresponding right-handed slab (with opposite
permittivity and permeability) so that leaky slab modes are excited for the same angle of
incidence of the impinging beam for both structures. Many negative giant lateral shifts
can be explained by the excitation of these leaky modes.
Keywords: Left-handed materials, leaky modes, complex plane analysis
1 Introduction
Left-handed materials [1] have long been considered a theoretical oddity. Since it has been
demonstrated that they could be produced using metamaterials [2], they have attracted
much attention. The basic physics of left-handed materials (LHM) is truly exotic and has
been completely ignored until recently, it renews the physics of lamellar structures to the
extend that a bare slab of LHM exhibits many surprising properties : it can for instance
support unusual guided modes [3,4] or behave as a perfect lens [5]. In this paper, we study
the exotic properties of the different types of leaky waves supported by a left-handed slab.
Given the importance of the left-handed slab for both fundamental and applied works,
there is obviously a need for a clear understanding of these properties.
We particularly show that two types of leaky waves are supported by such a structure
(i) leaky slab waves which are always backward due to negative refraction and (ii) leaky
surface waves which do not exist for a right-handed slab and which can be backward or
forward. The excitation of these modes leads to positive or negative giant lateral shifts,
the latter being rather exotic [6].
http://arxiv.org/abs/0704.0414v3
2 Leaky modes and giant lateral shifts
A leaky mode [6] is a solution of the wave equation which verifies the relation dispersion of
a structure but with a propagative solution above and (or) under the structure. Whereas
a guided mode has a real propagation constant, the propagation constant of a leaky
mode is complex because the energy of the waves leaks out of the structure and the
waves is attenuated. A leaky wave is thus a complex solution of the dispersion relation
and a complex plane analysis is thus particularly relevant for a thorough analysis of its
properties. Let us underline that a leaky mode may be either forward, which is common, or
backward, leading to a propagation constant which has a positive (respectively negative)
imaginary part.
Let us consider a slab characterized by ε2 and µ2 surrounded by right-handed media with
ε1 and µ1 (resp. ε3 and µ3) above (resp. under) the slab as shown figure 1. The values
we have chosen for ε2 and µ2 are arbitrary but realistic [7] so that this structure could be
realized using split-ring resonators and wires.
ε µ3 3
ε µ1 1
ε µ2 2
Figure 1: The LHM slab of thickness h surrounded by right-handed media.
We may assume that ε1 µ1 ≥ ε3 µ3 with no loss of generality.
The relation dispersion of such a structure can be written
r21 r23 = exp(−2iγ2 h) (1)
where γi =
εi µi k
0 − α2, k0 = ωc =
and rij =
κi−κj
κi+κj
with κi =
in TE polarization
(or κi =
in TM polarization). Since ε1 µ1 ≥ ε3 µ3 and we are concerned with leaky
waves, we will only consider values of α such that α <
ε1 µ1 k0, which means that the
solution will always be propagative at least in medium 1.
Let us now consider the reflection coefficient of a plane wave exp(i(αx+γ z−ω t)) coming
from upwards with an angle of incidence θ so that α = n k0 sin θ. Its reflection coefficient
can be written
r23 exp(2iγ2 h)− r21
1− r21 r23 exp(2iγ2 h)
using the above definitions.
It is obvious that when the relation dispersion is verified, then the reflection coefficient
presents a pole. A leaky mode thus corresponds to a pole of the reflection coefficient.
A zero, located on the other side of the real axis, corresponds to each pole. As we will
see in the following, a zone where the phase of r quickly varies lies between a pole and
its corresponding zero. This zone crosses the real axis, so that the presence of a pole is
responsible for a swift variation of the phase on the real axis.
When considering the reflection of a gaussian beam on a structure whose reflection co-
efficient has a modulus equal to one (so that it can be written r = eiφ), the lateral
displacement of the reflected beam’s barycenter along the interface is given by the well
known formula
δ = −dφ
. (3)
This lateral displacement is the sign that a leaky wave has been excited by the incident
beam. The reflected beam then has two components : the part which is reflected by
the first interface of the structure (whose barycenter is not particularly displaced) and
the leaky wave itself [6]. The reflected beam is heavily distorted by the leaky wave and
presents an exponentially decreasing tail so that its barycenter is largely displaced : this
is the so-called giant lateral shift.
This effect is sometimes called a giant Goos-Hänchen effect, but in this case the shift is
due to the excitation of a leaky mode [6] and not, as in the real Goos-Hänchen effect [8,9],
to the total reflection.
3 The left-handed slab
With left-handed materials, though, negative lateral shifts seem to be much more common
[10–14] than once expected [6]. Here we will consider the case of a left-handed slab (i.e.
if ε2 < 0 and µ2 < 0) and explain why the leaky modes supported by such a structure are
usually backward. Our explanations will be supported by a complex plane analysis of the
leaky modes.
Here the expression (2) of the reflection coefficient remains perfectly valid. We will now
distinguish two cases : the case when the solution is propagative in the left-handed medium
and the case when the solution is evanescent in region 2.
3.1 Leaky slab modes
When the field is propagative in the left-handed slab, large negative lateral shifts have
been reported but not interpreted [13]. These shifts are due to the excitation of leaky slab
modes or Perot-Fabry resonances of the slab at non normal incidence. Such leaky modes
have already been studied for a right-handed slab [15] and they can be considered as
constructive interferences of the multiple beams which are produced by reflections on the
interfaces of the slab. In the case of a left-handed slab, since the first beam undergoes a
negative refraction as shown figure 2 these constructive interferences will logically generate
a backward leaky mode. We may thus conclude that the existence of such a backward
leaky mode is linked to negative refraction.
Figure 2: Modulus of the field for a thick left-handed slab with ǫ1 = ǫ3 = µ1 = µ3 = 1,
ǫ2 = −3,µ2 = −1 and h = 60 λ using a gaussian incident beam with a waist of 20 λ and
an incidence angle of θ = 45.
This argument is not a proof, though : unexpected lateral shifts have been reported when
the beams interfere destructively [16]. But if the leaky modes are backward, then the
corresponding solutions of the dispersion relation and the poles of the reflection coefficient
should have a negative imaginary part. This is what is shown figure 3.
Figure 3: The phase of the reflection coefficient in a part of the complex plan [0, n1 k0] +
i[−k0
]. Each black point represents a pole and each white point a zero. The cut line is
clearly visible here. The rapid variation of the phase which is due to each pole is obvious.
Two types of leaky slab waves should be distinguished (i) L2 waves which are leaky in
both the upper and the lower media and (ii) L1 waves which are leaky only in the upper
medium and evanescent in the lower one. The latter correspond to the poles located under
the cut line.
Using complex plane analysis we will now try to show that all the solutions of the disper-
sion relation 1 are located in the lower part of the complex plane, meaning that all the
leaky modes are backward.
When the relation dispersion is satisfied, then the following condition holds :
|r23 r21| = e2 γ
h. (4)
As demonstrated in the annex, |rij > 1 whenever one of the media is left-handed. Since
medium 2 is left-handed then the condition
h > 1 (5)
should be satisfied, which is possible for γ′′2 > 0 and therefore for α
′′ < 0 (see the annex
for details). The fact that rij > 1 is thus the main reason why the poles of r are under
the axis and why the leaky slab modes are backward.
We must underline the fact that our demonstration is valid only for the first Riemann
sheet : our proof cannot exclude that there may be some poles on the other Riemann sheet
above the real axis, corresponding to forward L1 leaky slab waves when ε1 µ1 > ε3 µ3. But
we could not find any.
3.2 Leaky surface modes
Let us now consider the situation in which the field is evanescent in the left- handed
medium. Then γ2 is purely imaginary on the real axis. Since e
2 γ′′
h tends towards infinity
when h → +∞ then relation (4) can be verified only if r23 has a pole (r21 cannot have
one since the field is always propagative in the upper medium). This means that the
structure may support a leaky mode only if the interface between medium 2 and 3 can
support a guided mode. It is now well-known that such an interface actually supports a
surface mode [17,18] which can, depending on media 2 and 3, be backward (resp. forward)
corresponding to a pole under the real axis (resp. above the real axis but on the other
Riemann sheet). The leaky wave always has the same propagation direction as the surface
mode, whatever the thickness of the slab, as shown figure 4. In the case of a forward leaky
wave, only the zero belongs to the first Riemann sheet, just under the real axis. The pole
shown figure 4 belongs to the other Riemann sheet.
0.05
0.15
0.25
0.35
1.05 1.1 1.15 1.2 1.25 1.3 1.35 1.4 1.45 1.5
−0.45
−0.35
−0.25
−0.15
−0.05
1.7 1.75 1.8 1.85 1.9 1.95 2 2.05 2.1 2.15 2.2
Figure 4: Location of the poles in the α
complex plane for different values of h with
ε1 = 9, µ1 = µ3 = ε3 = 1 and (a) ε2 = −0.5 and µ2 = −1.5, showing a forward surface
mode and (b) ε2 = −5 and µ2 = −0.5, showing a backward surface mode.
Figure 5 finally shows the excitation of a backward leaky surface wave by a gaussian
beam. The chosen values of µ2 may be obtained with simple split ring resonators [19] for
instance.
Figure 5: Modulus of the field for a left-handed slab with ǫ1 = 9,ǫ3 = µ1 = µ3 = 1,
ǫ2 = −0.5,µ2 = −1.5 and h = 0.6 λ using a gaussian incident beam with a waist of 20 λ
and an incidence angle of θ = 21.496. The pole corresponding to the leaky mode is located
at αp = (1.0993 + 0.001267i) k0.
4 Fundamental property
Let us a consider a structure with left-handed materials. We will call corresponding
right-handed structure the structure obtained by replacing any left-handed medium by a
medium with opposite permittivity and permeability, without changing the geometrical
parameters.
In this section, we will concentrate on the link between the reflection coefficient of a
left-handed slab and the one of its corresponding right-handed structure.
Let us consider the interface between a right-handed medium labelled i and a left-handed
medium j. The reflection coefficient of such an interface is rij . We will now define
r+ij the reflection coefficient of an interface between medium i and right-handed medium
characterized by |εj| and |muj|. It is not difficult to see, from the expression of rij that
. (6)
This allows to understand why the Goos-Hänchen shift of an interface between a right-
and a left-handed medium is the opposite of the corresponding right-handed structure [11]
since the phases of both structures are opposite on the real axis.
The reflection coefficient r can now be written
e2iγ2 h
1− e2iγ2 h
r+23 e
−2iγ2 h − r+21
1− r+21 r+23 e−2iγ2 h
Since
except when z is on the cut line, then γ(z∗) = γ(z)∗ and hence r+ij(z)
∗) so that
r(z)∗ =
∗) e2iγ2(z
∗)h − r+21(z∗)
1− r+21(z∗) r+23(z∗) e2iγ2(z
, (10)
which can simply be written
r(z)∗ = r+(z∗), (11)
where r+ is the coefficient reflection of the corresponding right-handed slab. Note that
this relation does not hold on the cut line, but that it holds for the two Riemann sheets.
This means that the poles of the left-handed slab and the poles of the corresponding
right-handed slab are symmetrical with respect to the real axis. This means that L2
waves can be excited for the same incidence angle for both structures. This is not the
case for L1 modes : the function r on the real axis is continuous with the lower part of
the first Riemann sheet whatever the situation and the poles which are above the cut line
thus have no effect on the real axis.
As an example, we have computed the field in TE polarization inside and around the slab
when it is illuminated with a gaussian beam for the left-handed slab and its corresponding
right-handed structure. The results are shown figures 6 and 7.
Figure 6: Modulus of the field for a symmetrical slab with ǫ1 = ǫ3 = 9, µ1 = µ3 = 1,
ǫ2 = 1.5, µ2 = 1 and h = 1.3 λ using a gaussian incident beam with a waist of 20 λ and
an incidence angle of θ = 22.78.
Figure 7: Modulus of the field for a symmetrical slab with ǫ1 = ǫ3 = 9, µ1 = µ3 = 1,
ǫ2 = −1.5, µ2 = −1. and h = 1.3 λ using a gaussian incident beam with a waist of 20 λ
and an incidence angle of θ = 22.78. The pole corresponding to the leaky mode is located
at αp = (1.16823− 0.01125i) k0
5 The grounded left-handed slab
The grounded left-handed slab is a much more simple structure for (i) there is no need to
distinguish two different types of leaky slab modes and (ii) the structure can not support
any leaky surface mode. All the leaky modes are then slab modes and are found for
α < n2 k0. The reflection coefficient of the grounded slab is given by (2) with r23 = −1 for
the TE polarization and r23 = 1 for the TM polarization so that the relation dispersion
gives
|r12| = e2 γ
h. (12)
Since |r12| > 1 then all the solutions of the dispersion relation are located in the lower
part of the complex plane so that they are all backward.
It is then easy to show that the relation r+(z)∗ = r(z∗) still holds. As a consequence,
the leaky modes of a grounded left-handed slab and of its corresponding right-handed
structure can be excited for the same angle of incidence of the impinging beam.
6 Conclusion
In this paper, we have thoroughly studied the leaky modes of a left-handed slab for realistic
values of the permittivity and permeability of the left-handed medium [7,19,20] which can
be obtained using structures like split-ring resonators. Our results can be summarized as
follows. Left-handed slab may support two types of leaky modes :
• Leaky slab modes, which are always backward because of the negative refraction
phenomenon. When the transmission is not null, leaky modes of the left-handed
slab and of its corresponding right-handed structure are excited for the same angle
of incidence.
• Leaky surface modes, which may be backward or forward depending on the propa-
gation direction of the surface wave itself.
This work could help to interpret many giant lateral shifts as excitations of exotic leaky
waves [12, 13, 16]. Since the existence of backward slab waves is linked to the property
of negative refraction, and since these leaky waves constitute a signature of a left-handed
slab behavior we think that they could be used to characterize the left-handedness of
metamaterial or photonic crystal structures far better than other methods [21].
Acknowledgments
This work has been supported by the French National Agency for Research (ANR), project
030/POEM. The authors would like to thank Alexandru Cabuz and Kevin Vynck for their
help.
Annex
In this annex, we will clearly define the choice we have made for the definition of the
complex square root and prove that for z on the first Riemann sheet (but not on the cut
line) we have |rij(z)| > 1 when media i and j are not both right-handed.
Since the square root can be continued on the complex plane, r and rij can be continued
as well. We have chosen to take
2 with z = r eiθ and θ ∈]−π, π], as a definition
of the square root. This means that we have placed the cut line on the negative part of
the real axis and if x is a positive real,
−x = i
x. This defines the square root on the
entire complex plane, to which we refer as the first Riemann sheet. When we write that
z is on the second Riemann sheet, it will mean that we have taken the opposite of
defined above.
With this choice, we have (i) ℜ(
z) ≥ 0 (ii)
for z on both sheets but not on
the cut line (iii) if ℑ(z) < 0, ℑ(
z) < 0 and if ℑ(z) > 0, ℑ(
z) > 0 (iv) the function
γ(z) =
ǫ µ k20 − z2 has a cut line on the real axis (on ] − ∞,−n k0] ∪ [n k0,+∞] more
precisely) and the function γ on the real axis is continuous with the part of the complex
plane which is under the cut line : when z passes through the cut line from the first
Riemann sheet (coming from the lower part of the plane) to the second Riemann sheet,
γ(z) is continuous. When a function which can be written using γ(z) presents a pole,
it must be found either (i) for z on the first Riemann sheet and under the real axis (we
will say that the pole itself is on the first Riemann sheet in this case) or (ii) for z on the
second Riemann sheet but above the real axis.
We have
rij =
κi − κj
κi + κj
. (13)
The modulus of rij reads as
|rij|2 =
(κi − κj) (κ∗i − κ∗j )
(κi + κj) (κ
i + κ
|κi|2 + |κj|2 − 2 (κ′i κ′j + κ′′i κ′′j )
|κi|2 + |κj|2 + 2 (κ′i κ′j + κ′′i κ′′j )
, (15)
where κ = κ′ + i κ′′.
Let us define x and y the real and imaginary part of z = x + i y on the first Riemann
sheet. Let us assume that x > 0. We have
n2 k20 − z2 =
n2 k20 − x2 + y2 − 2 i x y. (17)
If y > 0, then x y > 0 and thus ℑ(n2 k20 − z2) < 0 so that finally ℑ(γ) < 0. If y < 0,
then x y < 0 so that ℑ(γ) > 0. Since γ(−z) = γ(z) the result will hold for x < 0 too and
for x = 0, γ(z) is real and positive so that the result obviously holds. So the imaginary
part of γ(z) is positive (resp. negative) when the imaginary part of z is negative (resp.
positive).
For any right-handed medium, κ has the same property than γ. For a left-handed medium,
since κ = γ
or κ = γ
depending on the polarization, the imaginary part of κ has the sign
of ℑ(z). Since i and j are not both right-handed, then κ′′i and κ′′j have not the same sign
and the product κ′′i κ
j is always negative. Since ℜ(
z) > 0 for all z on the first Riemann
sheet then κ′i κ
j is always negative too.
Finally, since κ′i κ
j + κ
j < 0, we have |rij| > 1 for all z except on the real axis. Please
note that rij is not, in the particular case of a left-handed medium, the reflection coefficient
on the interface [22].
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Introduction
Leaky modes and giant lateral shifts
The left-handed slab
Leaky slab modes
Leaky surface modes
Fundamental property
The grounded left-handed slab
Conclusion
|
0704.0415 | Coulomb blockade of field emission from nanoscale conductors | untitled
Coulomb blockade of field emission from nanoscale conductors
O. E. Raichev*
Institute of Semiconductor Physics, National Academy of Sciences of Ukraine, Prospekt Nauki 45, 03028, Kiev, Ukraine
�Received 9 February 2006�
Theoretical description of the field emission of electrons from nanoscale objects weakly coupled to the
cathode is presented. It is shown that the field-emission current increases in a steplike fashion due to single-
electron charging which leads to abrupt changes of the effective electric field responsible for the field emission.
A detailed consideration of the current-voltage characteristics is carried out for a nanocluster modeled by a
metallic spherical particle in the close vicinity of the cathode and for a cylindrical silicon nanowire grown on
the cathode surface.
PACS number�s�: 79.70.�q, 73.23.Hk, 73.40.Gk
I. INTRODUCTION
The discrete nature of electric charge reveals itself in the
transport of electrons through small conductors �nanopar-
ticles or other nanoscale objects� weakly coupled to the
source and drain electrodes �current-carrying leads� owing to
the Coulomb blockade effect. Numerous manifestations of
the charge quantization in transport properties, the most fa-
miliar of them are the Coulomb blockade oscillations of the
electric current as a function of the gate voltage and the
Coulomb staircase in the current-voltage characteristics,
have attracted considerable attention in the past years.1 Since
the fundamentals of the transport theory in the Coulomb
blockade regime have been established,2–4 the Coulomb
blockade-based physics has been applied to various issues of
electron transport in mesoscopic systems, and the field of its
applications expands in line with the advances in nanotech-
nology.
Usually, the influence of the Coulomb blockade on the
current in two-terminal devices is considered under assump-
tion that the coupling between the nanoscale object and the
leads is not sensitive to the number of electrons N determin-
ing the object charge eN. This corresponds to the introduc-
tion of ohmic �or nearly ohmic� effective resistances describ-
ing this coupling. Though this assumption often works well,
it can be violated, for example, in nanomechanical
systems,5–7 where charging of the object gives rise to its
displacement towards one of the leads thereby changing its
tunnel coupling to both leads. In this paper we study a situ-
ation when the sensitivity of the tunnel coupling to the num-
ber of electrons does not require a mechanical displacement
and is determined by the nature of tunneling. This implies a
device layout and conditions similar to those used in the
recent experiments on field emission of electrons from me-
tallic nanoclusters,8–10 silicon nanowires11–15 and
nanocones,16,17 and carbon nanotubes �see, for example,
Refs. 11 and 18–26�, when small �nanoscale� objects are
formed on the source electrode �cathode�, the latter is then
negatively biased with respect to the drain electrode �anode�
in vacuum. The current between the electrodes flows owing
to the field emission of electrons from nanoscale objects,
because the electric field F at the tips of the objects is higher
than in the other places of the device. The field-emission
current is described by the Fowler-Nordheim formula27
I = ASF2 exp�− F
�, F = 4�2m
3�e��
W3/2, �1�
where m is the free electron mass, W is the work function of
the emitting material, S is the effective emitting area, and A
is a constant expressed through the work function and Fermi
energy �F of the emitting material
�e�3��F/W
4�2 � ��F + W�
. �2�
The effective field F, which describes the tunnel coupling
between the nanoscale object and the anode, depends on the
object charge, which is induced by the applied voltage V
=V1−V2, where V1 and V2 are the cathode and anode poten-
tials, respectively. Under conditions of Coulomb blockade,
i.e., when the electric connection between the cathode and
the object is weak and the charging energy of the object
considerably exceeds the temperature T, the continuous
variation of the voltage V leads to discrete changes of the
object charge in units of e, and, consequently, to correspond-
ing discrete changes of the field F. Therefore, one may in-
troduce the field FN, which is a function of the discrete num-
ber N and continuous variable V. Next, if the current in the
device is limited by the field emission, the single-electron
tunneling processes become important. This means that, at a
fixed voltage V, the object stays mostly in the states with N
and N−1 electrons, the number N is determined by the volt-
age. In the N-electron state, no electrons can come to the
object from the cathode until an electron leaves the object by
tunneling through the barrier, see Fig. 1�a�. Then the object
appears in the N−1-electron state and returns to the
N-electron state before the next Fowler-Nordheim tunneling
event takes place. The field-emission current in these condi-
tions is given by Eq. �1� with F=FN and can be denoted as
IN. If the bias eV increases, the state with N+1 electrons
becomes more favorable, and the current changes in a step-
like fashion from IN to IN+1. This leads to staircaselike
current-voltage characteristics, which may look similar to the
usual Coulomb staircases.28–30 However, since the sensitivity
of the tunneling to the number of electrons is involved, the
staircaselike current-voltage characteristics can exist under
rather peculiar conditions, when the source-drain bias is or-
ders of magnitude larger than the charging energy.
The rest of the paper is devoted to quantitative studies
based on the physical idea outlined above. In Sec. II we give
the basic equations and calculate the current in the simplest
case of an idealized emitter shown in Fig. 1�b�. In Sec. III we
calculate the current from a nanocluster modeled by a spheri-
cal particle on the metallic cathode surface and from a semi-
conductor wire �nanowhisker� grown perpendicular to the
cathode surface. The discussion and concluding remarks are
given in the last section.
II. GENERAL CONSIDERATION
We consider the case of classical �or metallic� Coulomb
blockade, when the electron energy level separation in the
nanoscale object can be neglected in comparison to both
temperature and charging energy. Since the object is assumed
to be weakly coupled to the cathode, we study the sequential
tunneling process and not the coherent one. It is convenient
to investigate the electron transport by applying the kinetic
equation2 �Master equation� for the distribution function PN
describing the probability for the object to be in the state
with N electrons. Assuming that the electric connection be-
tween the cathode and the object is characterized by the con-
ductance G, this equation is written as
= QN+1 − QN, �3�
where
1 − exp�− �EN/T�
�PN − PN−1 exp�− �EN/T�
+ PNIN/�e� . �4�
Here �EN= �e
2 /C��N−1/2−C2V /e is the difference in Cou-
lomb energies for the objects with N and N−1 electrons, C is
the total capacitance, and C2 is the capacitance of the object
with respect to the anode �the capacitance with respect to the
cathode is given by C1=C−C2�. The first term in expression
�4� has the usual form2 and corresponds to the current be-
tween the object and the cathode. It is written as a difference
of the contributions describing the departure of an electron
from the object in the N-electron state and arrival of an elec-
tron at the object in the N−1-electron state. The second term
corresponds to the field-emission current from the object in
the N-electron state. Since no electrons come to the object
from the anode, this term does not contain a contribution
describing arrival of electrons. In the stationary case, Eq. �3�
is reduced to the form QN=const, where the constant can be
chosen equal to zero. After determining PN from the equation
QN=0 with the use of the normalization condition
NPN=1,
the total current is given by
PNIN. �5�
Under the condition GT� �e � IN, which means that the object
is in thermal equilibrium with the cathode, the stationary
solution of Eq. �3� is written as PN=Z
−1 exp�−EN /T�, where
EN= �e
2 /2C��N−C2V /e
2 is the Coulomb energy, and Z
N exp�−EN /T� is the partition function. The current in this
case is determined by the expression
J = Z−1
IN exp�− EN/T� . �6�
Let us apply the solution �6� to the idealized model of
emitter, Fig. 1�b�, when the emission takes place from a
spherical nanoparticle of radius R, placed at a distance d
from the cathode. The distance between the cathode and an-
ode is L. The connection c-p denotes a low-transparent con-
tact �for example, tunnel barrier� between the particle and the
cathode, which does not contribute to the field-emission
properties and electrostatics of the device. Assuming d�R,
we have C=R, C2=Rd /L, and neglect the charge polariza-
tion of the particle because this polarization is small in com-
parison to the total charge eN induced by the applied voltage.
The number of electrons is estimated as N�C2V /e
=RdF0 / �e�, where F0=−V /L is the applied electric field. The
effective field for the nanoparticle with N electrons is FN
= �e �N /R2, and the partial currents IN in these conditions are
given by
IN = AS�eN/R
2�2exp�− FR2/�e�N� , �7�
where the emitting area S, in the idealized model considered
here, can be approximated by the total surface area of the
nanoparticle, S=4�R2. In Fig. 2 we plot the current-voltage
characteristics of the idealized emitter, calculated according
to Eqs. �6� and �7�, where A is given by Eq. �2� with W
=5.1 eV and �F=5.5 eV �taken for Au�, and the geometrical
parameters are chosen as R=5 nm and d=0.5 �m. The char-
acteristics look like staircases with flat regions �plateaus� be-
tween the steps, which are visible even at room temperature.
It is possible to estimate the relative heights of the steps by
calculating the ratio of the currents IN and IN−1 emitted from
the nanoparticle with N and N−1 electrons
� exp� FR2
�e�N�N − 1�
� . �8�
In spite of the fact that the charged nanoparticle typically
contains a large number of electrons, N
100, one can al-
FIG. 1. �a� The mechanism of single-electron tunneling in the
Fowler-Nordheim regime. �b� Schematic representation of the ide-
alized emitter.
ways find a regime when the ratio IN / IN−1 is not small in
comparison to unity. This necessarily implies a weak Fowler-
Nordheim tunneling, when F /F=FR2 / �e �N�1.
In the calculations described above, the applicability of
the Fowler-Nordheim formula requires R�W / �e �F, which is
rewritten as R�e2N /W, or, according to N�RdF0 / �e�, as
�e �F0�W /d, independent of the nanoparticle radius. This
condition is satisfied at high enough applied voltages. If
�e �F0=eV /L�W /d, the approximation of a triangular poten-
tial barrier is not quite good, and one should consider the
tunneling through the barrier described by the potential en-
ergy W−e2N�1/R−1/r� at r�R, where r is the distance
from the center of the spherical nanoparticle; the tunneling
through the potential barrier of this form is described in Ref.
31. Even under the condition �e �F0�W /d, which is satisfied
in the calculations shown in Fig. 2, the relative change of the
current per one step, IN / IN−1−1, appears to be significant,
because the exponent FR2 / �e �N�N−1� in Eq. �8� is estimated
as c�W / �e �F0d�
2, where the dimensionless constant c
=4/3�2me4 /�2W is noticeably larger than unity.
If the current is high enough, the field emission cannot
remain the bottleneck for the electron transfer from the cath-
ode to the anode, and a finite resistance G−1 becomes essen-
tial. The nanoparticle in these conditions is no longer in equi-
librium with the cathode. This means that the distribution PN
is established kinetically, and several states with different
charges coexist at a fixed voltage �see the inset in Fig. 2�. As
a consequence, the Coulomb blockade features are washed
out. This case requires a numerical solution of the equation
QN=0. The corresponding current-voltage characteristics of
the idealized emitter calculated by using the RC time C /G
=100 ps are also shown in Fig. 2. The degradation of the
current steps appears to be stronger with increasing voltage,
because the current increases and the nanoparticle-cathode
link becomes more important. The shape of the steps in this
case resembles the usual Coulomb staircase.
III. MORE COMPLEX EXAMPLES
After demonstrating the possibility of the Coulomb-
blockade staircase of the field emission on a model example,
it is worth to consider more complex cases. Indeed, the
model example discussed above has certain disadvantages.
First of all, it is hardly possible to connect a particle placed
far from the cathode surface by a link �c-p in Fig. 1�b�
which does not contribute to the electrostatic properties of
the device. Second, the model of uniform charging is insuf-
ficient: the charge polarization of the nanoscale object ap-
pears to be important and should be always taken into ac-
count, see below in this section. Therefore, the model shown
in Fig. 1�b� is suitable only for the purposes of illustration of
the basic physics described by Eqs. �3�–�6�. To have a closer
approach to reality, we point out that the nanoscale objects
investigated in the above-cited experiments on field emission
can be roughly divided into two classes: the objects whose
dimensions in all directions are comparable �nanoclusters or
nanoparticles�, and the objects whose length in the direction
of the applied field is much larger than their transverse size
�nanowires or nanowhiskers�. The following consideration is
carried out for the cases of nanoclusters and nanowires of the
simplest geometries, when the electric fields FN and the ca-
pacitances C and C2 can be determined consistently by solv-
ing corresponding electrostatic problems. The current is cal-
culated according to Eq. �6�, under the assumption that the
objects are in equilibrium with the cathode.
A. Field emission from nanoclusters
Below we consider the field emission from a nanocluster
modeled by a spherical metallic particle of radius R depos-
ited on the flat cathode surface. To provide a finite capaci-
tance C, one should assume a finite separation d-R between
the particle and the metallic cathode plate �for instance, one
can imagine that the particle resides on an oxidized surface�,
see the inset to Fig. 3. Besides, this assumption provides
electrical isolation of the particle from the cathode, which is
FIG. 2. Current from the idealized emitter as a function of the
applied field F0=−V /L for the case of small C /G �nanoparticle in
thermal equilibrium with the cathode, upper curves� and for the
case of C /G=100 ps �lower curves�, at the temperatures T=77 K
�solid� and T=293 K �dashed�. The inset shows the distribution
function PN at F0=5 10
5 V/cm for the second case.
FIG. 3. Charge density per unit length in z direction for a spheri-
cal metallic nanocluster placed at the distance 0.1 R from the me-
tallic cathode. Here
0=F0R /2. The inset shows the geometry of
the problem and the directions of the field emission �arrows�.
a necessary condition for the Coulomb blockade. The elec-
trostatics of the plane-sphere system is known, and the field
and charge distributions in this case can be found in the form
of rapidly converging infinite series arising from the poten-
tials of image point charges and point dipoles.32 Such a con-
sideration allows one to present the distribution of the elec-
trostatic potential energy near the particle in the approximate
U�r,�� � W + e��F0R�1 +
�cos � − 1�
− ��eN − C2V /C�
r − R
, �9�
where r and � are the radial and azimuthal coordinates of the
spherical coordinate system with the origin at the center of
the particle, and �,
, and � are the dimensionless constants
of the order of unity, which are to be determined from nu-
merical calculations. Such calculations also give us the ca-
pacitances C and C2.
33 Note that if the charge quantization is
neglected �so that N=C2V /e when the particle is in equilib-
rium with the cathode�, � is identified with the field enhance-
ment factor conventionally used in the physics of field emis-
sion. The expression �9� provides an excellent description of
the electrostatic potential at r−R�R /2 and at small �. It
allows one to take into account deviations of the potential
energy from the linear form W− �e �F�r−R� and, therefore, to
find corrections to the Fowler-Nordheim tunneling exponent.
Neglecting such corrections in the prefactor, we obtain the
following expression for the partial currents
IN = ASFN
2 exp�− F
�� �e�FNR
��x� =
� x2�x − 1��2 − arctan �x − 1� − x� , �10�
where A is given by Eq. �2�, the dimensionless function ��x�
describes the corrections to the tunneling exponent, and the
effective emitting area S=2�R2�FN
2 /
F�F0�
�2�R2�FN /
F� is reduced due to the angular dependence of
the radial field described by Eq. �9�. The field FN is given by
FN = �F0 + �
�e�N − C̃2F0
, �11�
where the quantity C̃2=C2L does not depend on the distance
L between the cathode and anode. Note that, since we always
assume that L is much larger than any dimension of the
nanoscale object, the capacitance C2 is always proportional
to 1 /L, and it is more convenient to replace C2 �V� by C̃2F0.
This substitution also allows us to represent the Coulomb
energy standing in Eq. �6� as
�N − C̃2F0/�e�
2. �12�
Further calculations are done for the separation d−R
=0.1R, when C=2.16R, C̃2=1.74R
2, �=4.32, �=1.22, and
=0.66. Figure 3 shows the distribution of negative charges
on the surface of the spherical particle staying in equilibrium
with the cathode for this case ��e �N= C̃2F0 is assumed�. The
distribution of the radial field F�z� at the surface of the par-
ticle is given by the same dependence, F�z� /F0=
�z� /
The field-emission current from the nanocluster described
above has been calculated according to Eqs. �6� and �10�–
�12� at R=5 nm. The results of the calculations shown in Fig.
4 demonstrate the staircaselike behavior caused by the Cou-
lomb blockade. However, in contrast to the staircases shown
in Fig. 2, the current continues to increase between the steps.
This occurs because of electrostatic polarization of the nano-
particle. According to Eq. �11�, when the particle charge is
constant, the increase in the applied field F0 leads to an in-
crease in the effective field FN because the factor �
−�C̃2 /CR is positive. For the chosen particle radius, the
steps of the current are clearly visible at liquid nitrogen tem-
perature but poorly visible at room temperature. Neverthe-
less, the Coulomb blockade features at room temperature be-
come quite distinct in the plots of the derivative of the
current, as shown in the inset to Fig. 4.
B. Field emission from nanowires
Let us consider the field emission from a small semicon-
ductor wire modelled by a cylinder of radius R and length d,
which ends with a hemispherical tip of the same radius, see
the inset to Fig. 5. The cathode substrate upon which the
wire is grown is assumed to be a metal �or a heavily doped
semiconductor� so that one can use the method of image
charges instead of solving the electrostatic problem in the
whole space. The electric isolation of the wire from the cath-
ode in this case takes place in a natural way, because a
Schottky barrier is formed between the wire and the metallic
cathode �in the case of semiconducting cathode there can be
a heterobarrier or an interband p-n barrier�. In other words,
FIG. 4. Current from the spherical nanocluster of radius R
=5 nm as a function of the applied field F0=−V /L at T=4.2 K
�solid� and 77 K �dashed�. The inset shows the derivative of the
current at T=293 K.
the wire region adjacent to the cathode becomes depleted of
electrons and positively charged because of the presence
of donors �we assume that the wire is uniformly doped
with bulk donor density nD�. When a bias eV is applied
between the cathode and anode, the wire acquires a con-
siderable negative charge because of tunneling or thermi-
onic emission of electrons from the cathode through the
barrier. When the field emission from the wire of nan-
oscale radius becomes essential, the density of induced
negative charges per unit length of the wire appears to be
much larger than the equilibrium charge density
=�R2 �e �nD even if nD is of the order of 10
18 cm−3. For
this reason, one can use the “metallic” approximation
assuming that the charges in the wire are placed mostly on its
surface. This means that the electron density distribution
n�� ,z�, which depends on the radial coordinate � of the cy-
lindrical coordinate system connected with the wire, is given
by n�� ,z�= �2�� �e � �−1���−R�
�z�+nD for z�d and n�� ,z�
= �2�� �e � �−1���−�R2− �z−d�2
�z�+nD for d�z�d+R,
where
�z� is the density of negative charges on the surface
per unit length. Since this approximation is based on the
assumption that the screening length is small in comparison
with the wire radius, it works better for wider wires. For
silicon wires, whose field-emission properties are currently
the subject of investigations,11–15 the metallic approximation
remains suitable even for the radius of several nanometers,
because, owing to the large effective masses and six-valley
degeneracy, the density of electron states in n-Si appears to
be high enough to provide the Thomas-Fermi screening
length less than one nanometer for Fermi energies �F
�0.01 eV. The metallic approximation, of course, fails to
describe the region of the wire in the close vicinity of the
cathode, where the depletion occurs. Nevertheless, since this
region is a small part of the whole wire, see the charge dis-
tribution in Fig. 5, its presence cannot considerably modify
the parameters calculated as described below.
According to the discussion given here, we search for the
charge distribution
�z� satisfying the integral equation
U�z� = U0 − �e�F0z + �
dz�K�z,z��
�z�� , �13�
where U�z� is the potential energy counted from the Fermi
level in the cathode material, U0 is the barrier height, and
K�z ,z�� is the potential of interaction between the electrons
in the points z and z� of the wire surface in the presence of
the cathode plate, see the Appendix. Equation �13� is accom-
panied with additional requirements: U�z�=0 at z�z0 and
�z�=−
D at z�z0, where z0 is the depletion edge coordi-
nate, which is to be found self-consistently. The first of these
requirements corresponds to a full screening of the bare po-
tential U0− �e �F0z by the induced charges of the wire, while
the second one models the presence of the positive charges in
the depletion region z�z0. Once the distribution
�z� is
found, the total charge of the wire, −�0
d+Rdz
�z�, as well as
the distribution of electric field around the wire, can be cal-
culated. To find the capacitance C and describe modification
of the effective field under single-electron charging, one may
calculate the variation of the total charge and the field at the
tip �at z=d+R� with respect to a small variation of U0. Equa-
tion �13� is solved numerically by using the method of itera-
tions. The dependence of the effective field FN on F0 and N
can be represented in the form similar to Eq. �11�
FN = ��F0�F0 + �
�e��N + B� − C̃2F0
C�F0�R
, �14�
while the Coulomb energy is written as
2C�F0�
�N + B − C̃2F0/�e�
2. �15�
These equations take into account a finite �though weak�
dependence of the capacitance C and field enhancement fac-
tor � on the applied field F0. The dependence of the param-
eters C̃2 and � on F0 appears to be much weaker and can be
neglected. The positive dimensionless constant B reflects the
fact that the average number of induced charges is smaller
than C̃2F0 / �e�. These features appear because the system un-
der consideration is not entirely metallic and contains a
depletion region whose length changes with F0.
The numerical calculations leading to the results pre-
sented below are done for U0=0.7 eV, which approximately
corresponds to the Schottky barrier height for n-Si in contact
with Al.34 The chosen donor density is nD=2 10
18 cm−3.
The parameters standing in Eqs. �14� and �15�, however, are
not sensitive to nD, except for the capacitance C, which
changes within 10% when nD varies from 10
18 cm−3 to 2
1018 cm−3. Figure 5 shows the charge density distribu-
tion for the wire of radius R=5 nm and length d=0.1 �m
at F0=10
6 and 2 106 V/cm. The charge density shows
a nearly linear growth through the main part of the wire
and a sharp enhancement at the hemispherical tip from
which the field emission occurs. The dependence of the
field enhancement factor and capacitance on the applied
electric field is shown in Fig. 6, and the other calcu-
lated parameters are C̃2=2.44 dR, �=0.414, and B=12.14.
FIG. 5. Charge density per unit length for the cylindrical wire
whose geometry is shown in the inset �see parameters in the text�.
The plots of the field-emission current calculated with the
use of the parameters listed here are given in Fig. 7. The
calculations are done according to Eqs. �6�, �14�, and �15�,
and the Fowler-Nordheim formula for the partial current,
IN=ASFN
2 exp�−F /FN�. Since the calculated radial electric
field in the region of the tip weakly depends on z �in contrast
to the case of the nanocluster studied above� and sharply
decreases in the region of transition to the cylindrical part of
the wire, the effective emitting area S is estimated by the
total area of the hemispherical tip, S=2�R2. The work func-
tion is taken for silicon, W=4.2 eV. Next, the Fermi energy
standing in the expression for A, see Eq. �2�, is estimated
from the equation �F��e �FinrTF, where Fin��F0 /� is the
field inside the semiconductor near the end of the tip, rTF is
the Thomas-Fermi screening length, and � is the dielectric
constant of the semiconductor. Such an estimate, carried out
for n-Si, leads to �F�0.1 eV at F0�2 10
6 V/cm. The pic-
ture of Coulomb staircase shown in Fig. 7 is basically the
same as that in Fig. 4. Again, the increase of the current with
the applied field is determined by the increase of the effec-
tive field �14� due to both single-electron charging �steps�
and charge polarization under a constant charge �regions be-
tween the steps�. The main difference is that the interval of
the applied field needed for addition of one electron to the
wire is considerably reduced, owing to the larger capacitance
C2, and appears to be of about 1.2 V/�m �further reduction
of this interval takes place with the increase of the wire
length, see below�. Next, since the capacitance C increases
considerably in comparison to the case of nanocluster of the
same radius, the Coulomb blockade features at room tem-
perature are poorly visible even in the derivative plot, see the
inset. Nevertheless, these features remain pronounced at T
=77 K.
With the increase of the wire length d, the parameters
entering Eqs. �14� and �15� are modified as shown in Fig. 8.
The field enhancement factor and the capacitances increase
nearly in a linear way, while the parameter �, which charac-
terizes relative contribution of charging to the effective field,
slightly decreases �for comparison, the idealized emitter con-
sidered in the previous section is described by the parameters
�=d /R, �=1, C=R, and C̃2=dR, where d is the distance
from the cathode to the emitting sphere�. The increase of the
total capacitance C makes it difficult to observe the Coulomb
staircase in long wires. For example, at d=1 �m one should
have temperatures considerably lower than 77 K. The inter-
val of the applied field corresponding to the addition of one
electron is inversely proportional to C̃2. This interval de-
creases very fast with the increase of d and becomes equal to
2.5 102 V/cm at d=1 �m.
FIG. 6. Field dependence of the enhancement factor and capaci-
tance for the cylindrical wire with R=5 nm and d=0.1 �m.
FIG. 7. Current from the cylindrical wire of radius R=5 nm and
length 0.1 �m as a function of the applied field F0=−V /L at T
=4.2 K �solid� and 77 K �dashed�. The inset shows the derivative of
the current at T=293 K.
FIG. 8. Dependence of the parameters �, �, C, and C̃2 on the
length of the wire for R=5 nm and F0=5 10
5 V/cm.
IV. CONCLUSIONS
The key point of the presented theoretical study is the
possibility of noticeable modification of the effective electric
field causing the field emission from a nanoscale conductor
by addition of just one electron to this conductor. Formally,
this modification is described by introducing the effective
field FN, which determines the partial current IN, and by
evaluating the dependence of this field on the bias applied
between the cathode and anode, see Eqs. �11� and �14�. As a
result of this effect, the current-voltage characteristics of the
field emission show steps in the Coulomb blockade regime.
In other words, the steplike current-voltage characteristics
related to single-electron charging �Coulomb staircases� may
exist even under the conditions of field-emission experi-
ments, when the applied bias is orders of magnitude larger
than the charging energy. The steps on the current-voltage
characteristics can be visible at 77 K in the case of field
emission from nanoclusters and nanowires of 10 nm diam-
eter and submicron length. In the regions between the steps,
where the total charge of the nanoscale object is constant, the
current increases with the increase of the applied bias owing
to charge polarization.
The staircases described in this work are similar to the
usual Coulomb staircases obtained in the transport through
small metallic islands28–30 or quantum dots �see Ref. 35 for
review� with strong asymmetry in the barriers. In both cases,
each step of the current is associated with addition of an
electron to the nanoscale object, and the applied source-drain
voltage drops mostly across the low-transparent barrier �the
barrier between the object and the drain�. Therefore, the pe-
riodicity of the steps in both cases is determined by the
object-drain capacitance C2. However, the steps in the sec-
ond case are formed due to shifts of effective �N-dependent�
electrochemical potential of the object with respect to elec-
trochemical potentials of the source and drain. For this rea-
son, the usual Coulomb staircase shows well-defined steps
when C2 is greater than the object-source capacitance C1,
while in the opposite situation, C1�C2, the steps are sup-
pressed and the current-voltage characteristic approaches to a
linear dependence.29,30 In contrast, in the case described in
this work the steps are formed due to changes in the prob-
ability of Fowler-Nordheim tunneling from the object to the
drain �anode�. That is why the steps are clearly visible under
the condition C1�C2, which is imposed by the field-
emission layout considered in this paper. To summarize, the
sensitivity of the field emission to the number of electrons in
the nanoscale object makes it possible to obtain the Coulomb
staircases under the conditions when such staircases cannot
be observed in the transport through small metallic islands or
quantum dots.
The quantitative consideration has been applied here to
some simple models of the nanoscale objects, whose electro-
static properties necessary for description of the field en-
hancement and charging have been determined consistently.
Consequently, the number of geometrical parameters charac-
terizing the objects has been minimized. For example, the
nanowire has been characterized only by its length d and
radius R. In reality, the geometrical structure of objects is
more complicated. For example, their tips may contain sharp
regions which provide a more efficient field emission. In
fact, high field-emission currents from nanoscale objects are
typically observed at the applied fields of the order of
105 V/cm, which requires the field enhancement factors
much larger than those calculated in this paper. On the other
hand, the presence of sharp tips cannot strongly modify the
capacitances of the objects. The general picture of the single-
electron tunneling under the field-emission regime also re-
mains valid. For possible application to experiments, the
field enhancement due to charging can be described by equa-
tions of the kind of Eqs. �11� and �14�, where � and � should
be considered as parameters to be determined experimen-
tally.
At the present time, there is no experimental evidence of
the Coulomb staircase phenomenon under the field emission.
Though the current-voltage characteristics sometimes show
steplike features, see, for example, Ref. 11, these features are
not regular and, most probably, should be attributed to insta-
bilities of the emission process and burning out of the emit-
ting material. There are numerous reasons which make ob-
servation of the phenomena considered in this paper difficult.
First of all, in most cases the nanoscale objects on the cath-
ode surface form dense arrays. This means that the field
emission takes place from a macroscopic number of objects
which are electrostatically coupled. The charging and field-
emission properties appear to be considerably different36
from those of individual objects. The Coulomb blockade
phenomena in this case should be dramatically suppressed by
the size dispersion of the objects and by the effects of mutual
screening. Investigation of field emission from individual ob-
jects is possible in the cases of metallic nanoclusters8–10 and
carbon nanotubes.26 However, there exists the problem of
electric isolation of these objects from the cathode, which is
one of the necessary conditions for Coulomb blockade. No
special attempts to achieve such an isolation in the field-
emission experiments have been undertaken so far, except
for the nanomechanical system investigated in Ref. 7, where
the electron emission from an isolated Au island to a
submicron-sized electrode has been observed. Most of the
experiments on field emission are carried out at room tem-
perature, though existing experimental techniques also allow
measurements at liquid nitrogen temperature. This means
that the Coulomb blockade phenomena can be observed only
for small-sized objects whose capacitances are low enough
�see the results of Sec. III�. Besides, the interval of the ap-
plied field corresponding to addition of one electron strongly
decreases in the case of emission from long nanowires,
which requires high resolution with respect to field. In sum-
mary, a search for the Coulomb blockade features in the
field-emission current would require a special planning of
experiment. The author hopes that the presented theoretical
study will stimulate experimental investigations in this direc-
tion.
ACKNOWLEDGMENTS
The author is grateful to A. I. Klimovskaya for stimulat-
ing discussions.
APPENDIX: KERNEL OF EQUATION (13)
If z�d and z��d, K�z ,z��=K0�z ,z��−K0�−z ,z��, where
K0�z,z�� = �
��z − z��2 + 2R2�1 − cos ��
. �A1�
If z�d and z��d, K�z ,z��=K0�z ,z��−K0�−z ,z��, where
K0�z,z�� = �
��d − z�2 + 2R�d − z�cos �� + 2R2�1 − sin �� cos ��
. �A2�
If z�d and z��d, K�z ,z��=K0�z ,z��−K0�z ,−z��, where
K0�z,z�� = �
��d − z��2 + 2R�d − z��cos � + 2R2�1 − sin � cos ��
. �A3�
Finally, if z�d and z��d,
K�z,z�� = �
� �e��2R2�1 − cos � cos �� − sin � sin �� cos ����
�4d2 + 4dR�cos � + cos ��� + 2R2�1 + cos � cos �� − sin � sin �� cos ��� . �A4�
In Eqs. �A2�–�A4�, cos �= �z−d� /R and cos ��= �z�−d� /R, so � and �� are the azimuthal angles. The integrals are taken over
the polar angle �. The function K�z ,z�� is also representable in the form of full elliptic integrals.
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|
0704.0416 | Origamis with non congruence Veech groups | Origamis with non congruence Veech groups
Gabriela Schmithüsen
In this article we give an introduction to origamis (often also called square-tiled
surfaces) and their Veech groups. As main theorem we prove that in each genus
there exist origamis, whose Veech groups are non congruence subgroups of SL2(Z).
The basic idea of an origami is to obtain a topological surface from a few combina-
torial data by gluing finitely many Euclidean unit squares according to specified
rules. These surfaces come with a natural translation structure. One assigns in
general to a translation surface a subgroup of GL2(R) called the Veech group. In
the case of surfaces defined by origamis, the Veech groups are finite index sub-
groups of SL2(Z). These groups are the objects we study in this article.
One motivation to be interested in Veech groups is their relation to Teichmüller
disks and Teichmüller curves, see e.g. the article [H 06] of F. Herrlich in the
same volume: A translation surface of genus g defines in a geometric way an
embedding of the upper half plane into the Teichmüller space Tg of closed Rie-
mann surfaces of genus g. The image is called Teichmüller disk. Its projection to
the moduli space Mg is sometimes a complex algebraic curve, called Teichmüller
curve. More precisely this happens, if and only if the Veech group is a lattice in
SL2(R). In this case the algebraic curve can be determined from the Veech group
up to birationality.
It is hard to determine the Veech group for a general translation surface. How-
ever, if the translation surface comes from an origami there is a special approach
to this problem. It is based on the idea of describing origamis by finite index sub-
groups of F2, the free group in two generators. This leads to a characterization of
origami Veech groups as the images in SL2(Z) of certain subgroups of Aut(F2),
the automorphism group of F2.
Using this approach we will calculate Veech groups of two origamis explicitly.
They turn out to be non congruence groups. Starting from these examples we
obtain infinite sequences of origamis all of whose Veech groups are non congruence
groups. This leads to the following theorem.
Theorem 1. Each moduli space Mg (g ≥ 2) contains an origami curve whose
Veech group is a non congruence group.
In Section 1 we introduce origamis and present different equivalent ways to de-
scribe them. In Section 2 we give a glance on the mathematical context. We
describe, how an origami defines a family of translation surfaces and explain
roughly , how one obtains a Teichmüller curve in moduli space starting from
http://arxiv.org/abs/0704.0416v1
an origami. We introduce Veech groups and shortly point out their relation to
Teichmüller curves. In Section 3 we turn to Veech groups of origamis and present
a characterization of them in terms of automorphisms of the free group F2 in two
generators. We use this characterization to calculate two examples explicitly.
Finally, in Section 4 we show that these two examples produce Veech groups that
are non congruence groups and give a method to construct out of them infinite
sequences of Veech groups that are again non congruence groups.
The first part (Section 1 - Section 3) of this article is meant to give a handy in-
troduction to origamis and an overview on some of our results about their Veech
groups. In the second part we state and prove Theorem 1 based on the results in
the PhD thesis [S 05] of the author.
For a broader introduction and overview on origamis and Teichmüller curves
as well as for references to the larger context, we refer the the reader e.g. to
[HeSc 06], [S 04] and [S 05].
Acknowledgments: I would like to thank Frank Herrlich for his support in
respect of the content and for his proof reading, Stefan Kühnlein for helpful dis-
cussions and suggestions especially on non congruence groups and the organizers
of the conference for giving me the opportunity to contribute to these proceedings.
This work was partially supported by a fellowship within the Postdoc-Programme
of the German Academic Exchange Service (DAAD).
1 ORIGAMIS 3
1 Origamis
There are several ways to define origamis. We start with the somehow playful
description that we have learned from [Lo 05], where also the name origami was
introduced: An origami is obtained by gluing the edges of finitely many copies
Q1, . . . , Qd of the Euclidean square Q via translations according to the following
rules:
• Each left edge shall be identified to a right edge and vice versa.
• Similarly, each upper edge shall be identified to a lower one.
• The arising closed surface X shall be connected.
We only study what is called oriented origamis in [Lo 05] and call them just
origamis.
Example 1.1.
a) The simplest example is the origami that is made from only one square.
There is precisely one possibility to glue its edges according to the rules.
One obtains a torus E. We call this origami the trivial origami O0.
Figure 1: The trivial origami. Opposite edges are glued.
Observe that the four vertices of the square are all identified and become
one point on the closed surface E. We call this point ∞.
b) We now consider an origami made from four squares, see Figure 2. Some
identifications of the edges are already done in the picture. For all other
edges those having same labels are glued. The origami is called L(2, 3) for
obvious reasons.
2 3 4
a b c
Figure 2: The origami L(2, 3). Opposite edges are glued.
1 ORIGAMIS 4
Observe that in this case the vertices labeled with • and the vertices labeled
with ◦ are respectively identified and become two points on the closed surface
X. By calculating the Euler characteristic one obtains, that the genus of
the surface X is 2.
c) Finally, we consider an example with five squares, see Figure 3. Here,
edges with same labels are identified. For the unlabeled edges, those which
are opposite to each other are glued. We call the origami D.
1 2 3
• • •
Figure 3: The origami D. Edges with the same label and unlabeled edges
that are opposite are glued.
In this case, we obtain the three identification classes ◦, ⋆ and • for the
vertices. The genus of the closed surface X is again 2.
Origamis as coverings of a torus
Observe, that the trivial origami O0 from Example 1.1 a) is universal in the
following sense: If X is the closed surface that arises from an arbitrary origami
O and E the torus that arises from O0, then we have a natural map X → E by
mapping each of the unit squares of the origami O that form the surface X to
the one unit square of O0 that forms the torus E. This map is a covering that is
unramified except over the one point ∞ ∈ E. Conversely, given a closed surface
X together with such a covering p : X → E, we obtain a decomposition of X
into squares by cutting X along the preimages of the edges of the one square of
O0 that forms E. This motivates the following definition of origamis.
Definition 1.2. An origami O of genus g and degree d is a covering p : X → E
of degree d from a closed, oriented (topological) surface X of genus g to the torus
E that is ramified over at most one marked point ∞ ∈ E.
Remark that we have fixed here one torus E and one point ∞ ∈ E. In particular
we may furthermore fix a point M 6= ∞ on E and a set of standard generators
of the fundamental group π1(E,M) that do not pass through ∞. That way we
obtain a fixed isomorphism
∗) ∼= F2, (1)
1 ORIGAMIS 5
where E∗ = E−∞ and F2 = F2(x, y) is the free group in two generators x and y.
Describing E by gluing the edges of the unit square via translations, we choose
M to be the midpoint of the unit square and the standard generators to be the
horizontal and the vertical simply closed curve through M , see Figure 4.
Figure 4: Generators of π1(E
Example 1.3. In Example 1.1, in a) the covering is the identity id : E → E.
In b) we have a covering p : X → E of degree 4 that is ramified in the two points
labeled by • and ◦. Recall that the genus of X is 2.
In c) we have a covering p : X → E of degree 5 ramified in the two points labeled
by • and ⋆. Observe that though the point on X labeled by ◦ is a preimage of ∞,
the covering is not ramified in this point. The genus of X is again 2.
Definition 1.4. We say that two origamis O1 = (p1 : X1 → E) and O2 =
(p2 : X2 → E) are equivalent, if there is a homeomorphism ϕ : X1 → X2 with
p1 = p2 ◦ ϕ.
Description by a pair of permutations
An origami O = p : (X → E) of degree d defines (up to conjugation in Sd)
• a homomorphism m : F2 = F2(x, y) → Sd or equivalently
• a pair of permutations (σa, σb) in Sd
as follows:
Let M1, . . . , Md be the preimages of the point M (defined as above) under p.
Furthermore, let
m : π1(E
∗,M) → Sym(M1, . . . ,Md)
be the monodromy map defined by p, i.e. for the closed path c ∈ π1(E
∗,M) the
point Mi is mapped to Mj by m(c) if and only if the lift of the curve c to X via
p, that starts in Mi, ends in Mj.
Choosing an isomorphism Sym(M1, . . . ,Md) ∼= Sd and using the isomorphism
∗) ∼= F2 fixed in (1) makes m into a homomorphism from F2 to Sd. We set
σa = m(x) and σb = m(y).
Observe that this homomorphism depends on the chosen isomorphism to Sd and
on the choice of the origami in its equivalence class only up to conjugation in Sd.
Therefore we consider two homomorphisms m1 : F2 → Sd and m2 : F2 → Sd to
be equivalent, if they are conjugated by an element in Sd. Similarly we call two
pairs (σa, σb) and (σ
b) in Sd equivalent, if they are simultaneously conjugated,
i.e. there is some s ∈ Sd such that σa = sσ
−1 and σb = sσ
1 ORIGAMIS 6
Example 1.5. In Example 1.1 we obtain for the origami L(2, 3) in b) the mon-
odromy homomorphism
m : F2 → S4, x 7→ (2 3 4) and y 7→ (2 1),
and thus σa = (2 3 4) and σb = (2 1).
For the origami D in c) we similarly obtain the permutations
σa = (1 2 3) and σb = (1 4 5)(2 3).
Description as finite index subgroups of F2
Origamis can be equivalently described as finite index subgroups of F2, the free
group in two generators, as stated in the following remark. The characterization
of the Veech groups of origamis is mainly based on this observation.
Remark 1.6. We have a one-to-one correspondence:
origamis up to equivalence ↔ finite index subgroups of F2 up to conjugacy.
More precisely, this correspondence is given as follows:
Let O = (p : X → E) be an origami. Define E∗ = E − {∞} and X∗ =
X − p−1(∞). Thus we may restrict p to the unramified covering p : X∗ → E∗.
This defines an embedding of the corresponding fundamental groups:
U = π1(X
∗) ⊆ π1(E
∗) ∼= F2
Again we use the fixed isomorphism in (1), see also Figure 4. Changing the
origami in its equivalence class leads to a conjugation of U with an element in
F2. The index of the subgroup of F2 is the degree d of the covering p.
Conversely, given a finite index subgroup U of F2 we retrieve the origami in
the following way: Let v : Ẽ∗ → E∗ be a universal covering of E∗. By the
theorem of the universal covering, π1(E
∗) is isomorphic to Deck(Ẽ∗/E∗), the
group of deck transformations of Ẽ∗/E∗. Furthermore, the finite index subgroup
U of Deck(Ẽ∗/E∗) corresponds to an unramified covering p : X∗ → E∗ of finite
degree. This can be extended to a covering X → E, where X is a closed surface.
Example 1.7. In Example 1.1, we obtain the following subgroups of F2:
In a), X∗ is the once punctured torus itself and U = F2.
In b), X∗ is a genus 2 surface with 2 punctures. Thus U = π1(X
∗) is a free group
of rank 5. Keeping in mind that we use the identification π1(E
∗) ∼= F2 = F2(x, y)
described in Figure 4, one can read off from the picture in Figure 2 that
U = < x3, xyx−1, x2yx−2, yxy−1, y2 >
In c), X∗ is a genus 2 surface with three punctures. Thus U is a free group of
rank 6. More precisely, we read off the picture in Figure 3, that
U = < x3, xyx−2, x2yx−1, yxy−1, y2xy−2, y3 >
2 TRANSLATION STRUCTURES AND VEECH GROUPS 7
Description as a finite graph
Finitely, sometimes it is convenient to describe an origami O = (p : X → E)
as a finite, oriented labeled graph: Namely, let U be the finite index subgroup
of F2 (unique up to conjugation) that corresponds to O as described in the last
paragraph. Then we represent the origami by the Cayley-Graph of U ⊆ F2:
The vertices of the graph are the coset representatives. They are labeled with a
representative of the coset. The edges are labeled with x and y. For each vertex
(with label w ∈ F2) there is an x-edge from it to the vertex that belongs to the
coset of wx. And similarly there is a y-edge to the vertex that belongs to the
coset wy.
Example 1.8. The following figure shows the Cayley-graph for the origami L(2, 3)
from Example 1.1:
?>=<89:;ȳ
GFED@ABCīd
x // ?>=<89:;x̄ x //
�� GFED@ABC
Figure 5: Graph for O = L(2, 3).
2 Translation structures and Veech groups
Translation structures
Recall that an atlas on a surface is called translation atlas, if all transition maps
are translations. An origamiO = (p : X → E) naturally defines an SL2(R)-family
of translation structures µA (A ∈ SL2(R)) on X
∗ = X − p−1(∞) as follows:
• As first step, observe that each A ∈ SL2(R) naturally defines a translation
structure ηA on the torus E itself by identifying it with C/ΛA, where
and ΛA is the lattice <
> in C (2)
• Then define the translation structure µA on X
∗ by lifting ηA via p, i.e.
µA = p
Using the first description of an origami that we gave by gluing squares, we obtain
the translation structure µI (where I is the identity matrix), if we identify the
squares with the Euclidean unit square in C. We obtain µA for a general matrix
2 TRANSLATION STRUCTURES AND VEECH GROUPS 8
A ∈ SL2(R) from this by identifying the squares with the parallelogram spanned
by the two vectors
Thus the SL2(R)-variations of the translation structure µI can be thought of as
affine shearing of the unit squares, see Figure 6.
Figure 6: Sheared translation structure for the origami L(2, 3).
From an origami to a Teichmüller curve in the moduli space
By the SL2(R)-family of translation structures, the origami O = (p : X → E)
defines a specific complex algebraic curve called Teichmüller curve in the moduli
space Mg of closed Riemann surfaces of genus g. We state this construction here
only briefly as motivation and refer e.g. to the overview article [HeSc 06] for a
detailed description and links to references. A particular nice configuration of
such Teichmüller curves is described in [H 06] in this volume.
The Teichmüller curve in Mg is obtained from the origami in the following way:
• The translation structure µA described in the previous paragraph is in par-
ticular a complex structure on the surface X∗ which can be extended to the
closed surface X . The Riemann surface (X, µA) together with the identity
map id : X → X as marking then defines a point in the Teichmüller space
Tg. Thus we obtain the map: ι̂ : SL2(R) → Tg, A 7→ [(X, µA), id].
• If A ∈ SO2(R), then the affine map z 7→ A ·z is holomorphic. Thus the map
ι̂ factors through SO2(R). Furthermore using that SL2(R) modulo SO2(R)
is isomorphic to the upper half plane H, one obtains a map
ι : H ∼= SO2(R)\SL2(R) → Tg
In fact, this map is an embedding that is in the same time holomorphic and
isometric. A map with this property is called Teichmüller embedding and
its image ∆ in Teichmüller space is called a Teichmüller disk or a geodesic
disk.
2 TRANSLATION STRUCTURES AND VEECH GROUPS 9
• Finally, one may compose the map ι with the projection to the moduli space
Mg. The image of ∆ in Mg is a complex algebraic curve. A curve in Mg
that arises like this as the image of a Teichmüller disk is called Teichmüller
curve.
Note: More generally, one obtains a Teichmüller disk ∆ in a similar way starting
from an arbitrary translation surface (or a bit more general: from a flat surface).
However, the image of such a disk ∆ in moduli space is not always a complex
algebraic curve; in fact its Zariski closure tends to be of higher dimension. It
is an interesting question how to decide whether a translation surface leads to a
Teichmüller curve. One possible answer to this is given by the Veech group which
we introduce in the following paragraph.
Veech groups
Let X∗ be a connected surface and µ a translation structure on it. One assigns to
it a subgroup of GL2(R) called Veech group as described in the following: We con-
sider the group Aff+(X∗, µ) of all orientation preserving affine diffeomorphisms,
i.e. orientation preserving diffeomorphisms that are locally affine maps of the
plane C, see Figure 7. Here – and throughout the whole article – we identify C
with R2 by the map z 7→ (Re(z) , Im(z))t. Thus an affine diffeomorphism f can
be written in terms of local charts as
f : z = (Re(z), Im(z))t 7→ A · (Re(z), Im(z))t+ z0 with A ∈ GL2(R) and z0 ∈ C.
Observe that A does not depend on the chart, since µ is a translation structure.
Thus one obtains a well defined map
D : Aff+(X∗, µ) → GL2(R), f 7→ A
called Derivative map.
Definition 2.1. The Veech group Γ(X∗, µ) of the translation surface (X∗, µ) is
the image of the derivative map D:
Γ(X∗, µ) = D(Aff+(X∗, µ))
z 7→ Az + z0
Figure 7: An affine diffeomorphism of a translation surface
3 VEECH GROUPS OF ORIGAMIS 10
Example 2.2. Let (X∗, µ) be C/ΛI with the natural translation structure. Here
I is the identity matrix and ΛI is the corresponding lattice as defined in (2).
An affine diffeomorphisms of C/ΛI lifts to an affine diffeomorphism of C respect-
ing the lattice. Conversely, each such diffeomorphism descends to C/ΛI. Thus,
we have in this case
Γ(X∗, µ) = SL2(Z).
Veech groups and Teichmüller curves
As indicated in the paragraph about Teichmüller curves, the Veech group “knows”
whether a translation surface defines a Teichmüller curve in moduli space or not.
More precisely, one has the following statement:
Fact: Let X be a surface of genus g and X∗ = X−{P1, . . . , Pn} for finitely many
points P1, . . . , Pn on X . Furthermore let µ be a translation structure on X
Then (X∗, µ) defines a Teichmüller curve C if and only if the Veech group Γ(X∗, µ)
is a lattice in SL2(R). In this case, the curve C is (antiholomorphic) birational
to H/Γ(X∗, µ).
We describe the relation to Teichmüller curves here just as motivation and in
order to give a glance at the general frame. We have therefore resumed theorems
contributed by several authors condensed in what is here called “fact”. A good
access to it can be found e.g. in [EG 97] or [Z 06]. A broader overview on
Veech groups of translation surfaces is given e.g. in [HuSc 01] and in [Le 02].
Teichmüller disks, Teichmüller curves and Veech groups have intensively been
studied by numerous authors, starting from Thurston [T 88] and Veech himself
[V 89]. We refer to [S 04] and [HeSc 06] for more comprehensive overviews on
references.
3 Veech groups of origamis
Let O = p : (X → E) be an origami. We have seen in Section 2 that O defines an
SL2(R)-family of translation structures µA (A ∈ SL2(R)) on X
∗ = X − p−1(∞).
The corresponding Veech groups are not very different. In fact, they are all
conjugated to each other. More precisely, we have:
Γ(X∗, µA) = AΓ(X
∗, µI)A
Thus, we may restrict to the case where A = I which justifies the following
definition.
Definition 3.1. The Veech group Γ(O) of the origamiO is defined to be Γ(X∗, µI).
From Example 2.2 it follows that the Veech group of the trivial origami O0 (de-
fined in Example 1.1) is SL2(Z). For a general origami one can show that Γ(O)
3 VEECH GROUPS OF ORIGAMIS 11
is a finite index subgroup of SL2(Z). In fact, also the converse is true as it was
shown by Gutkin and Judge in [GJ 00]: A Veech group is a finite index subgroup
of SL2(Z) if and only if it comes from an origami.
From this it follows in particular by the Fact presented in Section 2 on page 10
that an origami always defines a Teichmüller curve in the moduli space.
Characterization of origami Veech groups
Recall from Section 1 that an origami O corresponds (up to equivalence) to a
finite index subgroup U of F2 = F2(x, y), the free group in two generators (up to
conjugation). This description enables us to give a characterization of its Veech
group entirely in terms of F2 and its automorphisms.
For this we need the following two ingredients:
• Let β̂ : Aut(F2) → Out(F2) ∼= GL2(Z) be the natural projection. The
fact that we only consider orientation preserving diffeomorphisms applies
to only taking automorphisms of Aut(F2) that are mapped to elements in
SL2(Z). We denote Aut
+(F2) = β̂
−1(SL2(Z)) and restrict to the map
β̂ : Aut+(F2) → SL2(Z).
• Let Stab(U) = {γ ∈ Aut+(F2)|γ(U) = U}
Using these ingredients, it was shown in [S 04] that Veech groups of origamis can
be described as stated in the following theorem.
Theorem 2 (Proposition 1 in [S 04]). For the Veech group Γ(O) of the origami
O holds:
Γ(O) = β̂(Stab(U))
Let us make two comments on this description:
One consequence is, that one obtains an algorithm that can calculate the Veech
group of an arbitrary origami explicitly. This algorithm is described in detail in
[S 04].
As an other consequence, we have now a characterization of all origami Veech
groups as stated in the following corollary.
Corollary 3.2. A finite index subgroup of SL2(Z) occurs as origami Veech group
if and only if it is the image of the stabilizing group Stab(U) ⊆ Aut+(F2) for
some finite index subgroup U in F2.
Thus the question, which finite index subgroups of SL2(Z) are Veech groups be-
comes roughly speaking the same as the question which subgroups of Aut+(F2)
are such stabilizing groups. So far, there is no general answer known.
3 VEECH GROUPS OF ORIGAMIS 12
In [S 05] it was shown that many congruence subgroups of SL2(Z) are Veech
groups. Recall that a congruence group of level n is a subgroup of SL2(Z) that is
the full preimage of some subgroup of SL2(Z/nZ) under the natural homomor-
phism SL2(Z) → SL2(Z/nZ) and n shall be minimal with this property. For prime
level congruence groups the following statement is shown in [S 05, Theorem 4]
Theorem 3. Let p be prime. All congruence groups Γ of level p are Veech groups
except possibly p ∈ {2, 3, 5, 7, 11} and Γ has index p in SL2(Z).
This result is generalized to a statement for arbitrary n in [S 05, Theorem 5]
Presenting the Veech group Γ and the quotient H/Γ for an origami
As mentioned above, using Theorem 2 the Veech group of an origami can be
calculated explicitly. The Veech groups are described as subgroups of SL2(Z)
by generators and coset representatives. We use for the notation that SL2(Z) is
generated by S and T , with
and T =
Recall furthermore from the discussion on Veech groups and Teichmüller curves
in Section 2 on page 10 that for a Veech group Γ we are in particular interested
in the quotient H/Γ, since this quotient is birational to the corresponding Te-
ichmüller curve. Here Γ acts as Fuchsian group on the upper half plane H, which
is endowed with the Poincaré metric.
Since an origami Veech group Γ is a finite index subgroup of SL2(Z), the quotient
H/Γ comes with a natural triangulation. More precisely, we choose the funda-
mental domain for the action of SL2(Z) on H that is the geodesic pseudo-triangle
∆ with vertices P = −1
i, Q = 1
i and P∞ = ∞.
Figure 8: Fundamental domain of SL2(Z).
The surface H/SL2(Z) is obtained by identifying the vertical edges P∞ and Q∞
via T and the edge PQ with itself (with fixed point i) via S.
For an arbitrary subgroup Γ of SL2(Z) of finite index we obtain a fundamental
domain as a union of translates of the triangle ∆: for each coset A we take the
3 VEECH GROUPS OF ORIGAMIS 13
triangle A(∆), where A is a representative of the coset. The identification of
the edges is given by the elements in Γ. Gluing the edges gives the quotient
surface H/Γ, filling in the cusps leads to a closed Riemann surface endowed with
a triangulation. We draw stylized pictures of the fundamental domains that
indicate the triangles (see Figure 9 and 10). The triangles are labeled with a
coset representative, edges that are identified are labeled with the same letter
and vertices that are identified with the same number. Vertices that come from
cusps (i.e. points at ∞) are marked with •.
In particular, one can read off from these stylized pictures the genus and the
number of cusps of the quotient surface H/Γ.
Two examples: the origami L(2,3) and the origami D
The origami L(2,3):
In [S 04, Example 3.5] the Veech group is calculated as follows:
Γ(L(2, 3)) = <
More precisely, one obtains the generators presented as products of S and T as
well as a list of coset representatives.
• List of generators:
= T 3,
= TST 2ST−1T−1,
= TSTST−1S,
= T 2STST−1S−1T−2,
• List of representatives:
I, T, S, T 2, TS, ST, T 2S, TST, T 2ST
Hence, Γ(L(2, 3)) is a subgroup of index 9 in SL2(Z).
The stylized picture of the quotient H/Γ(L(2, 3)) is determined in [S 04, Example
3.6] and is shown here in Figure 9.
3 VEECH GROUPS OF ORIGAMIS 14
TTSTT
Figure 9: Fundamental domain of Γ(L(2, 3)).
From this one can read off that the genus of the quotient H/Γ(L(2, 3)) is 0 and
that it has 3 cusps, namely the vertices labeled by 1,4 and 5. It follows in par-
ticular that the corresponding Teichmüller curve has genus 0.
The origami D:
The Veech group of the origami D is calculated in [S 05, Section 7.3.2]. It has
index 24 in SL2(Z) and the following generators:
= −I, A1 =
= T 3,
= ST 6S−1, A3 =
−7 16
= (T 2S)T 4(T 2S)−1.
= (TS)T 4(TS)−1, A5 =
−20 11
= (TST 2S)T 5(TST 2S)−1,
−18 −5
= (ST 3S)T 2(ST 3S)−1,
The following is a system of cosets representatives:
I , T , S , T 2 , TS , ST , T 2S , TST , ST 2 , STS , T 2ST , TST 2 ,
ST 5 , ST 3 , T 2S , TST 3 , TST 2S , ST 4 , ST 3S , TST 2ST−1 ,
TST 2ST−2 , TST 2ST−3 ; TST 2ST−4 , ST 3ST
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 15
The corresponding origami curve C(D) has genus 0. It is shown with its natu-
ral triangulation in Figure 10. It has six cusps, namely C1, C2, C3, C4, C5 and C6.
TT TTS
TTSTT
STTSTT
TSTTS
TSTTST−1
3ST ST
Figure 10: The origami curve to D.
4 Veech groups that are non congruence groups
Theorem 3 implies that there are many congruence groups which are Veech
groups. How about non congruence groups? In this section we will see that
the Veech groups for the two examples, the origami L(2, 3) and the origami D,
studied in the last paragraph are both non congruence groups. Furthermore,
we give a construction that produces for both of them an infinite sequence of
origamis whose Veech group is a non congruence group. We use this in order to
prove our main theorem.
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 16
An other generalization of the example L(2, 3) was given by Hubert and Lelièvre
in [HL 05], where they show for certain “L-shaped” origamis or square-tiled sur-
faces, how they are called there, that their Veech groups are non congruence
groups. These surfaces are all of genus 2, hence it follows that there are infinitely
many origamis of genus 2 whose Veech group is a non congruence group.
Recall that a group is a congruence group, whose level is a divisor of n, if and
only if it contains the principal congruence group
Γ(n) = {
mod n} = kernel(proj : SL2(Z) → SL2(Z/nZ))
In [S 04, Proposition 3.8] it was shown using a proof of Stefan Kühnlein that the
Veech group of L(2, 3) is a non congruence group. The basic tool for this is the
general level that is defined for any subgroup Γ of SL2(Z) as follows: For each
cusp we define its amplitude to be the smallest natural number n such that there
is an element of Γ conjugated in SL2(Z) to the matrix
which fixes the cusp. Observe that this is equal to the number of triangles around
the vertex that represents the cusp in our stylized picture of the quotient surface
(see Figures 9 and 10). The general level of Γ is the least common multiple of the
amplitudes of all its cusps. A theorem of Wohlfahrt [W 64, Theorem 2] states
that the level and the general level of a congruence group coincide.
The amplitude of the three cusps of H/Γ(L(2, 3)) labeled with 1, 4 and 5 in Fig-
ure 9 is 3, 2 and 4 respectively. Hence, the general level of Γ(L(2, 3)) is 12. Then
it is shown in the proof that Γ(L(2, 3)) does not contain Γ(12) which gives the
contradiction.
The same method can be used in order to show that Γ(D) is a non congruence
group. We here carry out the proof for it. Observe from Figure 10 that the six
cusps C1, . . . , C6 have the amplitude 3, 6, 4, 4, 5 and 2, respectively. Thus the
general level is 60.
Proposition 4.1. The Veech group Γ(D) is a non congruence group.
Proof. Suppose that Γ = Γ(D) is a congruence group. Since the general level of
Γ is 60, we have by the theorem of Wohlfahrt mentioned above, that Γ(60) is a
subgroup of Γ.
We will use the following facts, which can be checked e.g. in Figure 10:
∈ Γ, A6 =
−18 −5
∈ Γ and T =
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 17
In order to verify this in Figure 10, use that
A1 = T
3 and A6 = S
−1T 2S−1T−1S−1TS−1T−3S−1.
We will find an element in Γ whose projection to SL2(Z/60Z) is equal to that of
T , which gives us the desired contradiction.
Recall that
SL2(Z/60Z) ∼= SL2(Z/4Z)× SL2(Z/3Z)× SL2(Z/5Z).
We identify in the following these two groups. Furthermore we denote by p4, p3,
p5 and p60 the projection from SL2(Z) to SL2(Z/4Z), SL2(Z/3Z), SL2(Z/5Z) and
SL2(Z/60Z), respectively. Then p60 = p4 × p3 × p5.
We have
p60(A1) = (
) and
p60(A6) = (
The order of p4(A1) in SL2(Z/4Z) is 4, the order of p3(A1) in SL2(Z/3Z) is 1
and the order of p5(A1) in SL2(Z/5Z) is 5. We also say: The order of p60(A1) is
(4, 1, 5). Since 7 ≡ 3 mod 4 and 7 ≡ 2 mod 5 we have
p60(A
1) = (
) (4)
Furthermore:
p60(A
6) = (
and with the same notation as above p60(A
6) has the order (1, 3, 5). Thus
p60(A
6 ) = (
). (5)
From (4) and (5) it follows that
p60(A
6 · A
1) = (
) = p60(
) = p60(T )
But A206 · A
1 ∈ Γ and T /∈ Γ, thus Γ(60) = ker(p60) cannot be contained in Γ.
Therefore, Γ cannot be a congruence group of level 60. Contradiction!
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 18
Sequences of origamis with non congruence Veech groups
Starting from the origamis L(2, 3) and D we will define respectively a sequence
On, such that for each n ∈ N the Veech group Γ(On) again is a non congruence
group. The basic idea is to “copy and paste”: we will cut the origami along a
segment, take n copies of it and glue them along the cuts.
In Figure 11 we show the origami On for L(2, 3):
1 3 4
5 7 8
. . . 4n-7 4n-5 4n-4
4n-3 4n-1 4n
Figure 11: n copies of L(2, 3). Opposite edges are glued.
Using the description of an origami by a pair of permutations from Section 1, On
is given as:
σa = (1 3 4 5 7 8 9 11 12 . . . 4n−3 4n−1 4n), σb = (1 2)(5 6) . . . (4n−3 4n−2).
Observe that the genus of On is n + 1 and it has 2n cusps: n of order 3 (all n
marked by • in Figure 11), and n of order 1 (all n marked by ◦ in Figure 11).
Finally, we want to present the origami On by the finite index subgroup Hn =
∗) of F2, that corresponds to On by Remark 1.6.
Recall from Example 1.7 that for O1 = L(2, 3), we obtain the free group of rank 5:
U = H1 =< g1 = x
3, g2 = xyx
−1, g3 = x
2yx−2, g4 = yxy
−1, g5 = y
2 >= F5.
The group Hn is obtained as as follows:
Hn = < g
1 , g
1 gj g
1 ∈ F5 | i ∈ {0, . . . , n− 1} and j ∈ {2, . . . , 5} >
In Figure 12, we show the origami Dn:
1 2 3
6 7 8
. . . . . . . . .
5n-4 5n-3 5n-2
b1 a2
b2 an
• • •
* • • •
• • •
Figure 12: n copies of D. Edges with the same label or
unlabeled opposite edges are glued.
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 19
The pair of permutations describing Dn is:
σa = (1 2 3 6 7 8 . . . 5n− 4 5n− 3 5n− 2),
σb = (1 4 5)(6 9 10) . . . (5n− 4 5n− 1 5n)(2 3)(7 8) . . . (5n− 3 5n− 2)
The genus of Dn is 2n and it has n+ 2 cusps: 2 of order 2n (marked as • and ⋆)
and n of order 1 (all n marked by ◦).
Again, we present On by the corresponding finite index subgroup Hn of F2. We
have from Example 1.7 that U = H1 = F6, the free group of rank 6:
U =< g′1 = x
3, g′2 = xyx
−2, g′3 = x
2yx−1, g′4 = yxy
−1, g′5 = y
2xy−2, g′6 = y
3 > = F6
And similarly as above, we obtain:
Hn = < g
1 , g
1 ∈ F6 | i ∈ {0, . . . , n− 1} and j ∈ {2, . . . , 6} >
We will see in the following that for both sequences all Veech groups Γ(On) are
non congruence groups. More precisely, we will show:
Proposition 4.2. For both sequences On the following holds:
• Γ(On) ⊆ Γ(O1), which is for both sequences a non congruence group.
• More generally one has:
n divides m ⇒ Γ(Om) ⊆ Γ(On).
• Different origamis in one sequence have different Veech groups, i.e.:
Γ(On) 6= Γ(Om) for n 6= m.
To prove this, let us detect that we are in the following more general setting.
Setting A:
• Let U be a finite index subgroup of F2. Then U is a free group of rank k
for some k ≥ 2, i.e.
U = < g1, . . . , gk > = Fk
• Let α : Fk → Z be the projection w 7→ ♯g1w
where ♯g1w is the number of g1 in the word w = w(g1, . . . , gk) with g
counted as −1.
• Let Hn be the kernel of pn ◦ α, where pn : Z → Z/nZ is the natural
projection, i.e.
Hn = < g
1 , g
1 gj g
1 ∈ Fk | i ∈ {0, . . . , n− 1} and j ∈ {2, . . . , k} > .
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 20
• Finally, let H0 be the kernel of α, i.e.:
Hn = << g2, . . . , gk >>U ,
is the normal subgroup in U generated by g2, . . . , gk.
Observe that we are in this setting with
U = π1(X
∗) =< x3, xyx−1, x2yx−2, yxy−1, y2 > for the origami L(2, 3) and
U = π1(X
∗) =< x3, xyx−2, x2yx−1, yxy−1, y2xy−2, y3 > for the origami D.
In order to prove the properties in Proposition 4.2, we will need that U fulfills
the following a bit technical condition:
Property B: Let U =< g1, . . . , gk > (k ≥ 2) be as above a finite index
subgroup of F2 of rank k and {wi}i∈I a system of coset representatives with
w1 = id. Suppose that U has the following property:
∀ j ∈ I − {1} : wj << g2, . . . , gk >>U w
j 6⊆ U.
One can check by hand that for both origamis, L(2, 3) and D, this property is
fulfilled. In this setting we obtain the following conclusions.
Proposition 4.3. Let n ∈ N ∪ {0}. Let U be a finite index subgroup of F2
fulfilling property B. With the notations from Setting A, we have:
a) The normalizer of Hn in F2 is equal to U : NormF2(Hn) = U
b) Stab
(Hn) ⊆ StabAut+(F2)
c) Recall that U = Fk, the free group in k generators.
Let βn : Aut(Fk) → GLk(Z/nZ) be the natural projection.
Then Stab
(Hn) is equal to
β−1n ({A = (ai,j)1≤i,j≤k ∈ GLk(Z/nZ)| a1,2 = . . . = a1,k = 0}) ∩ G.
Here we use the notation Z/(0Z) = Z thus β0 is the natural projection
Aut(Fk) → GLk(Z).
Proof.
By definition Hn is normal in U , i.e. U ⊆ NormF2(Hn).
Let now w be an element of F2\U . Hence, w = wj ·u for some j ∈ I−{1}, u ∈ U .
By Property B, there exists some h0 ∈ << g2, . . . , gk >>U = H0, such that
wjh0w
j 6∈ U . Therefore we have w(u
−1h0u)w
−1 6∈ U . But u−1h0u ∈ H0 ⊆ Hn,
since H0 is normal in U . This shows that w 6∈ NormF2(Hn).
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 21
This follows from a), since for a subgroup H of F2 in general holds:
(H) ⊆ Stab
(NormF2(H)), see e.g. [S 06, Remark 3.1].
Define M = {A = (ai,j)1≤i,j≤k ∈ GLk(Z/nZ)| a1,2 = . . . = a1,k = 0}.
Let γ ∈ G. We have to show that γ(Hn) = Hn if and only if βn(γ) ∈ M .
Let furthermore pkn : Fk → (Z/nZ)
k be the natural projection.
Consider the following commutative diagram:
Hn = p
n(Hn) ⊆ (Z/nZ)
βn(γ) // (Z/nZ)k ⊇ pkn(Hn) = Hn
Since pkn is surjective and Hn is the full preimage of Hn = p
n(Hn), it follows that
γ(Hn) = Hn if and only if βn(γ)(Hn) = Hn.
Observe finally that:
Hn = {(0, x2, . . . , xk) ∈ (Z/nZ)
k} and
StabGLk(Z/nZ)(Hn) = { A = (ai,j)1≤i,j≤k ∈ GLk(Z/nZ)|
(y1, . . . , yk) = A · (0, x2, . . . , xk) ⇒ y1 = 0 }
= {A = (ai,j) ∈ GLk(Z/nZ)| a1,2 = . . . = a1,k = 0}
Theorem 2 suggests the following notation.
Definition 4.4. Let U be a subgroup of F2.
With β̂ : Aut+(F2) → SL2(Z) as in Theorem 2, we define
Γ(U) = β̂(Stab
and call Γ(U) the Veech group of U .
We now obtain from Proposition 4.3 the following conclusions.
Corollary 4.5. Suppose that we are in the same situation as in Proposition 4.3,
in particular that U is a finite index subgroup of F2 fulfilling property B. Then
we have for all n ∈ N:
a) Stab
(H0) ⊆ StabAut+(F2)
(Hn) and Γ(H0) ⊆ Γ(Hn).
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 22
b) If m ∈ N with n|m, then:
(Hm) ⊆ StabAut+(F2)
(Hn) and Γ(Hm) ⊆ Γ(Hn).
(H0) =
(Hn) and Γ(H0) =
Γ(Hn)
Proof.
a) and b):
Let γ ∈ G. By Proposition 4.3 we have that
∀n ∈ N : γ ∈ Stab
(Hn) ⇔ βn(γ) = A = (ai,j)
with a1,2 ≡ . . . ≡ a1,k ≡ 0 mod n
and γ ∈ Stab
(H0) ⇔ β0(γ) = A = (ai,j)
with a1,2 = . . . = a1,k = 0.
Thus we have for all n ∈ N and for all m ∈ N with n|m, that
(H0) ⊆ StabAut+(F2)(Hm) ⊆ StabAut+(F2)(Hn).
We have in particular by the definition of the Veech group of a subgroup of F2:
Γ(H0) ⊆ Γ(Hm) ⊆ Γ(Hn).
⊆ follows from a). ⊇ follows from Remark [S 06, Remark 3.1].
We now return to the language of origamis: Let O be an origami, U the corre-
sponding subgroup of F2. Define for U the subgroups Hn (n ∈ N) as in Setting
A and let On be the origamis corresponding to the groups Hn.
By Corollary 4.5 and Theorem 2 we obtain immediately the following result.
Proposition 4.6. If U has the Property B, then
∀n ∈ N : Γ(On) ⊆ Γ(O) and ∀n,m ∈ N : n|m ⇒ Γ(Om) ⊆ Γ(On).
In particular, if Γ(O) is a non congruence group, each Γ(On) is a non congruence
group. Thus in this case, we obtain infinitely many origamis whose Veech group
is a non congruence group.
In order to conclude Proposition 4.2, it is now just left to prove the last item. But
this follows , since we have (see [S 05]) for both sequences On, the one coming
from the origami L(2, 3) and the one coming from the origami D, that
∈ Γ(On) ⇔ 3n divides s. (6)
4 VEECH GROUPS THAT ARE NON CONGRUENCE GROUPS 23
This finishes the proof of Proposition 4.2.
Furthermore, Theorem 1 follows from Proposition 4.2.
Remark: From Corollary 4.5 and (6) it follows that Γ(H0) has infinite index in
SL2(Z). Furthermore it is non trivial, since it contains
for L(2, 3) respectively B3 =
for D.
REFERENCES 24
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Origamis
Translation structures and Veech groups
Veech groups of origamis
Veech groups that are non congruence groups
|
0704.0417 | Thermodynamic Stability - A note on a footnote in Ruelle's book | UWThPh-2007-09
Thermodynamic Stability – A note on a
footnote in Ruelle’s book
Bernhard Baumgartner1
Institut für Theoretische Physik, Universität Wien
Boltzmanngasse 5, A-1090 Vienna, Austria
April 3, 2007
Abstract
Thermodynamic stable interaction pair potentials which are not of the
form “positive function + real continuous function of positive type” are
presented in dimension one. Construction of such a potential in dimension
two is sketched. These constructions use only elementary calculations. The
mathematical background is discussed separately.
PACS numbers: 05.20.-y, 02.20.-a, 02.40.Ft
Keywords: thermodynamic stability, convex cone
1 Introduction
In Ruelle’s book [R69] on statistical mechanics, in section 3.2 concerning one
species of classical particles in Rν , you can read:
1 PROPOSITION. If the pair potential Φ can be written in the form
Φ = Φ1 + Φ2 (1)
where Φ1 is positive, and Φ2 is a real continuous function of positive type, then
Φ is stable.
[email protected]
http://arxiv.org/abs/0704.0417v1
TDS BB April 3, 2007 2
“Positive” is meant here and throughout this paper as nowhere negative,
“stable” means
∃E0 ∈ R such that ∀N, ∀{x1...xN} ⊂ Rν : U(x1 · · ·xN) ≧ N · E0, (2)
where U(x1 · · ·xN ) =
i 6=j
Φ(xj − xi). (3)
This proposition is accompanied by the
2 FOOTNOTE. It seems to be an open problem to construct a stable potential
which is not of the form (1).
We solve this problem in dimension 1, considering particles either in Z or in
R, giving a detailed proof. In dimension 2 the problem can also be solved, but
we give only a sketch of the ideas. 1
To make it simple, we consider only pair potentials which are bounded con-
tinuous functions and state the stability property as
3 DEFINITION. A bounded continuous real valued function V on Rν is stable,
E(ρ) :=
ρ(x)V (x− y)ρ(y)dνx dνy ≥ 0 (4)
for every positive finite measure ρ(x)dνx on Rν. A bounded real valued function
V on Zν is stable, if
E(ρ) :=
ρ(~m)V (~m− ~n)ρ(~n) ≥ 0 (5)
for every positive bounded function ρ(~m) on Zν .
The stability property used in Ruelle’s Theorem is an immediate consequence.
With ρ =
i=1 δ(xi − xj) put into equation (4) one gets
U(x1 · · ·xN ) = E(ρ)−N · V (0) ≥ −N · V (0).
The main result of our considerations is stated as
4 THEOREM. Each of the following functions is a stable pair potential, but
not a sum of a positive and a real valued positive definite function.
1. The function V : Z → R, defined as
V (0) = V (2) = V (−2) = 1, V (1) = V (−1) = −1, (6)
V (n) = 0 ∀n with |n| ≥ 3,
2. The function W : R → R, defined as
W (x) =
V (n)f(n− x+ y)f(y)dy, (7)
with f a positive continuous function (−1
) → R and V as defined in (6).
1 Construction in higher dimensions is still an open problem.
TDS BB April 3, 2007 3
2 Properties of the interaction potentials
Proof. Of part (1) of 4 Theorem.
Denote the distribution of particles on the chain by the “density” ρ, a function
Z → Z+. The interaction energy U becomes smaller, when the system is cut into
non-interacting pieces: If ρ(n) ≥ ρ(n+1) divide the chain, cutting between n+1
and n+ 2. Moving the pieces apart, one looses the energy
2[ρ(n)− ρ(n + 1)]ρ(n+ 2) + 2ρ(n+ 1)ρ(n + 3) ≥ 0.
The symmetric procedure of cutting between n − 2 and n − 1 lowers the energy
if ρ(n− 1) ≤ ρ(n).
Now there remains a set of pieces of no more than three lattice points, with
densities like
0 ≤ ρ(n− 1) ≤ ρ(n) ≥ ρ(n + 1) ≥ 0.
Including the “self-energies” N · V (0) one gets for each piece, centered around n,
E = ρ(n−1)2+ρ(n)2+ρ(n+1)2+2[ρ(n−1)ρ(n+1)−ρ(n−1)ρ(n)−ρ(n)ρ(n+1)]
= [ρ(n− 1)− ρ(n) + ρ(n+ 1)]2 ≥ 0.
Proving the stability of V .
If V were the sum of a positive and a positive definite function, it would give
V (n)µ(n) ≥ 0, (8)
for each µ being both positive and positive definite. Now consider
µ(5ν) = 1, µ(5ν ± 1) =
, µ(5ν ± 2) = 0, (9)
which is obviously positive. Positive definiteness is seen by using Bochner’s the-
orem [RN55] and calculating the Fourier-Transform, with α ∈ (−π,+π]:
µ̂(α) =
µ(n)e−inα
5 δ(α) +
),+δ(α+
> 0. (10)
But it does not give a positive value in (8):
V (n)µ(n) = 2−
5 < 0.
TDS BB April 3, 2007 4
The appearance of the numbers 5 and
5 may seem mysterious. Demystifying
is the next section, where we present the “origin” of these V and µ.
In this section we develop further use of these functions in R and in R2.
Proof. Of part (2) of 4 Theorem.
For N particles at x1 . . . xN consider the measure
ρ(x) =
δ(x− xj). (11)
Adding the self-energies N ·W (0), we study
ρ(x)W (x− y)ρ(y) dxdy
V (n)
ρf (x+ n)ρf (x) dx, (12)
with ρf(x) :=
f(x − y)ρ(y) dy. Splitting the integral in (12) into pieces of
intervals with unit length and defining ρf,x(m) = ρf(x+m) gives
V (n)ρf (x+m+ n)ρf(x+m)
ρf,x(p)V (p−m)ρf,x(m) ≥ 0,
by part (1) of the theorem. So the potential W is stable.
Now consider the distribution
µD(x) =
µ(m)δ(x−m), (13)
using the sequence µ defined in (9). This distribution is positive and positive
definite, as can be seen at its Fourier transform, which is (up to a factor) the
same as in (10), now with µ̂D(α+2π) = µ̂D(α) periodically extended to all α ∈ R.
This µD is used to show that the potential is not a sum of positive and positive
definite functions:
W (x)µD(x)dx
V (n)
∫ + 1
dx δ(x−m)f(n− x+ y)f(y)
V (n)µ(n) ·
f 2(y)dy < 0. (14)
In the last step the final support of f is essential.
TDS BB April 3, 2007 5
Construction of a stable pair potential in R2 being a function of the particle
distances only may be done in the following way:
• Use W (x) defined in (7), now with an f supported on (−1
), convolute it
twice with the distribution
h(x) =
e−ǫ|n|δ(x− 5n) :
W1(x) =
h(x− y)W (y − z)h(z)dy dz.
• Take the mean value (times 2π) of all rotated versions: Wr(~x) = 1rW1(|~x|).
• Smoothen out Wr with a positive continuous function g(r) with support on
[0, 1
W2(~x) =
g(|~x− ~y|)
W1(|~y − ~z|)
|~y − ~z|
g(|~z|)d2y d2z.
That the stability is not destroyed by the double convolution with h follows
from a consideration as it is used in the equation (12). Written in a formal way:
〈ρ|W1 |ρ〉 = 〈ρ| h ∗W ∗ h− |ρ〉 = 〈ρ ∗ h|W |ρ ∗ h〉.
Considering only smooth densities ρ(~x) one may take W1(x1)δ(x2) as a stable
distribution in R2:
〈ρ|W1 · δ |ρ〉dim=2 =
〈ρy|W1 |ρy〉dim=1 d y ≥ 0.
Now rotating the axes and taking the mean value does not destroy the stability.
Once more a double convolution is done, now with g in order to get W2 as a
bounded continuous potential acting in R2.
〈ρ|W2 |ρ〉 = 〈ρ| g ∗W ∗ g− |ρ〉 = 〈ρ ∗ g|W |ρ ∗ g〉 ≥ 0.
Smoothing by convolution with g enables to consider again sets of particles rep-
resented by delta-functions in ρ.
To disprove the possibility of splitting W2 into a sum of a positive and a
positive definite function one may use the µD of equ. (13) embedded into R
µD(x, y) = µD(x)δ(y).
Due to the smoothing of Wr by g and due to its decrease given by the decrease
of h, the integral
W2µD is finite:
W2µD(x) dx = W2(0) + 2W2(1)µ(1) + 2 ·
W2(|5ν + n|)µ(n)
TDS BB April 3, 2007 6
The bounded support of f and g is needed here as it was in equ. (14). The
exponential decrease implies
W2(|5ν + n|) = const. · e−5ǫ ν
V (n) ·
1 +O(
The “const.” factor involves the integrals over f 2 and g2, the error term O( 1
gives the difference between e−5ǫ ν/5ν and e−ǫ (5ν+n)/(5ν + n). The summations
over ν and n give
≈ 2 · const. ·
V (n)µ(n) · log(1/ǫ) +O
e−5ǫ ν
The first part is negative and increases without limit when ǫ → 0, while the other
term remains finite. So W2 with small ǫ can not be a sum of positive and positive
definite functions.
3 Mathematical background
Only in applying Proposition 1 in statistical mechanics the Thermodynamic Limit
is considered, not yet in the investigations of “stability”. Moreover, in the re-
formulation in 3 Definition there is no mentioning of “particles”. What is used
of properties of space are: A distance relation between points and an invariant
measure. This allows for a more general version of the definition, concerning
functions on groups. We keep the notation we used above: x and y are elements
of the group, their “group product” is x+ y, the “inverse” of x is −x.
5 DEFINITION. Consider a bounded continuous real valued function V on a
locally compact abelian group G which has the Haar measure dx. V is stable, if
〈ρ|V |ρ〉 :=
ρ(x)V (x− y)ρ(y)dx dy ≥ 0 (15)
for every finite positive Borel measure ρ(x)dx.
Stable functions can be added, multiplied by positive numbers, and limits
may be formed. So they form a closed convex cone, which we call STB. This
cone STB contains POS, the cone of positive functions, also PDF, the cone of
positive definite functions and sums thereof.
STB ⊃ POS + PDF (16)
An investigation of the relations between these cones may proceed via investi-
gation of the dual cones (see [V64, R62, G03]). The dual cones are subsets of
V ′, the space of finite Borel measures µ(x)dx, which is the dual space to V, the
TDS BB April 3, 2007 7
Banach space of bounded continuous functions. The dual cone to POS is POS′,
the set of finite positive Borel measures, dual to PDF is PDF′, the set of finite
positive definite Borel measures. The cone STB′ is given as the closure of the
cone of convex combinations of “correlation measures”
µ(x) =
ρ(x)ρ(y + x)dy, (17)
i.e. convolutions of finite positive Borel measures ρ(x)dx with their reflected ver-
sion ρ(−x)dx. These correlation measures are both positive and positive definite:
STB′ ⊂ POS′ ∩ PDF′ (18)
Now the question of equality or inequality in this relation is related to the
central problem which is our concern in this investigation, the question of equality
or inequality in (16). If the closed cone POS′∩PDF′ contains an element µ which
is not in the closed cone STB′, then, by definition of “dual cone”, there exists
an element V ∈ STB such that
V µ < 0, incompatible with a decompostion
V = f + g, f ∈ POS, g ∈ PDF.
For the groups Z2, Z3, Z4 there is equality in the equations (16) and (18), but
not for Z5.
6 PROPOSITION. The intersection of POS′ ∩ PDF′ with the plane
{(µ(−2) . . . µ(2))|µ(0) = 1} is completely characterized by its extremal points
(0, 0, 1, 0, 0), (0, γ, 1, γ, 0), (γ, 0, 1, 0, γ), (1, 1, 1, 1, 1), with γ = (
5 − 1)/2 =
1/(2| cos 4π/5|).
Proof. By using Bochner’s theorem and analyzing the Fourier transform
µ̂(k) =
µ(n)e−2π k n/5. (19)
On the other hand there is a bound for STB′ which cuts off a triangular subset
of this convex quadrangle:
7 LEMMA. Each element of STB′ obeys the inequality
µ(1) ≤
µ(n)/4. (20)
Proof. STB′ is defined by its extremal rays, formed as correlation measures of
positive densities.
µ ∈ STB′, µ extremal ⇔ ∃ρ ≥ 0, µ(n) =
ρ(m)ρ(m + n).
TDS BB April 3, 2007 8
Assume, w.l.o.g., that ρ(−1) ≥ ρ(−2). then
µ(1) = [ρ(−1) + ρ(1)] · [ρ(−2) + ρ(0) + ρ(2)]− [ρ(−1)− ρ(−2)] ρ(2)− ρ(−2)ρ(1)
− x)(s
+ x) ≤ s
Here s =
m ρ(m), x = [ρ(−2) + ρ(0) + ρ(2)− ρ(−1)− ρ(1)] /2.
Observe
n µ(n) = s
Remark: Also µ(2) obeys this inequality and µ(−1) = µ(1), µ(−2) = µ(2).
Closer inspection reveals moreover two rounded edges of STB′.
Now the extremal point with µ(n) as in equation (9) with ν = 0 is outside
this boundary. And V (n) as in equation (6) is an element of STB, but outside of
POS+PDF.
4 Conclusion
For pair potentials which are bounded continuous functions the property of being
“stable” can be reformulated without mention of particles. In this way it can be
studied for abstract abelian groups. At the heart of the present investigation is
the observation of a function V in Z5 which is stable, but indecomposable into a
sum of positive and positive definite functions. This function V can also be used
on Z. With some smoothing it can be used on R, and in damped periodically
extended, rotationally symmetrized and again smoothed form on R2. Of course
it is possible find sets of other examples nearby. So V (−1) = V (1) in Theorem 4
could be a little bit higher than −1. Only at −(
5+1)/4 ≈ −0.8 does it become
decomposable.
The construction of a rotationally invariant example for dimension two is not
so simple. A nicer one, or one for higher dimension, is not yet known.
References
[R69] D. Ruelle: Statistical Mechanics: Rigorous Results (W. A. Benjamin, inc.,
New York) 1969.
[RN55] F. Riesz and B. Sz. Nagy: Functional Analysis (Ungar, New York) 1955
[V64] Frederick A. Valentine: Convex sets (McGraw-Hill, NY (McGraw-Hill se-
ries in higher mathematics)) 1964
[R62] W. Rudin: Fourier Analysis on Groups (Interscience, New York) 1962
[G03] “Convex cones and their faces” Chapter 3 in: H. Glöckner: Positive Def-
inite Functions on Infinite-Dimensional Convex Cones; Memoirs AMS,
166, Number 789, 2003
Introduction
Properties of the interaction potentials
Mathematical background
Conclusion
|
0704.0418 | Entanglement entropy at infinite randomness fixed points in higher
dimensions | Entanglement entropy at infinite randomness fixed points in higher dimensions
Yu-Cheng Lin1, Ferenc Iglói2,3 and Heiko Rieger1
Theoretische Physik, Universität des Saarlandes, 66041 Saarbrücken, Germany
Research Institute for Solid State Physics and Optics, H-1525 Budapest, Hungary
Institute of Theoretical Physics, Szeged University, H-6720 Szeged, Hungary
(Dated: November 3, 2018)
The entanglement entropy of the two-dimensional random transverse Ising model is studied with
a numerical implementation of the strong disorder renormalization group. The asymptotic behavior
of the entropy per surface area diverges at, and only at, the quantum phase transition that is
governed by an infinite randomness fixed point. Here we identify a double-logarithmic multiplicative
correction to the area law for the entanglement entropy. This contrasts with the pure area law valid
at the infinite randomness fixed point in the diluted transverse Ising model in higher dimensions.
PACS numbers: Valid PACS appear here
Extensive studies have been devoted recently to under-
stand ground state entanglement in quantum many-body
systems [1]. In particular, the behavior of various entan-
glement measures at/near quantum phase transitions has
been of special interest. One of the widely used entan-
glement measures is the von Neumann entropy, which
quantifies entanglement of a pure quantum state in a bi-
partite system. Critical ground states in one dimension
(1D) are known to have entanglement entropy that di-
verges logarithmically in the subsystem size with a uni-
versal coefficient determined by the central charge of the
associated conformal field theory [2]. Away from the crit-
ical point, the entanglement entropy saturates to a finite
value, which is related to the finite correlation length.
In higher dimensions, the scaling behavior of the entan-
glement entropy is far less clear. A standard expectation
is that non-critical entanglement entropy scales as the
area of the boundary between the subsystems, known as
the ”area law” [3, 4]. This area-relationship is known to
be violated for gapless fermionic systems [5] in which a
logarithmic multiplicative correction is found. One might
suspect that whether the area law holds or not depends
on whether the correlation length is finite or diverges.
However, it has turned out that the situation is more
complex: numerical findings [7] and a recent analytical
study [8] have shown that the area law holds even for
critical bosonic systems, despite a divergent correlation
length. This indicates that the length scale associated
with entanglement may differ from the correlation length.
Another ongoing research activity for entanglement in
higher spatial dimensions is to understand topological
contributions to the entanglement entropy [9].
The nature of quantum phase transitions with
quenched randomness is in many systems quite different
from the pure case. For instance, in a class of systems
the critical behavior is governed by a so-called infinite-
randomness fixed point (IRFP), at which the energy scale
ǫ and the length scale L are related as: ln ǫ ∼ Lψ with
0 < ψ < 1. In these systems the off-critical regions
are also gapless and the excitation energies in these so-
called Griffiths phases scale as ǫ ∼ L−z with a nonuni-
versal dynamical exponent z <∞. Even so, certain ran-
dom critical points in 1D are shown to have logarith-
mic divergences of entanglement entropy with universal
coefficients, as in the pure case; these include infinite-
randomness fixed points in the random-singlet universal-
ity class [12, 13, 14, 15, 16] and a class of aperiodic singlet
phases [17].
In this paper we consider the random quantum Ising
model in two dimensions (2D), and examine the disorder-
averaged entanglement entropy. The critical behavior of
this system is governed by an IRFP [10, 11] implying that
the disorder strength grows without limit as the system is
coarse grained in the renormalization group (RG) sense.
In our study, the ground state of the system and the
entanglement entropy are numerically calculated using a
strong-disorder RG method [18, 19], which yields asymp-
totically exact results at an IRFP. To our knowledge this
is the first study of entanglement in higher dimensional
interacting quantum systems with disorder.
The random transverse Ising model is defined by the
Hamiltonian
H = −
〈i,j〉
i . (1)
Here the {σαi } are spin-1/2 Pauli matrices at site i of an
L × L square lattice with periodic boundary conditions.
The nearest neighbor bonds Jij(≥ 0) are independent
random variables, while the transverse fields hi(≥ 0) are
random or constant. For a given realization of random-
ness we consider a square block A of linear size ℓ, and
calculate the entanglement between A and the rest of
the system B, which is quantified by the von Neumann
entropy of the reduced density matrix for either subsys-
tems:
S = −Tr(ρA log2 ρA) = −Tr(ρB log2 ρB). (2)
The basic idea of the strong disorder RG (SDRG) ap-
proach is as follows [18, 19]: The ground state of the sys-
tem is calculated by successively eliminating the largest
http://arxiv.org/abs/0704.0418v2
FIG. 1: (color online). An example of typical ground state
in the random quantum Ising model (a) in 1D, and (b) in 2D;
it contains a collection of spin clusters of various sizes, which
are formed and decimated during the RG. The entanglement
of a block (shaded area) is give by the number of decimated
clusters (indicated by red loops) that connect the block with
the rest of the system.
local terms in the Hamiltonian and by generating a new
effective Hamiltonian in the frame of the perturbation
theory. If the strongest bond is Jij , the two spins at i and
j are combined into a ferromagnetic cluster with an effec-
tive transverse field h̃(ij) =
. If, on the other hand,
the largest term is the field hi, the spin at i is decimated
and an effective bond is generated between its neighbor-
ing sites, say j and k, with strength J̃jk =
JijJik
. After
decimating all degrees of freedom, we obtain the ground
state of the system, consisting of a collection of indepen-
dent ferromagnetic clusters of various sizes; each cluster
of n spins is frozen in an entangled state of the form:
(| ↑↑ · · · ↑︸ ︷︷ ︸
n times
〉+ | ↓↓ · · · ↓︸ ︷︷ ︸
n times
〉). (3)
In this representation, the entanglement entropy of a
block is given by the number of clusters that connect
sites inside to sites outside the block [Fig. 1]. We note
that correlations between remote sites also contribute to
the entropy due to long-range effective bonds generated
under renormalization.
In 1D the RG calculation can be carried out analyti-
cally and the disorder-averaged entropy Sℓ of a segment
of length ℓ has been obtained as Sℓ =
log2 ℓ [12].
In higher dimensions the RG method can only be imple-
mented numerically. The major complication in this case
is that the model is not self-dual and thus the location of
the critical point is not exactly known. To locate the crit-
-1.5 -1 -0.5
-1.5 -1 -0.5 0
-0.15 -0.1 -0.05 0
L = 16
L = 32
L = 64
-1.5 -1 -0.5
-1.5 -1 -0.5 0
(a) (b)
(d) (e)
=1.175 h
=1.175
=1.175
=1.18 h0=1.17
PSfrag replacements
ln ˜h
ln ˜h
ln ˜J
ln ˜J
ln ˜h
/L0.55
FIG. 2: (color online). The distribution of the last deci-
mated effective log-fields lneh∞, and the distribution of the
last decimated effective log-bonds ln eJ∞ in the RG calcula-
tions. At h0 = 1.175, the distributions, shown in (a) and (b),
get broader with increasing system sizes, indicating the RG
flow towards infinite randomness, i.e. the system is critical.
A scaling plot of the data in (a) using energy-length scaling
lneh∞ ∼ L
ψ with ψ = 0.55 is presented in (c). The solid line
is just a guide to the eye. The subfigures (d) and (e) show
the log-field distribution at h0 = 1.18 and the log-bond dis-
tribution at h0 = 1.17, respectively; the distributions show
a power-law decaying tail in the low energy region, which is
clear evidence that the system is in the Griffiths phases.
ical point, we can make use of the fact that the excitation
energy of the system has the scaling behavior ln ǫ ∼ Lψ
at criticality, while it follows ǫ ∼ L−z in the off-critical
regions. In the numerical implementation of the SDRG
method, the low-energy excitations of a given sample can
be identified with the effective transverse field h̃∞ of the
last decimated spin cluster, or with the effective coupling
J̃∞ of the last decimated cluster-pair.
In our implementation we set for convenience the
transverse fields to be a constant h0 and the random
bond variables were taken from a rectangular distribu-
tion centered at J = 1 with a width ∆ = 0.5. The critical
point was approached by varying the single control pa-
rameter h0. Although this initial disorder appears to be
= 1.170
= 1.175
= 1.180
= 1.185
= 1.190
L = 16
L = 24
L = 32
L = 40
L = 64
= 1.175
PSfrag replacements
ln ℓℓ
FIG. 3: (color online). Left panel: The disorder averaged
block entropy per surface unit Sℓ/ℓ vs. the linear size of the
block ℓ for a system size L = 64 for various values of h0. We
observe that the entropy for ℓ = L/2 reaches its maximum
at the critical point hc = 1.175 (cf. Fig 2). Right panel:
The block entropy per surface area vs. ln ℓ on a log-scale for
different system sizes L at the critical point. The data show
a straight line (guided by the dashed line), corresponding to
the scaling obeying the area law with a double-logarithmic
correction, as given in Eq. (4).
weak, the renormalized field and bond distributions be-
come extremely broad even on a logarithmic scale [Fig. 2]
at the critical point h0 = hc = 1.175. This indicates the
RG flow towards infinite randomness. Slightly away from
the critical point, both in the disordered Griffiths-phase
with h0 = 1.18 and in the ordered Griffiths-phase with
h0 = 1.17, the distributions have a finite width and obey
quantum-Griffiths scaling h∞ ∼ L−z. At the critical
point one has IRFP scaling lnh∞ ∼ Lψ and we estimate
the scaling exponent as ψ = 0.55, quite close to the value
ψ = 0.5 for the 1D case [18].
Now we consider the entanglement entropy near the in-
finite randomness critical point. To obtain the disorder-
averaged entanglement entropy Sℓ of a square block of
size ℓ, we averaged the entropies over blocks in different
positions of the whole system for a given disorder real-
ization and then averaged over a few thousand samples.
In Fig. 3 we show the entropy per surface unit Sℓ/ℓ = sℓ
for different values of h0. This average entropy density
is found to be saturated outside the critical point, which
corresponds to the area law. At the critical point sℓ in-
creases monotonously with ℓ, and the numerical data are
consistent with a log-log dependence:
Sℓ ∼ ℓ log2 log2 ℓ (4)
as illustrated in Fig. 3. In this way we have identified an
alternative route to locate the infinite randomness criti-
cal point: it is given by the field h0 for which the average
block entropy at ℓ = L/2 is maximal. Indeed the nu-
merical results in Fig. 3 predict the same value of hc as
obtained from the scaling of the gaps. We note that the
same quantity, the position of the maxima of the average
entropy, can be used for the random quantum Ising chain
to locate finite-size transition points [21].
The log-log size dependence of the average entropy in
Eq.(4) at criticality is completely new; it differs from the
scaling behavior observed in 2D pure systems, like the
area law, Sℓ ∼ ℓ, for critical bosonic systems [7, 8], or
a logarithmic multiplicative correction to the area law,
Sℓ ∼ ℓ log2 ℓ, as found in free fermions [5, 6, 7, 8]. This
double-logarithmic correction can be understood via a
SDRG argument: In the 1D case a characteristic length
scale r at a given RG step is identified with the aver-
age length of the effective bonds, i.e. the average size
of the effective clusters. At the scale r(< ℓ) the frac-
tion of the total number of spins, nr, that have not been
decimated is given by nr ∼ 1/r [18]; these active (i.e.,
undecimated) spins have a finite probability to form a
cluster across the boundary of the block (a segment ℓ in
the 1D case) and thus to give contributions to the en-
tanglement entropy. Repeating the renormalization until
the scale r ∼ ℓ, the contributions to the entropy are
summed up: Sℓ ∼
dr nr ∼ ln ℓ, leading to the log-
arithmic dependence of the 1D model [12]. For the 2D
case with the same type of RG transformation with a
length scale r < ℓ, the fraction of active spins in the
renormalized surface layer of the block is nr ∼ ℓ/r. Here
we have to consider the situation in which some of these
active surface spins would form clusters within the sur-
face layer and thus contribute zero entanglement entropy;
the number of the active spins that are already engaged
in clusters on the surface at RG scale r is proportional
to ln r, as known from the 1D case, and only O(1) of
the active surface spins would form clusters connecting
the block with the rest of the system. Consequently,
the entropy contribution in 2D can be estimated as:
dr nr/ ln r ∼ ℓ ln ln ℓ, i.e. a double-logarithmic ℓ-
dependence, as reflected by the numerical data in Fig. 3.
Based on the SDRG argument described above, the
double-logarithmic correction to the area law appears to
be applicable for a broad class of critical points in 2D with
infinite randomness. For instance, the critical points of
quantum Ising spin glasses are believed to belong to the
same universality class as ferromagnets since the frustra-
tion becomes irrelevant under RG transformation, and
the same type of cluster formations as observed in our
numerics for the ferromagnet is expected to be generated
during the action of the RG. The entanglement entropy
at the IRFP is completely determined by the cluster ge-
ometries occurring during the SDRG.
Another type of IRFP in higher dimensions occurs
in the bond-diluted quantum Ising ferromagnet: The
Hamiltonian is again given by (1), but now Jij = 0 with
probability p and Jij = J > 0 with probability 1− p. At
percolation threshold p = pc there is a quantum critical
line along small nonzero transverse fields, which is con-
p = 0.49
= 0.50
p = 0.52
1 10 10010
L = 128
L = 256
L = 512
= 0.5
PSfrag replacements
FIG. 4: (color online). The entropy per surface area
Sℓ/ℓ = sℓ vs. ℓ near the percolation threshold pc = 0.5 for
the 2D bond-diluted Ising model at small transverse fields for
L = 512. The curves converge to finite values for ℓ → ∞,
corresponding to the area law. The inset shows sℓ − s∞ as
a function of ℓ. s∞ is estimated from sL/2 at L = 512. The
dashed line corresponds to ℓ−1.
trolled by the classical percolation fixed point, and the
energy scaling across this transition line obeys ln ǫ ∼ Lψ,
implying an IRFP [20]. The ground state of the system
is given by a set of ordered clusters in the same geom-
etry as in the classical percolation model – only nearest
neighboring sites are combined into a cluster. In this
cluster structure, the block entropy, determined by the
number of the clusters connecting the block and the rest
of the system, is bounded by the area of the block, i.e.
Sℓ ∼ ℓd−1 with d being the dimensionality of the sys-
tem. To examine this, we determined the entanglement
entropy by analyzing the cluster geometry of the bond-
diluted transverse Ising model. Fig. 4 shows our results
for the square lattice, which follow a pure area-law with
an additive constant: Sℓ = aℓ+ b+O(1/ℓ).
To summarize, we have found that the entanglement
properties at quantum phase transitions of disordered
systems in dimensions larger than one can behave quite
differently. Generalizing our arguments for the 2D case,
we expect for the random bond transverse Ising systems
a multiplicative d-fold logarithmic correction to the area
law in d dimensions at the critical point, whereas for
diluted Ising model at small transverse fields the area
law will hold in any dimension d > 1 at the percolation
threshold. Although both critical points are described
by infinite randomness fixed points, the structure of the
strongly coupled clusters in both cases is fundamentally
different, reflecting the different degrees of quantum me-
chanical entanglement in the ground state of the two sys-
tems. This behavior appears to be in contrast to one-
dimensional systems governed by IRFPs [12].
Other disordered quantum systems in higher dimen-
sions might also display interesting entanglement prop-
erties: For instance, the numerical SDRG has also been
applied to higher dimensional random Heisenberg anti-
ferromagnets which do not display an IRFP [22]. The
ground states involve both singlet spins and clusters with
larger moments; therefore, we expect the correction to
the area law to be weaker than a multiplicative loga-
rithm and different from the valence bond entanglement
entropy in the Néel Phase [23].
Useful discussions with Cécile Monthus are gratefully
acknowledged. This work has been supported by the Na-
tional Office of Research and Technology under Grant
No. ASEP1111, by a German-Hungarian exchange pro-
gram (DAAD-MÖB), by the Hungarian National Re-
search Fund under grant No OTKA TO48721, K62588,
MO45596.
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http://arxiv.org/abs/quant-ph/0703044
http://arxiv.org/abs/cond-mat/0612503
http://arxiv.org/abs/cond-mat/0703027
|
0704.0419 | Ultrasound attenuation of superfluid $^{3}$He in aerogel | Ultrasound Attenuation of Superfluid 3He in Aerogel
H.C. Choi, N. Masuhara, B.H. Moon, P. Bhupathi, M.W. Meisel, and Y. Lee∗
Microkelvin Laboratory, Department of Physics, University of Florida, Gainesville, FL 32611-8440, USA
N. Mulders
Department of Physics and Astronomy, University of Delaware, Newark, DE 19716, USA
S. Higashitani, M. Miura, and K. Nagai
Faculty of IAS, Hiroshima University, Kagamiyama 1-7-1, Higashi-Hiroshima 739-8521, Japan
(Dated: October 27, 2018)
We have performed longitudinal ultrasound (9.5 MHz) attenuation measurements in the B-phase
of superfluid 3He in 98% porosity aerogel down to the zero temperature limit for a wide range of
pressures at zero magnetic field. The absolute attenuation was determined by direct transmission
of sound pulses. Compared to the bulk fluid, our results revealed a drastically different behavior in
attenuation, which is consistent with theoretical accounts with gapless excitations and a collision
drag effect.
Liquid 3He has attracted intense interest for many
decades in the field of low temperature physics [1]. In
its normal state, liquid 3He has served as a paradigm
for a Fermi liquid whose nature transcends 3He physics.
The superfluid phases of 3He exhibit exotic and intrigu-
ing features associated with the broken symmetries in the
condensate, having an unconventional structure of the or-
der parameter with spin triplet p-wave pairing. Liquid
3He is arguably the most well-understood system mainly
because of its extreme intrinsic pureness at low temper-
atures. Therefore, it has provided important insights
in understanding other unconventional superconductors
such as the high temperature superconductors, the heavy
fermion superconductors, and in particular the more re-
cently discovered Sr2RuO4, which is also thought to have
the p-wave symmetry [2]. However, the same virtue has
hampered the effort in pursuing answers to an important
overarching question: how does the nature of a quantum
condensate (spin triplet p-wave superfluid in this case)
respond to increasing impurity or disorder?
Observation of superfluid transitions in liquid 3He im-
pregnated in high porosity aerogel in 1995 [3, 4] opened
a novel path to introducing static disorder in liquid 3He.
Aerogel possesses a unique structure, whose topology is
at the antipode of widely studied porous media such as
Vycor glass and metallic sinters. Due to its open struc-
ture, there are no well-defined pores in aerogel and conse-
quently, the liquid is in the proximity to the bulk. Ninety
eight percent porosity aerogel, which has been used in
most of the studies including this work, offers a corre-
lated network of strand-like aggregates of SiO2 molecules
whose structure can be characterized by the geometrical
mean free path (ℓ ≃ 100 - 200 nm), the diameter of strand
(r ≈ 3 nm), and the average inter-strand distance (d ≃ 25
- 40 nm). The coherence length of pure superfluid 3He,
ξ0, which varies from 20 nm (34 bar) to 80 nm (0 bar),
is at least an order of magnitude larger than the strand
diameter but is comparable to ℓ and d. As a result, the
scattering off the aerogel strand would have a significant
influence on the superfluid. It is now well established
that the superfluid transition temperature is significantly
depressed from that of the bulk, and the effect of pair-
breaking is progressively magnified at lower pressures,
leading to the possibility of a quantum phase transition
at Pc ≈ 6 bars [5]. To date, three distinct superfluid
phases have been experimentally identified, namely the
A-like, B-like, and A1-like phases [4, 6, 7, 8, 9]. The B-like
phase and the A1-like phase in aerogel show striking sim-
ilarity to their counterparts in the bulk superfluid [9, 10].
Detailed NMR studies [7, 8, 10] suggest that the aerogel
B-phase has the same order parameter structure as the
bulk B-phase. The aerogel A1-phase only appears in the
presence of magnetic field as is the case in the bulk [9].
However, the aerogel A-phase exhibits quite a different
behavior from the bulk A-phase (e.g. in NMR frequency
shift and superfluid density), although the overwhelming
experimental evidence suggests that it is an equal spin
pairing state. Various interpretations or novel proposi-
tions on the possible order parameter structure have been
suggested for this phase [11, 12, 13].
Nuclear magnetic resonance and ultrasound spec-
troscopy have been used in concert to investigate the mi-
croscopic structure of the superfluid phases [1, 14]. These
two experimental methods encompass complementary in-
formation on the orbital (ultrasound) and spin (NMR)
structure of the Cooper pairs. Rich spectra of order
parameter collective modes in bulk superfluids, which
are the fingerprints of specific broken symmetries in the
system, have been mapped by ultrasound spectroscopic
techniques [14]. In 2000, Nomura et al. [15] performed ul-
trasound attenuation measurements on 98% aerogel us-
ing a 16.5 MHz cw acoustic impedance technique. Their
work was limited to a single pressure at 16 bars and down
to 0.6 mK. Although their technique was not adequate in
determining absolute attenuation, they managed to ex-
tract the absolute sound attenuation after making auxil-
http://arxiv.org/abs/0704.0419v1
Time (µs)
FIG. 1: Acoustic response from the receiver vs. time at 34
bars for select temperatures ranging from 0.3 mK to 2.5 mK.
The aerogel superfluid transition is marked by a small arrow.
iary assumptions. A Bayreuth group [16] performed ab-
solute sound attenuation measurements in aerogel (97%
porosity) using a direct sound transmission technique at
10 MHz. They experienced poor transducer response,
and observed self-heating and no depression in the aero-
gel superfluid transition. We conducted high frequency
sound transmission experiments in 98% porosity aero-
gel, covering the whole phase diagram of the superfluid
phases in aerogel, from 8 to 34 bars and from the transi-
tion temperatures to as low as 200 µK.
In this experiment, two matched LiNbO3 longitudi-
nal sound transducers with the fundamental resonance at
9.5 MHz were used as a transmitter and a receiver. The
6.3 mm diameter transducers were separated by a Ma-
cor spacer maintaining a 3.05 (± 0.02) mm sound path
between the transducers where the aerogel sample was
grown in situ. This scheme ensures the best contact be-
tween the transducer surface and the aerogel, which is
crucial for clean sound transmission at the boundaries.
A 1 µs pulse was generated by the transmitter and de-
tected by the receiver. Temperature was determined by
a melting pressure thermometer (MPT) for T ≥ 1 mK
and a Pt NMR thermometer for T ≤ 1 mK which was
calibrated against the MPT. No non-linear response or
self-heating was observed at the excitation level used in
this work. All the data presented here, except for 8 bars,
were taken while warming with a typical warming rate
of 3 µK/min. A detailed description on the experimen-
tal cell and experimental techniques can be found else-
where [17, 18].
The temporal responses of the receiver taken at 34 bars
are shown in Fig. 1 for select temperatures ranging from
0.3 to 2.5 mK. The primary response, which starts to rise
around 8 µs, shows a rather broad response due to ring-
ing of the high Q transducer (Q ∼ 103). The step-like
structure of the receiver signal is caused by the slight mis-
match in the spectra of the transducers [18]. Below the
aerogel superfluid transition (marked around 2.1 mK by
an arrow in Fig. 1) the primary response starts to grow
and the trailing echoes emerge from the background, as
the sound attenuation decreases in the superfluid. No
change in the receiver signal was observed at the bulk
superfluid transition. The multiple echoes follow a bona
fide exponential decay in time. Absolute sound attenua-
tion was obtained in the following manner [19]. First, the
relative attenuation at each temperature was calculated
using the area under the primary response curve by inte-
grating the signal from the rising edge to a fixed point in
time (23 µs point). The absolute attenuation at 0.4 mK
and 29 bars, obtained using the primary signal and the
echoes, was used as a reference point in converting the
relative attenuation into the absolute attenuation. Due
to a drastic mismatch in the acoustic impedance at the
the transducer-aerogel/3He boundary, the signal absorb-
tion at the surface of transducers was ignored [19]. The
possible background contributions to attenuation from
the quasi-particle scattering off the cavity wall [20] and
the non-parallel alignment of the two transducers are es-
timated to be negligible.
The absolute attenuations on warming for several pres-
sures are plotted as a function of temperature in Fig. 2(a).
The superfluid transition is marked by the smooth drop
in attenuation. Our aerogel superfluid transition tem-
peratures are in excellent agreement with the previously
reported values for all pressures [5, 21]. At 9.5 MHz in
the bulk B-phase, a strong attenuation peak appears right
below the superfluid transition. This peak is the result of
the combined contributions from pair-breaking and cou-
pling to the order parameter collective modes. Above the
polycritical pressure, the B to A transition on warming
is registered as a sharp step in attenuation. In aero-
gel, none of these features exist. However, we did ob-
serve a sharp step in attenuation on cooling for P > 14
bars, which implies the existence of the supercooled A-
phase [19]. We were able to identify a rather smooth B
to A transition on warming for 29 and 34 bars within ≈
150 µK below the superfluid transition. This observation
is consistent with the previous results obtained using a
transverse acoustic impedance technique [13]. Therefore,
most of the attenuation data presented here are in the
aerogel B-phase. In the bulk B-phase with a clean gap,
the attenuation follows α ∝ e−∆(T )/kBT below the atten-
uation peak, practically reaching zero attenuation below
T/Tc ≈ 0.6, due to thermally activated quasi-particles,
where ∆(T ) is the temperature dependent gap and kB
is the Boltzmann constant. In contrast, the attenuation
in aerogel decreases rather slowly with temperature and
remains high even at T/Tc ≈ 0.2. Furthermore, a pe-
culiar shoulder feature appears at T/Tc ≈ 0.6 for higher
pressures. This feature weakens gradually and eventually
disappears at lower pressures, Fig. 2(a).
Sound propagation for higher harmonics up to 96 MHz
was measured for several temperatures and pressures,
but no evidence of sound propagation was found above
30 MHz even at 0.3 mK, where the lowest attenuation
is expected. Below about 10 mK, the scattering process
is dominated by the temperature independent impurity
scattering off the aerogel, and at 9.5 MHz, ωτi ∼ 0.1 for
all pressures where τi = ℓ/vf (see below for ℓ). There-
fore, the sound mode should remain in the hydrodynamic
limit. This claim is bolstered by the observation of the
strong frequency dependence in attenuation and the ab-
sence of a temperature dependence in the normal fluid
attenuation [15]. The coupling between the normal com-
ponent of the superfluid 3He and the mass of the elas-
tic aerogel modifies the conventional two-fluid hydrody-
namic equations [22, 23]. This consideration leads to two
(slow and fast) longitudinal sound modes with different
sound speeds, cs = ca
ρsρa/ρ, and cf = c1
1+ρaρs/ρnρ
1+ρa/ρn
Here, cf(s) represents the speed of the fast (slow) mode,
ρn(s) is the normal fluid (superfluid) density (ρ = ρn+ρs),
ρa is the aerogel density, c1 is the speed of hydrodynamic
sound in 3He, and finally ca is the sound speed of the bare
aerogel. From the time of flight measurements, we found
the sound speed in aerogel consistently lower (by ≈ 20%)
than c1 for all pressures studied and in good agreement
with the values obtained using the expression above [24].
Detailed analysis of sound velocity for various pressures
will be presented in a separate publication.
Low mass density and the compliant nature of aero-
gel necessitate the consideration of effective momentum
transfer upon quasi-particle scattering off the aerogel,
which generates dragged motion of aerogel. Ichikawa et
al. [25] incorporated the collision drag effect in calculat-
ing the dispersion relation in the normal fluid. Their
model offered a successful explanation for the experimen-
tal results of the Northwestern group [15]. Recently, Hi-
gashitani et al. [26, 27] extended this model to study the
longitudinal sound (fast mode) propagation in superfluid
3He/aerogel within the framework of the two-fluid model.
The drag effect can be described phenomenologically by
a frictional force, ~Fd =
(~vn −~va), introducing an addi-
tional relaxation time τf , where ~vn(a) is the normal fluid
component (aerogel) velocity. This effect is of particu-
lar importance when ωτi < 1, and the total attenuation
(Eq. (130) of ref. [27]) is
ω2/2cf
1 + ρaρs/ρnρ
ρ2aτf/ρρn
1 + ρa/ρn
4η/3ρc21
1 + ρaρs/ρnρ
), (1)
where η is the shear viscosity of liquid 3He. The first
term (αf ) arises from the frictional damping caused by
the aerogel motion relative to the normal fluid compo-
nent, and the second term (αv) from the conventional
hydrodynamic sound damping associated with the viscos-
ity. This expression allows us to extract ℓ in this system
from our absolute attenuation at the transition temper-
ature, αc. The inset of Fig. 3 shows our results of αc for
0.0 0.2 0.4 0.6 0.8 1.0
0.0 0.5 1.0 1.5 2.0
10 bar
14 bar
20 bar
25 bar
33 bar
34 bar
T / T
10 bar
12 bar
14 bar
21 bar
25 bar
34 bar
T (mK)
FIG. 2: (a) Absolute attenuation for various pressures vs.
temperature (color in on-line version). Thin solid lines are
the results of a quadratic fit to the low temperature part
(T/Tc <∼ 0.4) of the data at each pressure. (b) Normalized
sound attenuation vs. normalized temperature. The results
of theoretical calculation (solid lines, color in on-line version)
are plotted along with the experimental results at 34 bars for
comparison.
various pressures. The solid lines are the result of calcu-
lation using Eq. (1) for three different mean free paths,
ℓ = 100, 120, and 140 nm. As can be seen, ℓ = 120 nm
produces an excellent fit to our data for the whole pres-
sure range, which is in good agreement with the val-
ues obtained from the thermal conductivity (90 nm) [28]
and spin diffusion (130 nm) [29] measurements. With
the knowledge of the mean free path, one can calculate
the full temperature dependence of sound attenuation in
the superfluid phase. The results of the calculation (in
the unitary limit) following the prescription described in
ref. [27] are displayed in Fig. 2(b) along with the experi-
mental results at 34 bars. The calculation reproduces all
the important features observed in our measurements. In
particular, the conspicuous shoulder structure appearing
near T/Tc ≈ 0.6 at 33 bars softens at lower pressures and
is completely absorbed in an almost linear temperature
dependence below 20 bars. This behavior is the charac-
teristic of αf [27]. A fast decrease in ρn right below Tc
produces the bump in αf , and αf → 0 as T → 0. On
the other hand, αv decreases monotonically and reaches
a finite value due to non-zero ρn and the impurity states
0 5 10 15 20 25 30 35
0 10 20 30
P (bar)
100 nm
120 nm
140 nm
P (bar)
FIG. 3: Normalized zero temperature attenuation vs. pres-
sure. The dashed line is a guide for eye. Inset: Pressure
dependence of sound attenuation at Tc. The solid lines (color
on-line) are the results of theoretical fit for ℓ = 100, 120, and
140 nm (see text).
induced inside the gap as T → 0. The quantitative agree-
ment between the theory and experiment, however, is not
yet satisfactory. The calculation utilizes the isotropic ho-
mogeneous scattering model (IHSM) [30], which tends to
overestimate ∆(T ) and ρs compared to the experimen-
tally determined values [3, 23]. As shown in ref. [31],
the inhomogeneity gives rise to the reduction of the aver-
age value of the order parameter and consequently yields
larger η and ρn, which in turn increases α0 but decreases
the frictional contribution. It is also expected that the
non s-wave scattering components make non-trivial con-
tributions to the viscous and frictional relaxation times
in a direction that improves the quantitative agreement.
Theoretical calculations based on the IHSM [27, 32]
predict that the impurity states would completely fill the
gap, leading to a gapless superfluid when τiTc < 1 for the
B-phase in the unitary limit. We estimate 0.3 < τiTc < 1
for 10 < P < 34 bars with ℓ = 120 nm. The normalized
zero temperature attenuation (α0/αc) obtained by ex-
trapolating the low temperature part of the attenuation
(solid lines in Fig. 2(a)) is plotted in Fig. 3, where α0/αc
increases as the sample pressure is reduced and seems
to approach unity near Pc ≈ 6 bars. Since the viscosity
ratio is directly related to the density of states at zero
energy through η(0)/η(Tc) = n(0)
z, z = {2,4} for the
{Born, unitary} limit where n(0) is the normalized den-
sity of states at zero energy [27], the finite α0/αc is strong
evidence of a finite n(0). The gapless behavior has been
experimentally suggested by recent thermal conductivity
(for P ≤ 10 bars) [28] and heat capacity (for 11 ≤ P ≤ 29
bars) [33] measurements. The pressure dependence of
α0/αc is in qualitative agreement with the combined re-
sults of Fisher et al. and Choi et al. Although all of
these experimental techniques (including ours) are lim-
ited to probe the impurity states near the Fermi level,
the behavior is consistent with the theoretical predictions
with gapless excitations. Unlike the thermodynamic and
transport measurements, the high frequency ultrasound
measurement has a potential to unveil a larger portion
of the impurity states profile from the frequency depen-
dence.
We acknowledge support from an Alfred P. Sloan Re-
search Fellowship (YL), NSF grants DMR-0239483 (YL),
DMR-0305371 (MWM), and a Grant-in-Aid for Scientific
Research on Priority Areas (No. 17071009) from MEXT
of Japan (SH and KN). We would like to thank J.-H.
Park for his technical assistance, and Jim Sauls, Peter
Wölfle, and Bill Halperin for useful discussions.
∗ [email protected]
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mailto:[email protected]
|
0704.0420 | The Hourglass - Consequences of Pure Hamiltonian Evolution of a
Radiating System | The Hourglass—Consequences of Pure
Hamiltonian Evolution of a Radiating System
Donald McCartor
ABSTRACT
Hourglass is the name given here to a formal isolated quantum system that can
radiate. Starting from a time when it defines the system it represents clearly
and no radiation is present, it is given straightforward Hamiltonian evolution.
The question of what significance hourglasses have is raised, and this question
is proposed to be more consequential than the measurement problem.
1 Hourglasses
2 Physics without true histories
3 But histories are sometimes good
4 Phlogiston and oxygen
5 A closer look at quantum engineering
6 Conclusion
But I want to know the particular go of it
– the plea of James Clerk Maxwell as a young child
concerning, among many things, the bell-wires that
ring the bells that summon servants. [Mahon]
1 Hourglasses
Suppose that theory develops in such a way that quantum fields can be handled
like nonrelativistic quantum mechanics. Then if we are interested in something,
perhaps gooseberry bushes, we can model one as we would conceive it to be at
some instant and then follow its development through time. And not only the
atoms and molecules would be modeled, but also the radiation.
This is a scheme for the imagination. The gooseberry bush, though not
isolated, would grow within a suitable environment that would be an isolated
system, complete in itself. We do learn well from isolated systems, both real
ones in the laboratory and those envisaged in our theoretical musings.
We will provide the bush with air, earth, and water. And there can be
life-giving sunlight shining on it. As for the light that had been reflected or
emitted from the bush before the present time, we will leave that out. Such
light goes off and away, so it could only matter as information about what the
http://arxiv.org/abs/0704.0420v1
bush had been doing. We will take the bush just as it is now.
The gooseberry bush is then developed forward in time. Lagrange or
Hamilton would have recognized what we are doing, for we are doing physics
the classical way. We have an initial condition and we are finding out what
will happen next.
As we move toward the future, light shoots out from the bush, as we
expect. But, disconcertingly, the bush starts to lose definition. Its parts lose
their precise places. Within a few weeks it is a scarcely recognizable mess.
Let’s go back in time, then. This is terrifyingly worse. The bush has been the
subject of a vast conspiracy. Light has been streaming in on it from the entire
universe. The bush swallows it up. Then at the present time this suddenly all
stops. Time symmetry of the Hamiltonian makes it happen like that.
This is the hourglass. It is really more like a cone, with the light streaming
in before the set-up time forming one nappe and the light streaming out after
it the other. But hourglass is a more colorful name.
What to do? We will try to bring quantum mechanics to the rescue.
We will make what is conventionally called a measurement, but cautiously. A
place is chosen well outside the gooseberry bush, and a time chosen that is
later than when we set the state of the bush up. A check is made of whether
there are at this place and time any photons coming from the direction of
the bush. By doing things this way, we won’t disturb the bush at all, and
we don’t care if we disturb the escaping light. We get from this, of course, a
probability distribution over various possibilities for photons at this place and
time. Encouraged by this small success, we choose another place and time and
do the same. And this is what is nice: the two measurements are compatible.
Thus we get correlations between them, too. Emboldened by this opportunity,
we do millions of them, which all formally combine into a single measurement
with a single set of possible results. Each possible result of the single, combined
measurement is a combination of results of all the individual measurements
of light made at the various times and places. Thus each combined result
constitutes a kind of movie of the gooseberry bush.
What will the most probable of these results be like? This is the problem
of the hourglass. To begin with, however, it may be that there is no hourglass.
The deepest quantum theory might not provide a system with a state and
its evolution. Or if it does, it could still be objected that the Hamiltonian
evolution should not have been allowed to run on unchecked. There should
have been many quantum jumps. By leaving them out, quantum mechanics
has been misused, and what results is no matter.
But Lagrange and Hamilton and would have been best pleased if these
objections did not hold. And surely we would then hope to see in each of the
most probable results something like a movie of a bush producing gooseberries:
physics working right. The bushes in these movies would look much alike at
the start but then gradually differ, as chance has it. We would learn something
about how gooseberry bushes grow gooseberries!
Certainly Lagrange and Hamilton would have thought the problem of the
hourglass a leading one, if they had known of quantum mechanics. Indeed,
every physicist might like to take a stab at guessing its solution, just to orient
themselves in their science. Does the hourglass fail, and if so, where and why?
Or if it does produce movies true to our world, but not from a developing
quantum state that might be the true history of a gooseberry bush, rather
from a “history” that does at one time represent a gooseberry bush well, but
soon is unlike anything that ever did exist, then how can this be?
2 Physics without true histories
Here is what I think about it. But before we go into that, see if you don’t agree
that the hourglass question has gravity, and this regardless of the ideas that I
or anyone might have for its answer.
Now my guess is that quantum mechanics will give us movies of ripening
gooseberries, produced by hourglasses through the means described or some-
thing rather like that. And I think that to understand hourglasses, not to
solve the measurement problem, is the central question for the understanding
of quantum mechanics.
For the measurement problem begs a question, which makes it futile. It
assumes that we learn from physics simply because physics describes well those
things that exist. Like this example from classical physics. There exist in a
gas a multitude of zipping molecules. At any given moment, each particle has
its particular position and momentum, and over time this forms their history.
Physics has told us what a gas is—precisely what exists there. This is what
lets us learn about gases. Undoubtedly this is how Boltzmann saw it.
But when we look at the statistical mechanics he produced, and even
more at that of Gibbs, a person will acquire deep qualms about this view-
point. Boltzmann’s analysis of the collision of molecules seems like straightfor-
ward common sense. He is looking at what they are likely to do. But when
Loschmidt’s reversibility objection is brought forward, the lucidity vanishes.
Gibbs’s more abstract statistical mechanics made the problem even starker.
Gibbs found beautiful mathematical form in Boltzmann’s (and Maxwell’s)
work, which he generalized. He held that thermodynamic systems should be
represented as being in states that have the form of certain probability dis-
tributions over classical states. Gibbs could not well understand what these
probabilities were about, but he saw that his theory was good nevertheless. To
keep this lack of clear comprehension from poisoning work with the theory, he
devised a work-around. The axioms of probability theory are reflected in the
axioms of finite set theory. One can effectively solve problems of probability
by thinking about finite sets. So Gibbs suggested that we simply think about
these probabilities in terms of sets. The word he used was ensembles.
Gibbs described his intent in these words: “The application of this prin-
ciple is not limited to cases in which there is a formal and explicit reference to
an ensemble of systems. Yet the conception of such an ensemble may serve to
give precision to notions of probability. It is in fact customary in the discus-
sion of probabilities to describe anything which is imperfectly known as some-
thing taken at random from a great number of things which are completely
described.”[Gibbs]
But physicists have never been able to accept gracefully that they don’t
understand the elements of their science. So they have been moved to think
that they do understand Gibbs’s probabilities somehow, and this has led to
two missteps.
One has been to regard the probabilities in Gibbs’s theory as being the
result of our ignorance of the detailed state of the system we are considering.
But when a probability distribution is useful, this is a very great step up in
order from chaos. Ignorance cannot create order. If water always boils at the
same temperature, it is not our fault. Rather than being so explained, for it
is not, Gibbs’s theory shows that there is something deeply wrong with classi-
cal mechanics. Classical statistical mechanics is not really a form of classical
mechanics. It is quantum mechanics being born.
The following words of Gibbs seem to show that Gibbs himself took the
view just scotched. “The states of the bodies which we handle are certainly not
known to us exactly. What we know about a body can generally be described
most accurately and most simply by saying that it is one taken at random from
a great number (ensemble) of bodies which are completely described.”[Gibbs]
The impression that I get, though, is that Gibbs is cautiously hedging. He is
not saying plainly, as he might have, that a body we handle will be in some
completely described state, so that if we describe it with an ensemble, the
probabilities in the ensemble simply represent our partial ignorance about that
state. He does say plainly that his method seems to work.
The other misstep has come about because quantum theory is a mirror
of Gibbs’s statistical mechanics in the sense that it is based on what are prob-
abilities in form (in other words, sets of non-negative real numbers that add
up to one) and we don’t know what they mean in general. It is true that we
can make good sense of them as real probabilities in various special cases. For
instance, when quantum mechanics is applied to the Stern-Gerlach experiment,
to see the detector react is like seeing a coin tossed. But in the general case no
such kind of experience is directly implied by these probability forms. There
are, for example, canonical distributions in quantum mechanics too, and we
don’t ever expect to see a detector pick a pure state out of a hot cup of coffee.
We then sometimes think about these formal probabilities in terms of
ensembles, just as Gibbs did, and for the same reason. Where the formal
probabilities are highest and the members of the ensemble most numerous,
there the greatest significance will lie, whatever it may be. This is fine. But
quite often physicists say that ensembles (that is to say, Gibbs’s work-around)
provide the means to understand quantum theory. This is clearly wrong.
But to get back to the measurement problem. As you well know, but
for explicitness I will say it anyway, to see a problem in measurement is to
suppose that quantum mechanics can describe the equipment in the lab as it
exists at the start of an experiment, but when the representation is continued,
the equipment becomes entangled with the microscopic systems it is examining
and gets smeared. Then quantum mechanics has stopped describing what we
know exists in the lab and needs to be corrected so that it will continue to
describe what exists.
But it isn’t so that quantum mechanics, if it is to show us some pre-
dictability in nature, must provide us directly with histories of the existence of
things, as by a developing wave packet. As evidence, I offer the hourglass.
3 But histories are sometimes good
If physics does not work simply because it describes what exists, and if, rather,
the way of the hourglass is right, then a corollary is that how we learn about
nature necessarily becomes more indirect. We are given such information as
radiation provides about something, not directly told what exists there. And
for the purpose of inferring useful rules of nature’s behavior, what we deal with
are imagined situations that we think typical of what we want to learn about,
not faithful descriptions of actual things. No real radiating system is like an
hourglass, except momentarily near the hourglass’s neck.
But quantum engineering may temper the truth of that judgment just a
bit. For there is also an engineering use of quantum mechanics where, somewhat
as classical mechanics does it, for a time we can use a wave packet to represent
the development of an actual situation we are dealing with. But this is rather
more special, for we must take care to set things up so that this will work. The
vacuum must be excellent, etc. Isolation is important.
A simple example of quantum engineering is an ion that alternately blinks
for a spell and remains dark for a spell while sitting in an ion trap that is irradi-
ated by lasers. You can picture the ion well enough by thinking of Schrödinger
evolution of a wave packet with occasional quantum jumps interspersed. You
might then be tempted to think that everything can be handled effectively in
the same way, at least in principle. We have just not been clever enough to
find New York City’s wave packet and its measurement collapses.
This is trouble. The worst of it is that you will be led to ignore hourglasses
and what they imply, since clearly hourglasses cannot represent the history
of things in the same manner that you have advantageously represented the
history of the blinking ion.
On the other hand, imagine that decades ago physicists had taken hour-
glasses to their hearts, as well I think they might have. Then they could have
been tempted to look upon representing an ion in a trap by Schrödinger evolu-
tion of a wave packet with quantum jumps as ‘following the wrong philosophy’
(by trying to represent the actual histories of things with wave packets), and
might have disdained to do so. There is a lesson here. Don’t take your philo-
sophical ideas too seriously, we’re not good enough for that.
I believe, though, that from hourglasses you would be able to infer that
Schrödinger evolution with jumps is a simple and effective (not perfect) way
to regard a blinking ion in a trap. The hourglasses would then be in this sense
the more fundamental theory.
4 Phlogiston and oxygen
But what is a quantum jump? Here is where I think the community of physi-
cists has been careless in the use of words, perhaps mixed with real misun-
derstanding. Two principles of quantum physics have been formulated. The
first principle (promoted by Dirac and von Neumann) is that when a measure-
ment is made on a system, an immediately following measurement will give the
same result. Therefore, right after any measurement the system must be in the
eigenstate corresponding to the value found.
The second principle is that if the probabilities of the possible results of
all the measurements that may be made on a system are defined, then there
will be a (unique) quantum state that the system may be said to be in that will
yield these probabilities. Add to this that sometimes two measurements may
be made on a system without interfering with each other. Then when one of the
two measurements has a certain result this will define a conditional probability
for any result of the other measurement (simply divide the probability that
both results occur by the probability that this result of the first measurement
occurs). According to the second principle, then, there will be a quantum state
that yields these probabilities (for the possible results of any measurement that
may be made without interfering with, or suffering interference from, a given
measurement that has had a certain result).
Please notice that the argument above assumes that the set of all the
measurements compatible with a given measurement effectively constitutes ‘all
the measurements that may be made on a system’ as needed by the second
principle.
Now consider a system A in the state α. It is composed of two subsystems,
B and C, in reduced states β and γ respectively. A measurement is made on
subsystem B and it has a result. By the first principle, there is a quantum
state β′ that will yield the probabilities of the possible results of any immedi-
ately following measurement that might be made on subsystem B. And by the
second principle, there is a quantum state γ′ that will yield the probabilities of
the possible results of any compatible measurement made on subsystem C.
For the supplanting in one’s considerations of β by β′ there is the historical
name ‘collapse of the wave packet’. For the supplanting in one’s considerations
of α by γ′ most physicists use the same phrase (or any of its several synonyms).
It would easier to think about these things if different names were used for the
two. ‘Collapse of the wave packet’ might be retained for the first and, say,
‘conditioning of the wave packet’ adopted for the second.
This is all the more important because the first principle is an out and
out mistake by Dirac and von Neumann, whereas the second is an inalienable
part of quantum mechanics. To those two mathematically minded, and so
logically minded, people, the dignity of quantum mechanics required that there
be measurements, so that quantum mechanics might be real physics. And
since quantum mechanics did not say that a system had to have, before the
measurement, the value found in the measurement, the dignity of measurement
required that it at least have that value afterward, or what sort of measurement
was this anyway?
Tacked on to this was the fact that so distressed Schrödinger: wave pack-
ets spread interminably. If a developing wave packet were to represent the
history of a system, which they assumed to be necessary, then the spreading
had to be checked, and an occasional quantum jump such as their measurement
theory presupposed might do that.
And experiment lent some support. Above all, if an electron went splat
somewhere on a screen, which they regarded as a measurement by the experi-
menter of the electron’s position, then conservation of charge suggested strongly
that the electron could be found subsequently thereabout. This was the origin
of the phrase ‘collapse of the wave packet’. Too, the famous Stern-Gerlach
experiment allows a following measurement of spin, which will give the same
result as the first if the first measurement’s detection has been delicate enough.
But the idea of a quickly following measurement is just not well-defined
in general. And there are cases where the principle must prove false under any
reasonable definition of a following measurement. For example, a particle might
lose most of its energy in those collisions that measured its energy. Or if the
momentum of a charged particle were measured by the curvature of its path in a
magnetic field, the particle might end up going in the wrong direction, although
this is, to be sure, correctable. Those events called “measurements” are what
they are, and if they fall short of truly being measurements of properties, so
be it!
If the first principle is an error, then that leaves us with only one principle,
the second, and people might then be inclined to continue to use the traditional
phrase ‘collapse of the wave packet’, but now meaning the replacements the
second principle defines. This would result in the transfer, in the course of
history, of the meaning of the phrase from the first principle to the second. I
think that this would have the same unhappy effect as if Lavoisier, not wishing
to burden the world with a neologism, had instead given to the word phlogiston
a new sense.
5 A closer look at quantum engineering
The second principle has a very different flavor from the first. For it leads to
conditional probabilities, and these lend themselves to imaginative thinking. In
this mind-set you are free to take up points of view according to what you wish
to learn. The first principle, however, leads to probabilities that are thought
to be the properties of real events, such as an actual toss of a coin. You are
now in a reality mind-set. That probability is as much a part of the coin toss
as is the silver of the coin, and you must deal with it. You have no choice.
But I don’t mean to say that this is an absolute difference between the two
principles. Rather, they tend to lead us into these respective modes of thought,
and vice versa. Bearing this in mind, let us look at the hourglass and quantum
engineering.
First consider the hourglass that represents a gooseberry bush. By choos-
ing one among the more probable of the results of the course of observation
of light, we will select what is in effect a likely movie of such a bush. We can
look at the movie, and the marvelous algorithms of our brains will construct
an idea of a gooseberry bush and follow it through its history. We have gotten
something good out of this, and we have made no use of the conditional proba-
bilities offered by the second principle at all. However, if we are not limited to
one movie then we can use conditional probabilities as they are normally used,
to explore various interesting possibilities while taking into account how likely
they are when we are supplied with certain information.
Notice that we have been thinking imaginatively. No one would suppose
that we have directly grasped the reality of a gooseberry bush in our garden in
this way, particularly because real gooseberry bushes do not start to exist at a
special time.
Now consider quantum engineering. By means of careful construction
of the equipment a clearly defined situation can be set up where the power
of wave packets to give understanding will be enhanced. On the other hand,
here there can be significant entanglement. The power of our minds to achieve
understanding through their everyday methods will be set at nought.
Then for quantum engineering, a history formed by wave packet develop-
ment with occasional saltations may be a quite good route to understanding.
We would take up this idea of what exists simply because it is good enough to
help us with the job at hand. And for this case, where we find it fit to think
that we are dealing with an actual system that is an evolving wave packet,
and with saltations that we regard as actual events, but in a way so different
from that intended by Dirac and von Neumann, then perhaps a third term,
say, ‘change of the wave packet’, would be appropriate for the saltations.
These changes of the wave packet would differ from collapses of the wave
packet because, although they would be thought of as real events just as col-
lapses have been, they would be derived from conditioning of wave packets, in
the following manner. When a system sends out radiation (or anything else)
that will not return, in one way you can consider the system of interest to be
the whole, including the radiation, and in another way you can consider it to
be the reduced system that does not include the radiation. Upon observation
of the radiation you will derive from the result and from the wave packet of
the whole system a wave packet for the system less the radiation, and this we
have called conditioning of the wave packet. But if before the observation your
interest had been focussed on the system less the radiation, and thus on its
reduced wave packet, then you will have gone from one wave packet to another
wave packet for the system less the radiation. And since you are reckoning
these wave packets as being portions of the system’s history, this looks like a
quantum jump. This is what is meant by a change of the wave packet. There
is no need to define any such change of the wave packet precisely, of course.
No more is there need to suppose that it can be defined with precision.
6 Conclusion
If hourglasses cannot be true histories, how can it be that we can learn from
them? What lets them tell us how gooseberry bushes grow, when they are only
momentarily like a gooseberry bush? I haven’t said a word about this yet.
First of all, there is an assumption hidden behind this puzzlement of ours.
The assumption is that we have no reason to be perplexed that we can learn
from things that can be true histories. For if it did not seem so perfectly natural
to us that we learn from true histories, then it would not appear unnatural to
learn from what clearly cannot be a true history. But I think this assumption
of ours is thoughtless, and I will try to explain why.
We make judgments about when we are better informed and when less so.
The ideas we hold true when thought to be better informed are compared with
those that we held when not so well informed. In this way, through the device
of taking the ideas that we presently have most confidence in as trustworthy,
we try to gather how successfully our ideas tend to stack up against reality. It
is not quite so simple, however, since we know from sad experience that the
ideas we now trust may fail us. But we have the conviction, or hope, that if
such happens we can land on our feet again. We will search for still better
ideas until we find something that works.
We are apt to give to this situation a logical cast. Namely, by positing
that there is a best of all possible ideas in whose direction we are headed. This
posit can be helpful. It can give us greater confidence in our search for better
ideas. If we guess that this best idea will have a certain form, and we guess
well, it can guide our search. But there is no necessity for this posit; all we
really know is what was said above.
Another thing we like to do is to find where things are and when. Our
vision, touch, and hearing do this automatically all the time, and we often give
them some conscious help, say by turning the head. When we are a teenager
it is likely to occur to us that there must be a best of all possible such ideas, a
complete map of where everything is, and has been, and perhaps will be too.
A further thought may cross one’s mind. Maybe this is all that our world is.
For instance, if one person likes another, this should show up in that person’s
actions, which the map will completely define. Maybe the liking simply is those
actions.
Now I will propose some physics, the red dust theory. According to this
theory the world is made up of an exceedingly large number of very fine specks
of a scarlet dust. Because of its ruddiness, the dust is extremely beautiful, if
only we could see it, but we will not be concerned with that. The red dust
theory differs from most physics in that the flight of the particles does not have
to satisfy a differential equation, it is merely continuous.
The interpretation of the theory is quite simple. Where we find things
there will be a crowd of these specks, and where we find vacancy they will be
much sparser. But can our world be as this theory says? Surely it can. There
will be among its solutions one that maps the entire history of our universe with
extraordinary precision. The collisions of galaxies, the evolution of whales, the
experiments in laboratories, all will be there and rightly shown.
Now you may think that the red dust theory is hopelessly bad physics
and should be ignored. It may be hopelessly bad, but it should not be ignored.
It is a benchmark. If another physics theory is proposed, is it better than the
red dust theory, and if so, just why? This is especially pertinent if the other
theory intends, as does the red dust theory, to give a precise description of all
that exists. Bohmian quantum mechanics is an example.
But what I intend to put up against the benchmark is classical mechanics.
Everyone will agree that classical mechanics is far better than the red dust
theory. You can do things with classical mechanics; you can’t do anything
with the red dust theory. For instance, you can pull a pendulum to the side
and let it go. It will swing. Classical mechanics can give you the history of that
swing ahead of time. The red dust theory has so many solutions compatible
with the way things are at the start that it won’t tell you anything useful about
how things will go.
Our experience with classical mechanics is that it is practical, but why
is this so? The most natural idea is that the world must at bottom be clas-
sical mechanical. Since we understood the pendulum by assigning a classical
mechanical state to it and evolving the state, there must then be an evolving
classical mechanical state that the whole world is in, and that would explain
why classical mechanics is so useful.
When we look at the history of our universe, however, and particularly at
the evolution of life over billions of years, and when we consider the resources
that it is likely that classical mechanics has to offer in its solutions, it doesn’t
really seem possible that there is any classical mechanical history that would
match our universe’s history, no matter how exquisitely the initial conditions
are chosen. For the more detailed structures of the classical representation
must in time dissolve into lasting chaos, and I would think rather quickly.
Still, this does depend on a point I don’t actually know the answer to.
For in order to make the universe behave as you wish, that is to say, give a
good account of continents rifting and hummingbirds feeding, it might be that
to obtain each additional second of the desired history it is always sufficient
to correctly calculate another, say, thousand decimal places for the positions
and momenta of the molecules in the initial state. Or to the contrary, the first
thousand decimal places might give you one second, the next thousand only a
further half second, then a fourth of a second, and so on.
Yet even if I am wrong in this, we would just go from Scylla to Charybdis.
For in that case classical mechanics must be like the red dust theory, where,
from our point of view, anything is possible, or too close to anything. In either
case the classical solution set would imply no structure such as we experience
in life. No sculpted dunes, no ants carting morsels, no shower of hail would
pop out of it. Nor can one imagine any reason why the solution set would show
a preference for depicting creatures learning classical mechanics, or if so doing
benefiting by it. In short, there is a total disconnect between the fact that
classical mechanics is useful and the hypothesis that the universe as a whole is
a classical mechanical system.
That leaves us with an unsolved mystery: why does classical mechanics
work for us? And classical mechanics is the archetype of the kind of physics
where we learn from what can be true histories of things.
To my mind, the hourglass with observation of its emitted light is deeply
conservative physics. It makes quantum mechanics as seamless a continuation
of the physics of the previous centuries as is at all possible. This is because
of the mathematical form of the hourglass, which is a continuous development
from initial conditions, as well as the form of the observations, which impinge
as little as can be. And when this leads to our being given movies rather than
direct histories, then I am surprised (and amused) by this, but accept it for the
sake of the qualities mentioned, which I consider to be virtues that promise.
Nature is teaching us another lesson.
Bohr’s old quantum theory was based on quantum jumps, and I think
this was a wonderful piece of exploration in the dark. When Heisenberg’s new
quantum mechanics came along, quantum jumps were kept. The jumps would
allow direct histories to be retained as the foundation of our physics, though
at the expense of the continuous Hamiltonian evolution of the wave packets
(and at the expense of clear definition, for no one has ever been able to specify
just when and where and what the quantum jumps are). Like Schrödinger,
I am jarred by this. If we are given the choice of preserving philosophical
principle or mathematical form, I think we should prefer mathematical form.
Isn’t this what Copernicus did?
A final thought: If learning from the movies provided by hourglasses is
how we do physics, then to know why quantum mechanics works would be to
know why all the inferences we might make from the movies will fit together
with sufficient coherence. But to know this would require that we know all the
things we might ever think of. It’s hopeless. Though we might nibble at the
problem, by showing that the hourglasses have some needed characteristics. So
I think the hourglasses will leave us with an essentially unfathomable mystery.
References
Gibbs, J. Willard [1981]: Elementary Principles in Statistical Mechanics , Wood-
bridge, CT: Ox Bow Press, p. 17 and p. 163.
Mahon, Basil [2003]: The Man Who Changed Everything, Chichester, UK:
John Wiley & Sons Ltd.
The hourglasses suggest that von Neumann’s measurement theory should be
recast for imaginative use rather than for the description of actual situations.
This gives one extra freedom in setting it up, and it can then work more
effectively. An outline is here:
McCartor, Donald [2004]: ‘Quantum Thought Experiments Can Define Na-
ture’, Concepts of Physics, Vol. I, no. 1, pp. 105–150 and quant-ph
0702192.
[email protected]
|
0704.0421 | The Sigma-D Relation for Planetary Nebulae: Preliminary Analysis | Serb. Astron. J. } 174 (2007), 73 - 76
Preliminary report
THE Σ − D RELATION FOR PLANETARY NEBULAE:
PRELIMINARY ANALYSIS
D. Urošević1, B. Vukotić2, B. Arbutina1,2 and D. Ilić1
1Department of Astronomy, Faculty of Mathematics, University of Belgrade
Studentski trg 16, 11000 Belgrade, Serbia
2Astronomical Observatory, Volgina 7, 11160 Belgrade 74, Serbia
(Received: February 22, 2007; Accepted: March 30, 2007)
SUMMARY: An analysis of the relation between radio surface brightness and
diameter, so-called Σ−D relation, for planetary nebulae (PNe) is presented: i) the
theoretical Σ − D relation for the evolution of bremsstrahlung surface brightness
is derived; ii) contrary to the results obtained earlier for the Galactic supernova
remnant (SNR) samples, our results show that the updated sample of Galactic
PNe does not severely suffer from volume selection effect - Malmquist bias (same
as for the extragalactic SNR samples) and; iii) we conclude that the empirical
Σ −D relation for PNe derived in this paper is not useful for valid determination
of distances for all observed PNe with unknown distances.
Key words. planetary nebulae: general – Radio continuum: ISM – Methods: ana-
lytical – Methods: statistical
1. INTRODUCTION
The relation between radio surface bright-
nesses and diameters of supernova remnants (SNRs),
the so-called Σ−D relation, has been subject of the
extensive discussions in the last more than fourty
years. Due to improvements of the observational
techniques (radio-interferometers), the several hun-
dreds planetary nebulae (PNe) were resolved in the
last two decades at radio frequencies, but the Σ−D
relation for PNe was not discussed until now. By
using radio data, some statistical methods were es-
tablished in order to determine distances to PNe.
The main method was related to the correlation be-
tween radius of PNe and brightness temperature –
R − Tb relation (Van de Steene and Zijlstra 1995,
Zhang 1995, Phillips 2002). The different samples
of Galactic PNe with known distances were defined
in these papers. All the obtained empirical R − Tb
relations were used for determination of distances to
PNe for which the independent distances (in order of
R− Tb dependence) were not obtained earlier.
The samples of Galactic PNe are better for
statistical analysis than the samples of Galactic
SNRs. The selection effects should be smaller in
the case of PN samples. However, the selection ef-
fects surely influence the Galactic PN samples and
the statistical determination of distances to Galactic
PNe has to be highly uncertain.
The main objectives of this paper are the fol-
lowing:
i) to derive a simple form of the theoretical
Σ−D relation for PNe by analyzing the evolution of
radio bremsstrahlung surface brightness,
ii) to discuss whether the updated sample of
radio PNe is affected by the selection effects, and,
iii) to check whether the Σ−D relation is valid
for determination of distances to PNe.
http://arxiv.org/abs/0704.0421v3
D. UROŠEVIĆ et al.
2. ANALYSIS AND RESULTS
2.1. Theoretical Σ−D relation for PNe
The thermal bremsstrahlung mechanism is re-
sponsible for radiation of HII regions at radio wave-
lengths. The bremsstrahlung volume emissivity εν of
a PN can be shown to be (Rohlfs and Wilson 1996):
εν [ergs s
−1 cm−3 Hz−1] ∝ n2T−1/2, (1)
where n is the volume density and T is the thermo-
dynamic temperature of interstellar medium (ISM).
The surface brightness can be expressed as:
Σν ∝ ενD, (2)
where D is the diameter of PN. Combining Eqs. (1)
and (2), we obtain:
Σν ∝ n
D. (3)
Our next step is to express dependance of n and T
on D. For a constant velocity mass flow the den-
sity distribution is ̺ = Ṁ
4πr2v
, i.e. n ∝ D−x, where
x = 2. Moreover, for the isothermal envelope with
a power-law electron density distribution there is re-
lationship between the shape of the density distri-
bution and the power-law index of the radio contin-
uum spectra (see Gruenwald and Aleman 2007, and
references therein). Supposing that n ∝ D−2 and
T=const. (HII regions are approximately isothermal
at T ∼ 104 K), we obtain the simplest form of the
theoretical Σ−D relation for PNe:
Σν ∝ D
−3. (4)
This is a standard power-law form of the Σ − D
relation which can be written in general form as
Σ = AD−β , that is the same as in the case of SNRs.
It is possible that x in density distribution is
slightly higher, x & 2, and that the temperature is
not strictly constant throughout the nebula. We can
expect to see temperature gradients in PNe arising
from radiation hardening. More energetic photons
will travel further and when they are absorbed by the
PN they will impart greater kinetic energy to the ions
thereby producing a higher temperature. Using the
numerical model results given by Evans and Dopita
(1985), we calculate the dependence between logTe
and logD and find the low slope (≈ 0.1). Therefore,
this only slightly changes the slope of the theoretical
Σ−D relation. The value β = 3 is then a theoretical
lower limit, and the Σ − D relation could only be
steeper, as one can see from Eq. (3).
2.2. The empirical Σ−D relation for PNe
The most important prerequisite for deriving
a proper empirical Σ − D relation is defining of a
representative sample of PNe. The distances to the
calibrators have to be determined by accurate meth-
ods, e.g. trigonometric or spectroscopic parallaxes of
central stars in PNe, or by a method that uses the
expansion of nebulae. On the other hand, all sam-
ples suffer from the severe selection effects that arise
from limitation in sensitivity and resolution, but the
most severe selection effect for the Galactic samples
of PNe is Malmquist bias; i.e. intrinsically bright
PNe are favored because they are sampled from a
larger spatial volume compared to any given flux lim-
ited survey. The result is a bias against low surface
brightness nebulae such as highly evolved old PNe.
In this paper we use the updated sample of PNe at
the distances less than 0.7 kpc collected by Phillips
(2002). The influence of Malmquist bias in this sam-
ple is limited because of the limitation in distances to
PNe. In addition, we assume that the distances are
accurately determined for this sample of relatively
close PNe. The empirical Σ − D relation at 5 GHz
for 44 calibrators with distances less than 0.7 kpc
(Phillips 2002) has the form:
Σ56Hz = 2.33
+0.88
−0.64 · 10
−2.07±0.19
. (5)
The parameters A and β are calculated by least-
squares fitting procedure with correlation coefficient
−0.86. The corresponding Σν −D diagram is shown
in Fig. 1.
0.01 0.1 1 10
D [pc]
Fig. 1. The Σ−D diagram at 5 GHz for 44 Galac-
tic PNe with distances less than 0.7 kpc.
The form of Eq. (5) is very close to the so-
called trivial Σ−D form with β = 2 (for details see
Arbutina et al. 2004). The additional test in order to
estimate the validity of Eq. (5) pertains to the possi-
ble dependence between the luminosity and diameter
of PNe. The Lν−D diagram is shown in Fig. 2. The
scatter in Lν − D plane shows that the correlation
between Lν and D is poor (correlation coefficient = -
0.06) and therefore the physical dependence between
L and D could not be confirmed by this statistical
procedure.
THE Σ−D RELATION FOR PLANETARY NEBULAE
0.01 0.1 1 10
D [pc]
Fig. 2. The L −D plot at 5 GHz for 44 Galactic
PNe with distances less than 0.7 kpc.
3. DISCUSSION
The theoretical Σν −D relation (Eq. (4)) for
PNe, derived in this paper, describes a trend of de-
creasing radio surface brightness with increasing di-
ameter of an object. The radiation mechanism used
in this simple derivation is thermal bremsstrahlung.
This is the basic process of production of the radio
radiation in HII regions. The theoretically derived
slope (β = 3) is steeper than the slope from the em-
pirical relation given by Eq. (5). This discrepancy
can be explained by the low quality of the sample of
Galactic PNe or by the assumptions used in deriva-
tion of theoretical relation. Due to small variation
in power-law density distribution with x & 2 (Gru-
enwald and Aleman 2007, and references therein)
and approximately constant temperature of expand-
ing envelope of PNe, theoretical slope can be slightly
steeper than in Eq. (4). Therefore, we conclude
that the theoretical relation has the correct form,
but our empirical relation is under influence of bi-
ases that could make the slope shallower. On the
other hand, there are some attempts to show that
evolution of PNe are not linear in log-log scales (e.g.
Phillips 2004). These different dependences cannot
be derived from the thermal bremsstrahlung radia-
tion formula (Eq. (1)).
A very interesting feature regarding the em-
pirical relation for Galactic PNe (Eq. (5)) is that
the slope is approximately equal to the slope of triv-
ial Σ − D relation. Therefore, we conclude that
Malmquist bias is not so severe as in cases of Galac-
tic SNR samples. This slope (β ≈ 2) was obtained
for the extragalactic samples of SNRs (except M82
sample) where Malmquist bias is small, because all
the SNRs are at the approximately same distance
(see Urošević 2002, Urošević et al. 2005).
The large scatter in Lν − D plane (Fig. 2)
suggests that the slope in Eq. (5) does not have
real and valid physical interpretation. It is a kind
of luminosity-diameter scattering artefact which pro-
duces the trivial Σ ∝ D−2 form. Therefore, the rela-
tion defined by Eq. (5) is not precise enough for de-
termination of valid distances to Galactic PNe. This
is due to the different biases: the limitations in sen-
sitivity and resolution of radio surveys, the source
confusion, Malmquist bias (in mild form), mixture
of different types of PNe in the same sample, and
insufficient precision in determining the distances to
the 44 calibrators.
4. SUMMARY
The main results of this paper may be sum-
marized as follows:
i) The theoretical Σν −D relation for the radio
evolution of thermal bremsstrahlung surface
brightness of PNe in form of Σν ∝ D
−3 is de-
rived.
ii) Our results show that the updated sample of
Galactic PNe does not severely suffer from vol-
ume selection effect - Malmquist bias (same
as in cases of the extragalactic SNR samples).
This is opposite to results obtained earlier for
the Galactic SNR samples.
iii) Due to analysis of the Lν −D dependence, we
conclude that the Σν −D relation for Galactic
PNe is not useful for reliable determination of
distances for all observed PNe with unknown
distances.
The above observation leads to the more gen-
eral comment that PNe may have very different ini-
tial conditions leading to independent evolutionary
paths. These paths could follow the same theoreti-
cal Σ−D curve but with varying intercepts, leading
to the scatter such as the one found in this paper.
Acknowledgements – The authors would like to thank
the referee Prof. Nebojsa Duric for valuable com-
ments which have improved this paper. This research
has been supported by the Ministry of Science and
Environmental Protection of the Republic of Serbia
(Projects: No 146002, No 146003, No 146012, No
146016).
D. UROŠEVIĆ et al.
REFERENCES
Arbutina, B., Urošević, D., Stanković, M. and Tešić,
Lj.: 2004, Mon. Not. R. Astron. Soc., 350,
Evans, I.N. and Dopita, M.A.: 1985, Astrophys. J.
Suppl. Series, 58, 125
Gruenwald, R. and Aleman, A.: 2007, Astron. As-
trophys., 461, 1019.
Phillips, J.P.: 2002, Astrophys. J. Suppl. Series,
139, 199.
Phillips, J.P.: 2004, Mon. Not. R. Astron. Soc.,
353, 589.
Rohlfs, K. and Wilson, T.L.: 1996, Tools of Radio
Astronomy, Springer
Urošević, D.: 2002, Serb. Astron. J., 165, 27
Urošević, D., Pannuti, T. G., Duric, N., Theodorou,
A.: 2005, Astron. Astrophys., 435, 437.
Van de Steene, G.C. and Zijlstra, A.A.: 1995, As-
tron. Astrophys., 293, 541.
Zhang, C.Y.: 1995, Astrophys. J. Suppl. Series, 98,
Σ − D RELACIJA ZA PLANETARNE
MAGLINE: PRELIMINARNA ANALIZA
D. Urošević1, B. Vukotić2, B. Arbutina1,2 and D. Ilić1
1Department of Astronomy, Faculty of Mathematics, University of Belgrade
Studentski trg 16, 11000 Belgrade, Serbia
2Astronomical Observatory, Volgina 7, 11160 Belgrade 74, Serbia
UDK 524.37–77–54
Prethodno saopxtenje
Prikazana je analiza tzv. Σ − D re-
lacije izme�u povrxinskog sjaja na radio-
frekvencijama i dijametra planetarnih
maglina (PM): i) izvedena je teorijska Σ −D
relacija za evoluciju povrxinskog sjaja
stvorenog zakoqnim zraqenjem; ii) suprotno
rezultatima dobijenim ranije za uzorke
saqinjene od Galaktiqkih ostataka super-
novih, naxi rezultati pokazuju da najnovije
formirani uzorak Galaktiqkih PM ne trpi
veliki uticaj zbog zapreminskog selekcionog
efekta, tzv. Malmkvistovog selekcionog
efekta (isto vaжi za vangalaktiqke uzorake
ostataka supernovih); i iii) zakljuqujemo da
empirijska Σ − D relacija za PM izvedena
u ovom radu nije upotrebljiva za pouzdana
odre�ivanja daljina do svih posmatranih PM
sa nepoznatim daljinama.
|
0704.0422 | Polarization conversion in a silica microsphere | Polarization conversion in a silica microsphere
Pablo Bianucci, Chris Fietz, John W. Robertson, Gennady Shvets, and Chih-Kang Shih∗
Physics Department, The University of Texas at Austin, Austin, Texas 78712
(Dated: May 22nd, 2007)
Abstract
We experimentally demonstrate controlled polarization-selective phenomena in a whispering
gallery mode resonator. We observed efficient (≈ 75%) polarization conversion of light in a silica
microsphere coupled to a tapered optical fiber with proper optimization of the polarization of the
propagating light. A simple model treating the microsphere as a ring resonator provides a good fit
to the observed behavior.
In the past few years, microresonators
have received a lot of attention1. Whis-
pering gallery mode (WGM) resonators2,
such as microspheres,3 microtoroids4 and
microrings5 have been the object of inten-
sive research, both in their fundamental prop-
erties (such as quality factors, non-linear
effects6,7 and coupling to quantum systems8
among many) and applications that include
lasers9,10, chemical11 and biological12 sens-
ing and photonic devices13. Microsphere res-
onators, particularly when coupled to a ta-
pered optical fiber14,15, are very useful to
characterize these properties and test new
ideas due to their high Q-factors and ease of
fabrication.
Recent reports have shown a further step,
taking into account the difference between
modes with different polarizations in micro-
spheres. In particular, changes in the output
polarization after coupling into the resonator
have been observed16 and transverse electric
(TE) and transverse magnetic (TM) modes
have been discriminated17.
Polarization conversion has been observed
in microrings5 and explained as a resonant
enhancement of polarization coupling caused
by waveguide bending. However, the mode
structure of microspheres makes it possible
to completely decouple the polarizations and
still obtain conversion. In this article, we re-
port on the observation of efficient, controlled
polarization conversion by using a silica mi-
crosphere resonator coupled to a tapered op-
tical fiber. We demonstrate that highly ef-
ficient polarization conversion (75% for our
particular case, higher for better optimized
conditions) is enabled by a specific orienta-
tion between the incoming light polarization
and fiber-resonator displacement. Specifi-
http://arxiv.org/abs/0704.0422v2
cally, for a horizontally stacked, strongly cou-
pled, fiber and resonator combination, a 45◦
incident polarization results in the largest
conversion. The conversion results in a
strong dip of the transmitted light with the
original polarization and a strong spike in the
orthogonally polarized transmission.
We fabricated the tapered fiber using the
“flame brush” technique18. This technique
involves mechanically stretching the optical
fiber while scanning a flame (oxy-hydrogen
in our case) over the region to be tapered.
Due to constraints in the maximum pulling
length, the fiber tapers are not completely
adiabatic, but typical losses are never larger
than 50%. SEM studies of the tapers reveal
a characteristic diameter close to 1 µm. The
microsphere was fabricated using a CO2 laser
to stretch and melt an optical fiber tip19. In
this way it is easy to obtain spheres with di-
ameters ranging from 10 µm to 200 µm. For
this particular experiment the sphere diame-
ter was measured using an optical microscope
to be 52 µm (corresponding to an estimated
free spectral range of 1.2 THz).
We mounted the microsphere on a piezo-
electric scanner which allowed us to finely
position the sphere over a range of a few
micrometers, and the stretched fiber taper
on a piezoelectric stick-slip walker permitting
both coarse and fine positioning of the fiber
taper next to the sphere. Both sphere and
PD PC2 PC1
Tapered
Fiber
Laser PR FC
Microsphere
10 um
FIG. 1: Experimental setup schematic. PR is
a polarization rotator, FC a fiber coupler, PC1
and 2 are fiber polarization controllers, P a po-
larizer and PD is an amplified photodiode. Inset:
Image of a sphere near a tapered fiber.
taper were then situated inside a compact,
closed chamber. We used an external cav-
ity tunable diode laser purchased from New
Focus as the excitation source, centered at
a wavelength near 927.85 nm. The polariza-
tion rotator set the polarization of the laser
which was then coupled into the optical fiber
using a free-space coupler. A polarizer and
an amplified photodiode at the fiber output
were used to analyze the transmitted light.
Space constraints in the chamber and lim-
itations on the arrangement of the optical
fiber caused bending of the fiber in differ-
ent locations and subsequent scrambling of
the input polarization. As a way to compen-
sate for these changes in the polarization, we
used two polarization controllers. The first
one preceded the fiber taper, compensating
for polarization changes up to the position of
the microsphere. The second controller was
placed after the fiber taper to ensure the lin-
earity of the output polarization. Figure 1
shows a schematic of this experimental setup.
We used the following procedure to mea-
sure the degree of polarization conversion.
First, the incoming polarization was selected
by using the polarization rotator. Then we
adjusted the first polarization controller to
ensure the polarization at the fiber taper was
linear and matched to one set of modes (“x-
polarized”). The next step was to uncouple
the taper from the sphere and make sure the
output polarization was linear (we achieved
this by turning the detection polarizer to its
position for minimum transmission and then
minimizing this transmission further with the
second polarization controller). This orienta-
tion of the detection polarizer is the one we
call “orthogonal”. Rotating the polarizer 90
degrees (the “parallel” orientation) resulted
in maximum transmission, with a contrast
of about 95%, confirming the linear polar-
ization of the output light. Finally, we posi-
tioned the sphere and the tapered fiber try-
ing to optimize the coupling, while measur-
ing transmission spectra for both orientations
of the detection polarizer. We repeated the
procedure for two other incoming polariza-
tions: one matched to the other set of sphere
modes (“y-polarized”) and another at 45◦ be-
tween the x- and y- polarization axis (“xy-
polarized”).
Figure 2 shows the resulting transmission
-30 -25 -20 -15 -10 -5 0 5 10 15 20 25 30
Frequency shift (GHz)
Input Detection
polariz. polariz.
FIG. 2: Transmission spectra for different in-
put polarizations. The resonant frequencies cor-
respond to modes with l ≈ 496. The x- and
y-polarizations are orthogonal and correspond
to the polarization eigenmodes of the resonator.
The xy-polarization is oriented at 45 degrees
from both x and y. The dark traces correspond
to the detection polarizer parallel to the input
polarization and the light traces correspond to a
crossed detection polarization.
spectra for the different configurations. The
cases for both the x- and y- polarized light
show the same behavior: a set of transmission
dips whenever the laser frequency hit a whis-
pering gallery resonance when the detection
polarization is parallel and no signal when
it is perpendicular. The xy-polarized case is
more interesting: the parallel detection po-
larization shows dips for both sets of modes,
while the orthogonal one shows transmission
peaks at the whispering gallery resonances.
At the highest peak, more than 70% of the
incident light had its polarization converted.
Most of the observed polarization conver-
sion can be understood by using a simple ring
resonator model for the whispering gallery
modes. In this model, the transmission of
polarized light through the resonator is given
by14,20
τ(φ) =
r − aeiφ
1− raeiφ
, (1)
where r is the field coupling coefficient be-
tween the resonator and the waveguide, a is
the attenuation due to the resonator intrinsic
losses and φ = 2π(ν−ν0)tRT is the phase shift
imposed by the resonator (ν and ν0 are the in-
coming light frequency and the resonant fre-
quency respectively, while tRT is the round-
trip time in the resonator). The model is
scalar, but we can include the polarization by
simply assuming that modes with orthogonal
polarizations are independent and neglecting
cross-polarized couplings (using an analysis
similar to that by Little and Chu21). In this
way we obtain the same expression, with pos-
sibly different parameters, for the transmis-
sion of both polarizations. In our particu-
lar case of whispering gallery modes in mi-
crospheres, we can safely assume that modes
with different polarizations are not degener-
ate, so one of the polarizations will be unaf-
fected by the presence of a resonance. This
differs from the case of microrings5, where the
conversion depends on coupling between TE
and TM modes.
The essence of the effect lies in the differ-
ent resonator response for each polarization.
For a strongly coupled fiber and microsphere,
|τ | ≈ 1, but the phase shift ψ = arg(τ)
is changed by ∆ψ = π as the frequency is
sweeped across the resonance. Because the
orthogonal polarization is transmitted unal-
tered, the transmitted polarization rotates by
as much as 90◦ for the initial xy-polarization.
When the fiber and the resonator are horizon-
tally stacked, the effect is maximized when
the incident polarization is at 45◦ degrees
with respect to the horizontal plane.
Conversion efficiencies of up to 25% can
be achieved if one of the polarizations is crit-
ically coupled to the ring, i.e. is completely
absorbed in+the resonator. Achieving higher
efficiencies requires increasing the resonator-
waveguide coupling to obtain a significant po-
larization dependent phase shift which will
change the final polarization state into one
closer to the desired one.
We can look in more detail at the data
by concentrating into a pair of modes show-
ing good conversion, now accounting for laser
frequency drift between scans using a Fabry-
Perot interferometer as a reference. This de-
tailed spectrum can be seen in in Fig. 3. The
resonance on the right side of Fig. 3, near a
shift of 31 GHz, shows a polarization conver-
sion of about 60%. The left-side resonance
shows a conversion near 75%. The higher ef-
ficiency is due to the leftmost mode being
26 27 28 29 30 31
Frequency shift (GHz)
30.5 31
26 26.5 27 27.5 28
FIG. 3: Detailed view of two modes showing
polarization conversion. The dashed lines are
fits using equations of the form of equation 1.
The fit parameters for the leftmost features are
a = 0.99997, r = 0.99977. The corresponding
ones for the rightmost feature are a = 0.99999,
r = 0.99993.
more strongly coupled (displaying a broader
feature) to the tapered fiber than the right-
most one. Consistent with theoretical predic-
tions, in both cases one of the polarizations
is over-coupled to the ring. The lack of a
shift in the center frequency of the features
also indicates that each pair of peak and dip
corresponds to a single resonant mode.
This phenomenon could be useful for po-
larization control in photonic devices, such
as narrowband polarization-dependent filter-
ing or switching, as shown in Fig. 4 or even
arbitrary polarization manipulation.
We have observed efficient polarization
conversion on a microsphere resonator cou-
PBSPBS
a) b)
FIG. 4: Schematic of a resonator working as
wavelength-selective polarization switch. a)
Two signals with different wavelengths+ (green
and blue) and orthogonal polarizations pass un-
changed through the waveguide and the un-
coupled resonator. A polarization beamsplitter
then routes the signals to different paths. b)
The polarization of the resonant signal (blue) is
converted by the coupled resonator, and both
signals are sent through the same path. The
resonator-waveguide coupling can be changed in
different ways, including mechanical or optical22
means.
pled to a tapered optical fiber and used a
simple theoretical model to understand the
phenomenon. The model does not involve di-
rect coupling of the orthogonal polarizations,
but rather a polarization-selective phase shift
induced by the resonator. This effect should
be common to all whispering gallery mode
resonators and could be useful for polariza-
tion control in photonic devices.
Acknowledgments
This work was supported by NSF-NIRT
(DMR-0210383), the Texas Advanced Tech-
nology program, and the W.M. Keck Foun-
dation. G.S. and C.F. acknowledge support
from ARO MURI grant no. W911NF-04-01-
0203.
∗ Electronic address: [email protected]
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|
0704.0423 | Limits on WIMP-nucleon interactions with CsI(Tl) crystal detectors | Limits on the WIMP-nucleon interactions with CsI(Tl) crystal detectors
H.S. Lee,1 H.C. Bhang,1 J.H. Choi,1 H. Dao,7 I.S. Hahn,4 M.J. Hwang,5 S.W. Jung,2 W.G. Kang,3 D.W.
Kim,1 H.J. Kim,2 S.C. Kim,1 S.K. Kim,1, ∗ Y.D. Kim,3 J.W. Kwak,1, † Y.J. Kwon,5 J. Lee,1, ‡ J.H. Lee,1 J.I.
Lee,3 M.J. Lee,1 S.J. Lee,1 J. Li,7 X. Li,7 Y.J. Li,7 S.S. Myung,1 S. Ryu,1 J.H. So,2 Q. Yue,7 and J.J. Zhu7
(KIMS Collaboration)
DMRC and Department of Physics and Astronomy, Seoul National University, Seoul, Korea
Department of Physics, Kyungpook National University, Daegu, Korea
Department of Physics, Sejong University, Seoul, Korea
Department of Science Education, Ewha Womans University, Seoul, Korea
Department of Physics, Younsei University, Seoul, Korea
Department of Engineering Physics, Tsinghua Universuty, Beijing, China
Department of Engineering Physics, Tsinghua University, Beijing, China
(Dated: November 4, 2018)
The Korea Invisible Mass Search (KIMS) experiment presents new limits on the WIMP-nucleon
cross section using data from an exposure of 3409 kg·d taken with low-background CsI(Tl) crystals
at Yangyang Underground Laboratory. The most stringent limit on the spin-dependent interaction
for a pure proton case is obtained. The DAMA signal region for both spin-independent and spin-
dependent interactions for the WIMP masses greater than 20 GeV/c2 is excluded by the single
experiment with crystal scintillators.
PACS numbers: 95.35.+d, 14.80.Ly
The existence of dark matter has been widely sup-
ported by many astronomical observations on vari-
ous scales [1][2][3]. Weakly interacting massive parti-
cles (WIMPs) are a good candidate for dark matter well
motivated by cosmology and supersymmetric models [4].
The Korea Invisible Mass Search (KIMS) experiment has
developed low-background CsI(Tl) crystals to detect the
signals from the elastic scattering of WIMP off the nu-
cleus [5][6][7]. Both 133Cs and 127I are sensitive to the
spin-independent (SI) and spin-dependent (SD) interac-
tions of WIMPs. Recently, the role of CsI in the direct
search for SD WIMP for pure proton coupling has been
pointed out [8]. It is worth noting that 127I is the dom-
inant target for the SI interactions in the DAMA exper-
iment. The pulse shape discrimination (PSD) technique
allows us to statistically separate nuclear recoil (NR) sig-
nals of WIMP interactions from the electron recoil (ER)
signals due to the gamma ray background [9][10].
The KIMS experiment is located at the Yangyang Un-
deground Laboratory (Y2L) at a depth of 700 m under an
earth overburden. Details of the KIMS experiment and
the first limit with 237 kg·d exposure data can be found
in the previous publication [11]. Four low-background
CsI(Tl) crystals are installed in the Y2L and operated
at a temperature of T = 0◦C. Throughout the exposure
period, the temperature of the detector was kept sta-
ble to within ±0.1◦C. Green-enhanced photomultiplier
tubes (PMTs) are mounted at both ends of each crystal.
The signals from the PMTs are amplified and recorded
by a 500 MHz FADC. Each event is recorded for a pe-
riod of 32 µs. Both PMTs on each crystal must have at
least two photoelectrons within a 2 µs window to form an
event trigger. We obtained 3409 kg·d WIMP search data
TABLE I: Crystals used in this analysis and amount of data
for each crystal
Crystal mass (kg) data (kg·days)
S0501A 8.7 1147
S0501B 8.7 1030
B0510A 8.7 616
B0510B 8.7 616
Total 34.8 3409
with four crystals, as shown in Table I. The energy is cali-
brated using 59.5 keV gamma rays from an 241Am source.
For calibration of the mean time, a variable used for the
PSD, NR events are obtained with small crystals ( 3 cm ×
3 cm × 3 cm ) using an Am-Be neutron source. Compton
scattering events taken with the WIMP search crystals
using the 137Cs source are used to determine the mean
time distribution of the gamma background. Compton
scattering events are also taken with the small crystals
to verify that the mean time ditributions for both the test
crystals and the WIMP search crystals are the same. In
order to understand the nature of the PMT background,
a dominant background at low energies, acrylic boxes are
mounted on the same PMTs used for the crystals. The
data obtained using this setup is used to develop the cuts
for the rejection of PMT background.
Since the decay time of the scintillation light in the
CsI(Tl) crystal is rather long, photoelectrons are well
separated at low energies and thereby enabling recon-
struction of each photoelectron. The time distribution of
photoelectrons in an event is fitted to a double exponen-
http://arxiv.org/abs/0704.0423v2
sec)µMean Time (
FIG. 1: (color online). MT distribution of NR events (open
squares), ER events (open circles) and WIMP search
data (filled triangles) of S0501A crystal in the 5-6 keV range.
Fitted PDF functions are overlayed. χ2/DOF =0.8 and 1.3
with DOF=38 and 35 for NR and ER events respectively.
tial function given by
f(t) =
−(t− t0)
−(t− t0)
where τf and τs are decay time constants of fast and
slow components, respectively, R is ratio between two
components, and t0 is the time of the first photoelectron
in the event. The mean time (MT ) of each event is then
calculated using these quantities as
t · f(t)dt/
f(t)dt.
With this method, an improvement in PSD is achieved
over the previous analysis where we used a simple math-
ematical mean [11]. In order to reject the PMT back-
ground, we applied cuts to the fit variable, τf . The ratio
between the maximum log likelihood value of the dou-
ble exponential fit and that of the single exponential fit
is also used to reject the PMT background, since PMT
background events tend to be shaped as single exponen-
tial decay. To reject the background that originates from
the radioactivity of the PMT, the asymmetry between
the signals from two PMTs is applied. Finally events in
which signals are recorded in more than one crystal are
rejected. The event selection efficiency was estimated
by applying the same analysis cuts to the neutron and
gamma calibration samples. The efficiency depends on
the measured energy and ranges from 30% at 3 keV to
60% above 5 keV.
The estimation of the NR event rate is performed in
each 1 keV bin from 3 to 11 keV for each crystal. TheMT
distributions of NR events and ER events are compared
with the WIMP search data in Fig. 1 for the 5-6 keV
energy range. The probability density functions (PDF)
for the ER and NR events are obtained by fitting these
distributions. An unbinned maximum likelihood fit is
Electron Equivalent Energy (keV)
3 4 5 6 7 8 9 10 11
FIG. 2: (color online). Extracted NR event rates of the
S0501A (open circles), S0501B (filled circles), B0510A (filled
squares), and B05010B (filled triangles) crystals and only sta-
tistical errors (1σ) are shown. The points are shifted with
respect to each other on the x-axis to avoid overlapping.
performed with the log(MT ) distribution of the WIMP
search data using the likelihood function,
× exp{−(NNR,i +NER,i)}
[NNR,iPDFNR,i(xk) +NER,iPDFER,i(xk)],
where the index i denotes the i-th energy bin; n =
NNR,i +NER,i is the total number of events; NNR,i and
NER,i are the numbers of NR and ER events, respec-
tively; PDFNR,i and PDFER,i are PDFs of NR and ER
events, respectively; and xk = log(MT ) for each event.
The NR event rates obtained for each bin and for each
crystal after efficiency correction are shown in Fig. 2.
The extracted NR event rates are consistent with a null
observation of the WIMP signal.
In order to obtain the expected measured energy spec-
trum of a WIMP signal including instrumental effects,
a Monte Carlo (MC) simulation with GEANT4 [12] is
used. A recoil energy spectrum is generated for each
WIMP mass with the differential cross section, form fac-
tor, and quenching factor, as described in Ref. [13]. The
spin-dependent form factor for 133Cs calculated by Toiva-
nen [14] is used, while for 127I, Ressell and Dean’s cal-
culation [15] is used. The photons generated with the
fitted decay function described above are propagated to
the PMT and digitized in the same manner as in the ex-
periment. Subsequently, the photoelectrons within given
time windows are counted to check the trigger condition
and to calculate energy. In this manner, the trigger ef-
ficieny and energy resolution is accounted for in the ex-
pected energy spectrum. The trigger efficiency is found
to be higher than 99% above 3 keV. The simulation is
verified with the energy spectrum obtained using 59.5
keV gamma rays from 241Am. The peak position and
TABLE II: Spin expectation values for 133Cs and 127I
Isotope J < Sp > < Sn > Reference
133Cs 7/2 -0.370 0.003 [16]
127I 5/2 0.309 0.075 [15]
width of the distribution are very well reproduced for
each crystal as described in Ref [11].
The total WIMP rate, R, for each WIMP mass is ob-
tained by fitting the measured energy spectrum to the
simulated one. The 90% confidence level (CL) limit on
R is calculated by the Feldman-Cousins’s approach in
the case of Gaussian with a boundary at the origin [17]
and then converted to the WIMP-nucleus cross section,
σW−A. Subsequently, the limits on WIMP-nucleon cross
section is obtained from Ref. [13][18] as follows:
σW−n = σW−A
where µn,A are the reduced masses of the WIMP-nucleon
and WIMP-target nucleus of mass number A. CA/Cn =
A2 for SI interactions and CA/Cn = 4/3{ap < Sp >
+an < Sn >}2(J + 1)/J for SD interactions. Here ap,
an are WIMP-proton and WIMP-neutron SD couplings
respectively. The spin expectation values used for this
analysis are shown in Table II. Following the “model-
independent” framework [18], we report the allowed re-
gion in two cases for SD interaction: one for an = 0,
and the other for ap = 0. We express the WIMP-nucleon
cross section as follows:
σSIW−n = σW−A
σSDW−n,p = σW−A
µ2n,p
(J + 1)
< Sn,p >2
where we indicate pure proton (p, an = 0) and pure
neutron (n, ap = 0) coupling for SD interaction. We also
present the allowed region in the ap − an plane with the
following relation [18]:
where GF is the Fermi coupling constant.
The uncertainty in the MT distribution results in the
uncertainty of the NR event rate. The limited statistics
of the calibration data and different crystals used for the
neutron calibration and WIMP search data are the ma-
jor sources of this uncertainty. The former is investigated
by varying the fitted parameters in PDF function within
errors. The lattter is estimated by changing the mean of
MT by the difference between the crystals. The system-
atic uncertainties from these two souces are combined in
quadrature resulting in 20-30% of statistical uncertain-
ties depending on the energy bins. In addition, there
WIMP Mass (GeV)
210 310 410
DAMA region
FIG. 3: (color online). Exclusion plot for the SD interaction
in the case of pure proton coupling (an = 0) at the 90%
confidence level
WIMP Mass (GeV)
210 310 410
DAMA region
FIG. 4: (color online). Exclusion plot for for the SD interac-
tion in the case of pure neutron coupling (ap = 0) at the 90%
confidence level
are uncertainties in the MC estimation of the expected
event rates due to the uncertainties in the quenching fac-
tors and the difference of energy resolution between the
MC simulation and the data. The systematic error from
the MC simulation is estimated to be 13.3% of the limits.
These systematic errors are combined with the statistical
error in quadrature in the presented results.
The limits on the SD interactions are shown in Fig. 3
and 4 in the cases of pure proton coupling and pure neu-
tron coupling, respectively. We also show the results ob-
tained from CDMS [19], NAIAD [20], SIMPLE [21], and
-6 -4 -2 0 2 4 6
FIG. 5: (color online). Allowed region (90% confidence level)
in ap − an plane by KIMS data (inside the solid line contour)
for 50 GeV WIMP mass. Results of CDMS [19](dotted line)
and NAIAD [20](dot-dashed line) are also shown.
WIMP Mass (GeV)
210 310 410
FIG. 6: (color online). Exclusion plot for the SI interactions
at the 90% confidence level.
PICASSO [22]. The DAMA signal region is taken from
Ref [23]. Our limit provides the lowest bound on the
SD interactions in the case of pure proton coupling for a
WIMP mass greater than 30 GeV/c2. The allowed region
in the ap − an plane for the WIMP mass of 50 GeV/c2 is
also shown in Fig. 5 together with the limits from CDMS
and NAIAD. The limit for the SI interactions is shown
in Fig. 6 together with the results of CDMS [24], EDEL-
WEISS [25], CRESST [26], ZEPLIN I [27], and the 3σ
signal region of DAMA (1-4) [28]. Although there are
several experiments that reject the DAMA signal region,
this is the first time that it is ruled out by a crystal de-
tector containing 127I, which is the dominant nucleus for
the SI interactions in the NaI(Tl) crystal.
In summary, we report new limits on the WIMP-
nucleon cross section with CsI(Tl) crystal detectors using
3409 kg·d exposure data. The DAMA signal regions for
both SI and SD interactions are excluded for the WIMP
masses higher than 20 GeV/c2 by the single experiment.
The most stringent limit on the SD interaction in the
case of purely WIMP-proton coupling is obtained.
The authors thank Dr. J. Toivanen and M. Korte-
lainen for the calculation of the SD form factor as well
as for the useful discussions. This work is supported by
the Creative Research Initiative Program of the Korea
Science and Engineering Foundation. We are grateful to
the Korea Middland Power Co. Ltd. and the staff mem-
bers of the YangYang Pumped Storage Power Plant for
providing us the underground laboratory space.
∗ [email protected]
† Current address: National Cancer Center, Ilsan, Korea
‡ Current address: Department of Physics, Ewha Womans
University, Seoul, Korea
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mailto:[email protected]
|
0704.0424 | Stopping effects in U+U collisions with a beam energy of 520 MeV/nucleon | Stopping effects in U+U collisions with a beam energy of 520 MeV/nucleon
Xiao-Feng Luo,1, ∗ Xin Dong,1 Ming Shao,1 Ke-Jun Wu,2 Cheng Li,1
Hong-Fang Chen,1 and Hu-Shan Xu3
University of Science and Technology of China, Hefei, Anhui 230026, China
Institute of Particle Physics, Hua-Zhong Normal University, Wuhan, Hubei 430079, China
Institute of Modern Physics, Chinese Academy of Sciences, LanZhou, Gansu 730000, China
(Dated: November 4, 2018)
A Relativistic Transport Model (ART1.0) is applied to simulate the stopping effects in tip-tip
and body-body U+U collisions, at a beam kinetic energy of 520 MeV/nucleon. Our simulation
results have demonstrated that both central collisions of the two extreme orientations can achieve
full stopping, and also form a bulk of hot, dense nuclear matter with a sufficiently large volume
and long duration, due to the largely deformed uranium nuclei. The nucleon sideward flow in
the tip-tip collisions is nearly 3 times larger than that in body-body ones at normalized impact
parameter b/bmax < 0.5, and that the body-body central collisions have a largest negative nucleon
elliptic flow v2 = −12% in contrast to zero in tip-tip ones. Thus the extreme circumstance and the
novel experimental observables in tip-tip and body-body collisions can provide a good condition and
sensitive probe to study the nuclear EoS, respectively. The Cooling Storage Ring (CSR) External
Target Facility (ETF) to be built at Lanzhou, China, delivering the uranium beam up to 520
MeV/nucleon is expected to make significant contribution to explore the nuclear equation of state
(EoS).
PACS numbers: 24.10.Lx,25.75.Ld,25.75.Nq,24.85.+p
I. INTRODUCTION
In recent years, the ultra-relativistic high energy heavy
ion collisions performed at SPS/CERN and RHIC/BNL
sNN ∼ 10 − 200 GeV) focus on high temperature
and low baryon density region in nuclear matter phase
diagram [1] to search a new form of matter with par-
tonic degree of freedom-the quark-gluon plasma (QGP)
[2, 3, 4, 5]. However, no dramatic changes of experimen-
tal observables, such as jet-quenching, elliptic flow and
strangeness enhancement, have been observed yet [6]. On
the other hand, the heavy ion collisions performed at the
BEVALAC/LBNL and SIS/GSI [7, 8] in last two decades
were used to produce hot and compressed nuclear mat-
ter to learn more about the nuclear equation of state
(EoS) [13, 14] at high baryon density and low temper-
ature region of the phase diagram. Although we have
made great efforts to study the nuclear EoS, theoreti-
cally and experimentally, a solid conclusion can hardly
be made. Then, it is still worthwhile to systematically
study on the collision dynamics as well as the EoS observ-
ables. Recently, for more understanding of the nuclear
matter phase diagram and EoS at high net-baryon den-
sity region, it is proposed to collide uranium on uranium
target at External Target Facility (ETF) of Cooling Stor-
age Ring (CSR) at Lanzhou, China with a beam kinetic
energy of 520 MeV/nucleon. [10].
Uranium is the largest deformed stable nucleus, and
has approximately an ellipsoid shape with the long and
short semi-axis given by Rl = R0(1 + 2δ/3) and Rs =
∗contact author: [email protected]
FIG. 1: (Color online) (a) body-body collisions (b) tip-tip
collisions
R0(1 − δ/3), respectively, where R0 = 7 fm is the effec-
tive spherical radius and δ = 0.27 is the deformation pa-
rameter [9]. Consequently, one has Rl/Rs = 1.3. In our
simulation, we consider two extreme orientations: the so-
called tip-tip and body-body patterns with the long and
short axes of two nuclei are aligned to the beam direc-
tion, respectively [12], see Fig. 1 for illustration. The
two types of orientations can be identified in random ori-
entations of U+U collisions by making proper cutoffs in
experimental data, such as the particle multiplicities, el-
liptic flow and so on [10, 11]. With the two extreme
collision orientations, some novel stopping effects which
are believed responsible for some significant experimental
observables, such as particle production, collective mo-
tion as well as attainable central densities, can be ob-
tained. Due to the large deformation of the uranium
nuclei [11, 12] , it is expected that the tip-tip collisions
can form a higher densities nuclear matter with longer
duration than in body-body or the spherical nuclei colli-
sions, which is considered to be a powerful tool to study
http://arxiv.org/abs/0704.0424v2
the nuclear matter phase transition at high baryon den-
sity [12], and the body-body central collisions may reveal
a largest out-of-plane elliptic flow (negative v2) at high
densities, which can be a sensitive probe to extract the
early EoS of the hot, dense nuclear matter [12, 17]. The
novel experimental observables can be effectively utilized
to study the possible nuclear matter phase transition and
the nuclear EoS [12, 13, 14, 15, 16, 17, 18, 19, 20]. For
comparing with tip-tip and body-body collisions, a type
of gedanken ”sphere-sphere” collisions without deforma-
tions of uranium nuclei are also included in the simula-
tion.
The ART1.0 model [21, 22] derived from Boltzmann-
Uehling-Uhlenbeck (BUU) model [23] has a better treat-
ment of mean field and Pauli-Blocking effects [23] than
cascade models [24]. The fragments production mech-
anism and partonic degree of freedom are not present
in the ART1.0 model. A soft EoS with compressibility
coefficient K = 200 MeV is used throughout the simu-
lation and the beam kinetic energy of uranium nuclei is
set to 520 MeV/nucleon if not specifically indicated. In
the next section, we discuss about the stopping power
ratio and selection of impact parameter b. In Sec. 3,
the evolution of baryon and energy densities as well as
thermalization of central collision systems are studied. In
Sec. 4, some experimental observables, such as nucleon
sideward flow and elliptic flow are also investigated. We
summarize our results in Sec. 5.
II. STOPPING POWER OF TIP-TIP AND
BODY-BODY COLLISIONS
Large stopping power can lead to remarkable pressure
gradient in the compressed dense matter. It is generally
also considered to be responsible for transverse collec-
tive motion [25], the maximum attainable baryon and
energy densities as well as thermalization of collision sys-
tems. Thus, the study of the stopping power in U+U
collisions may provide important information for under-
standing the nuclear EoS and collision dynamics.
A. Selection of impact parameter
The nuclear stopping power and geometric effects in
U+U collisions rely strongly on the impact parameter
b. Considering the conceptual design of the CSR-ETF
detector [10], two methods are invoked here to estimate
the impact parameter. The first one is the multiplicity
of forward neutrons with polar angle θ < 20o in the lab
frame which can be covered by a forward neutron wall.
The other method is to make use of the parameter Erat
[26], which is the ratio of the total transverse kinetic en-
ergy to the total longitudinal one. The particles are also
required to be within θ < 20o in the lab frame, while
the two qualities are calculated within the center of mass
0 0.2 0.4 0.6 0.8 1
tip-tip
body-body
0 0.2 0.4 0.6 0.8 10
maxb/b
FIG. 2: Upper: Forward neutron multiplicity and Lower:
Erat, as a function of normalized impact parameter b/bmax in
both tip-tip and body-body collisions.
system (c.m.s.).
Erat =
Ezi (1)
The normalized impact parameter b/bmax is used to rep-
resent centralities of tip-tip and body-body collisions and
the bmax of the two cases are quite different from each
other. As shown in Fig. 2, with either method, obvious
linear dependence of the normalized impact parameter
are demonstrated in both tip-tip and body-body near
central collisions. Then, the two methods can be com-
bined to determine the impact parameter to identify the
most central collision events in both tip-tip and body-
body collisions.
B. Stopping power ratio definition and evolution
It is difficult to obtain a universally accepted estimate
of the nuclear stopping power in heavy ion collisions due
to a proliferation of definitions of the concept [27]. The
stopping power ratio R [28] is employed to measure the
degree of stopping and defined as:
|Ptj |/
|Pzj | (2)
, the total nucleon transverse momentum |Ptj | divided by
the total absolute value of nucleon longitudinal momen-
tum |Pzj | in the c.m.s.. The ratio is wildly used to de-
scribe the degree of thermalization and nuclear stopping
by low and intermediate energies heavy ion collisions.
It’s a multi-particle observable on an event-by-event ba-
sis, which for an isotropic distribution is unity.
Fig. 3 shows the time and normalized impact param-
eter dependence of the stopping ratio R for three con-
ditions: tip-tip, body-body and sphere-sphere collisions.
0 10 20 30 40
=0max(a)b/b
tip-tip
body-body
sphere-sphere
0.2 0.4 0.6 0.8 1
(b)Minibias
tip-tip
body-body
sphere-sphere
t(fm/c) maxb/b
FIG. 3: (Color online) (a)The time evolution of the stopping
ratio R in tip-tip, body-body and sphere-sphere central colli-
sions, and (b) the stopping ratio R as a function of b/bmax in
minimum biased collisions.
When the ratio R reaches the value of 1, full stopping of
the collision system is considered to be achieved, and the
momenta is also isotropy, which are not sufficient but nec-
essary for thermal equilibrium of collision systems [28].
For R > 1, it can be explained by preponderance of mo-
mentum flow perpendicular to the beam direction [29]. It
is shown that all of the three conditions can achieve full
stopping when the stopping ratio R=1, the correspond-
ing time for body-body and tip-tip central collisions are
about 15 fm/c and 25 fm/c, respectively. Larger stopping
ratio and faster evolution to full stopping are observed
for body-body central collisions than tip-tip and sphere-
sphere ones at the early stage, which may indicate a more
violent colliding process for body-body central collisions
due to the sizable initial transverse overlap region. Al-
though the stopping ratio of tip-tip central collisions is
lowest than the other two cases at the early time, it raises
sharply later and even exceeding one. So, it means that
longer reaction and passage time can be obtained in tip-
tip central collisions than body-body and sphere-sphere
ones, which may indicate the nucleons in tip-tip colli-
sions can undergo more binary collisions to reach higher
transverse momentum.
In Fig. 3(b), the R of the three conditions are gradu-
ally decrease with the increase of the normalized impact
parameter. When b/bmax < 0.5, the ratio is always larger
for tip-tip collisions than the other two cases, while for
b/bmax > 0.5 all of the three conditions almost have the
same stopping power ratio.
III. BARYON, ENERGY DENSITY AND
THERMAL EQUILIBRIUM
Considering the discrepancy of stopping power be-
tween tip-tip and body-body collisions, it is interesting
4 (a)Baryon Density
tip-tip
body-body
Au-Au
0 10 20 30 40 50
(b)Energy Density
tip-tip
body-body
Au-Au
t(fm/c)
FIG. 4: The evolution of (a) baryon and (b) energy densities
in tip-tip, body-body and Au+Au central collisions.
to study further about the baryon and energy densities
evolution in both cases. As the full stopping and de-
formation effects in U+U collisions, it is believed higher
local baryon and energy densities system with long du-
ration can be created, which is considered to be a signif-
icant condition to study the nuclear EoS at high bayonic
density region.
A. The evolution of baryon and energy densities
The evolution of baryon and energy densities in the
central zone of tip-tip and body-body as well as Au+Au
central collisions are illustrated in Fig. 4.
In Fig. 4, it is observed the maximum attainable
baryon and energy densities for both tip-tip and body-
body central collisions are about 3.2 ρ0 and 0.8 GeV/fm
respectively, while the Au+Au one are about 2.6 ρ0 and
0.6 GeV/fm3. Both the baryon and energy densities in
U+U collisions are higher than the Au+Au one. Once a
baryon density threshold of ρ > 2.5 ρ0 is required, the
corresponding duration in tip-tip central collisions ∼ 20
fm/c (from ∼ 8 fm/c to ∼ 28 fm/c) is longer than ∼
10 fm/c ( from ∼ 8 fm/c to ∼ 18 fm/c ) of body-body
one, which is as predicted. But the peak densities have
no significant discrepancy between the two cases unlike
those at the energy region of the Alternating Gradient
Synchrotron (AGS) [12], which may be attribute to the
full stopping at the CSR energy.
B. Thermalization of the U+U collision systems
As mentioned before, the stopping ratio R = 1 is a
necessary but not sufficient condition for thermal equi-
librium of the collision system. In order to approach a
thermal equilibrium, a long duration of reaction is needed
for nucleons to undergo sufficient binary collisions. As
0 10 20 30 40 50
tip-tip
body-body
Au-Au
0 10 20 30 40 50
tip-tip
body-body
Au-Au
t(fm/c)
FIG. 5: The evolution of (a) volume with high density (ρ >
2.5ρ0) in tip-tip, body-body and Au+Au central collisions,
and (b) the scaled mean kinetic energy 2
< Ek >, within a
sphere of radius 2fm around the system mass center.
shown in Fig. 4(a), obvious long duration has been ob-
tained in both tip-tip and body-body central collisions.
It is therefore possible thermal equilibrium at the time
of freeze-out can be achieved.
The Fig. 5(a) is the evolution of the volume with the
high baryon density(ρ > 2.5 ρ0) for tip-tip, body-body
and Au+Au central collisions, respectively. Both tip-
tip and body-body central collisions have larger volumes
than Au+Au one at the same beam kinetic energy 520
MeV/nucleon. Although the maximum volume attain-
able for body-body central collisions(∼ 220 fm3) is about
two times larger than tip-tip one(∼ 120 fm3), the peak
volume of tip-tip central collisions lasts a much longer
time of ∼ 10 fm/c (from ∼ 15 fm/c to ∼ 25 fm/c) and
much more stable than body-body one. To estimate the
temperature at the freeze-out time, the scaled mean ki-
netic energy of all hadrons in a sphere of radius 2fm
around the system mass center is calculated as 2
< Ek >
[22], which is utilized to reflect the thermalization tem-
perature T of the collision system approximately. As
illustrated in Fig. 5 (b), both tip-tip and body-body
central collisions show a flat region about 75 MeV and
the corresponding time range are about 10 fm/c to 28
fm/c and 10 fm/c to 18 fm/c, respectively. Considering
the time range of the flat region in Fig. 5 (b) associat-
ing with the corresponding range in Fig. 5 (a) and also
looking back to Fig. 4, we obtain a large volume of hot,
dense nuclear matter in both tip-tip and body-body cen-
tral collisions. Consequently, the extreme circumstance
of sufficiently high temperature and density for a signif-
icant large volume and long duration [12, 22] has been
formed in tip-tip and body-body central collisions, which
can provide a good opportunity to study the nuclear EoS
as well as particles in medium properties, especially for
tip-tip case.
The time of freeze out should be cautiously determined
(a)tip-tip π+∆+*N
0 10 20 30 40 50 60
80 (b)Body-Body π+∆+*N
t(fm/c)
FIG. 6: Evolution of the multiplicity of the free pion, N∗ +
∆, N∗ +∆+ π in (a)tip-tip ,and (b)body-body central colli-
sions.
for estimating the thermalization temperature of colli-
sion system. In Fig. 6 the multiplicity evolution of free
pion which are not bounded in baryon resonances and
pion still bounded inside the excited baryon resonances
(∆, N∗) (unborn pion) are displayed. At the Lanzhou
CSR energy region (520 MeV/nucleon), the production
and destruction of the ∆ resonances are mainly through
NN ⇋ N∆ and ∆ → Nπ reactions in which the ∆ decay
rate is always higher than that of the formation of this
resonance and the production of pion is predominated by
the decay of the ∆ resonances (∆ → Nπ) [30]. The total
multiplicity of pion, ∆ and N∗ approaches a saturated
level after a period of evolution, indicating the freeze-
out time about t=28 fm/c and t=18 fm/c for tip-tip
and body-body central collisions, respectively. The larger
maximum attainable total multiplicity of pion, ∆ andN∗
and freeze out earlier indicates a existent of faster evo-
lution and more violently reaction process for the body-
body central collisions than tip-tip case consisting with
the discussing before.
The corresponding temperature about 75 MeV at
freeze-out time can be extracted from the Fig. 5 (b), for
both tip-tip and body-body central collisions. To further
confirm this estimation, both the energy spectrum of the
nucleon and negative-charged pion are studied within the
polar angle range of 900±100 in the c.m.s.. The thermo-
dynamic model [31] predicts that the energy spectra will
be represented by a temperature T which characterizes a
Maxwell-Boltzmann gas
PEdEdΩ
= const× e−Ekin/T (3)
, where P and E are the particle momentum and total
energy in the c.m.s.. Both the energy spectra and the
Boltzmann fit results are shown in Fig. 7. The inverse
slope (e.g. temperature T ) of the nucleons in tip-tip and
body-body central collisions are about 73 MeV and 70
0 0.2 0.4 0.6 0.8
310
tip-tip
body-body
(a)Nucleon
0 0.2 0.4 0.6 0.8
-π(b)
(GeV)kinE
FIG. 7: (a) Nucleon, and (b) negative-charged pion energy
spectrum at 900 ± 100 in the c.m.s. together with a Maxwell-
Boltzmann fit for both tip-tip and body-body central colli-
sions. The nucleon fit temperature for tip-tip and body-body
are about 73 MeV and 70 MeV, respectively and that of pion
are about 56 MeV and 52 MeV, respectively.
MeV, respectively, which are in good agreement with the
temperature extracted from the Fig. 5(b) at the freeze
out time. The spectra of negative-charged pion show a
different lower temperature than that of nucleon which
may be explained by considering an equilibrated N and
∆ system at thermal freeze out and taking into account
the kinematics of ∆ decay [32]. The nucleon temperature
closely reflects the freeze-out temperature of tip-tip and
body-body central collisions.
In conclusion, thermalization (or near thermalization)
of the collision system corresponding a freeze-out tem-
perature about 75 MeV is likely to be achieved in both
tip-tip and body-body central collisions. However, it’s
also possible that the collision system is still in a non-
equilibrium transport process on its path towards kinetic
equilibration [30].
IV. THE COLLECTIVE FLOW OF U+U
COLLISIONS
Stopping of nuclei in heavy ion collision can lead
to pressure gradient along different directions, result-
ing in collective motion as spectators bounce-off [34]
and participants squeeze-out effects [35]. Since last
two decades, at Bevalac/LBNL and SIS/GSI energies
the so-called ”collective flow” analysis has been estab-
lished [15, 34, 35, 36, 37] to study the collective mo-
tion of the products in heavy ion collisions. The collec-
tive flow resulting from bounce-off and squeeze-out ef-
fects, which can be explained well by the hydrodynamics
model [34, 38], and also be in good agreement with the
experimental data has been observed [39, 40]. Because of
the large deformation of the uranium nuclei, a novel col-
-1 -0.5 0 0.5 1
Soft:tip-tip
Soft:body-body
Cascade:tip-tip
Cascade:body-body
(a)b/b
-1 -0.5 0 0.5 1
(b)b/b
FIG. 8: The mean transverse momentum per nucleon pro-
jected into the reaction plane, < px/A >, as a function
of c.m.s. normalized rapidity is illustrated for tip-tip and
body-body collisions. With normalized impact parameter
cutoff:(a)b/bmax <= 0.5 (b)b/bmax > 0.5.
lective motion is expected [12], to be used to extract the
medium properties and nuclear matter EoS information.
[15, 16, 17, 18, 19, 20].
To perform flow analysis, it is necessary to construct a
imaginary reaction plane defined by direction of the beam
(z) and the impact parameter vector b [43, 45, 46]. In our
simulation, the x− z plane is just defined as the reaction
plane with the beam direction along z positive direction
and the impact parameter vector b along x positive direc-
tion. In last two decades, there are mainly two methods
to study the collective flow at the low and intermediate
energies. One is the sphericity method [28, 34, 41, 42]
which yields the flow angle relative to the beam axis of
the major axis of the best-fit kinetic energy ellipsoid, and
the other is to employ the mean transverse momentum
per nucleon projected into the reaction plane, < px/A >,
to perform nucleon sideward flow analysis [43, 44] which
reflects the spectator bounce-off effects in the reaction
plane. In recent years, it is usual to use an anisotropic
transverse flow analysis method. With a Fourier expan-
sion [47, 48] of the particle azimuthal angle φ distribu-
tion with respect to the reaction plane, different har-
monic coefficients can be extracted, among which the
first harmonic coefficient v1, called directed flow (simi-
lar to sideward flow) and the second harmonic coefficient
v2, called elliptic flow are mostly interested. The ellip-
tic flow reflects the anisotropy of emission particles in
the plane perpendicular to the reaction plane while the
directed flow describes the anisotropy in reaction plane.
The Fourier expansion can be expressed as
∼ 1 +
2vncos(nφ) (4)
Fig. 8 shows nucleon sideward flow, < px/A >, for
both tip-tip and body-body minimum biased collisions
(a)Nucleon Flow Parameter
tip-tip
body-body
Au-Au(500MeV/A)
0 0.2 0.4 0.6 0.8 1
10 2(b)Nucleon v
tip-tip
body-body
Au-Au(500MeV/A)
maxb/b
FIG. 9: (a)The nucleon flow parameter F and (b)the c.m.s.
mid-rapidity ( −0.5 < y0 < 0.5 ), nucleon elliptic flow v2 of
three collision conditions as a function of normalized impact
parameter b/bmax with soft EoS.
as a function of normalized rapidity, y(0) = Ycm/ycm, in
which Ycm represents the particle rapidity in c.m.s. and
ycm is the rapidity of the system mass center. To ex-
tract the nuclear EoS information and also demonstrate
the discrepancies of the nucleon sideward flow in tip-tip
and body-body collisions, the cascade events [49], which
neglect the mean field and pauli blocking effects are em-
ployed here to compare with the soft EoS case. In Fig.
8(a),(b), with a soft EoS, it is noted that either tip-tip
or body-body collisions show a spectator bounce-off ef-
fect revealing an obvious ”S” shape [15, 49] at the mid-
rapidity region of −0.5 < y0 < 0.5, while the cascade one
appear a almost vanishing nucleon sideward flow. It can
be understand by the nucleon sideward flow is related
to the mean field, which is mainly responsible for the
pressure gradient of the stopping nuclei, while the mean
field has a strong dependence of the nuclear EoS. There-
fore, the nucleon sideward flow is thought to be a good
indirect probe to extract the nuclear EoS information,
especially tip-tip case for its largely remarkable sideward
flow. A cutoff on normalized impact parameter is also
applied to explore the impact parameter dependence of
nucleon sideward flow. As shown in Fig. 8(b), when
b/bmax > 0.5 the curves of soft EoS and cascade are al-
most superposed with each other, while for b/bmax < 0.5
large discrepancy is observed. The situation is quite sim-
ilar to Fig. 3 (b), almost the same stopping power for
b/bmax > 0.5 and large discrepancy for b/bmax < 0.5 in
tip-tip and body-body minimum biased collisions, which
means there exists a correlation between nuclear stopping
power and sideward flow [33].
The normalized impact parameter dependence of the
collective flow of nucleon is further studied, by analyzing
the ”flow parameter” F [49] and also elliptic flow v2 for
both tip-tip and body-body as well as Au+Au minimum
biased collisions. The flow parameter F is a customarily
used quality to describe the nucleon sideward flow quan-
titatively defined as
d < px/A >
dy(0)
y(0)=0
the slope of the mean transverse momentum per nucleon
projected into the reaction plane at y(0) = 0.
In Fig. 9(a), with b/bmax > 0.5, the nucleon flow pa-
rameter F of tip-tip and body-body collisions are with
similar value. This similarity, along with the almost same
stopping ratioR in Fig. 3(b), indicates a existence of sim-
ilar pressure gradient effects on nucleon sideward flow in
the two collision orientations. While for b/bmax < 0.5,
the flow parameter F of tip-tip collisions is nearly 3 times
larger than that of body-body case. Even the sideward
flow of Au+Au collisions is larger than the body-body
one. It is further confirmed the tip-tip nucleon sideward
flow is a more sensitive probe to extract the information
of nuclear EoS than that of body-body one. The promi-
nence high of the nucleon sideward flow in tip-tip colli-
sions may be resulted from the stronger pressure gradient
between the participants and spectators in the reaction
plane than body-body one, due to the largely deformed
nuclei.
The normalized impact parameter dependent of nu-
cleon elliptic flow v2 at the mid-rapidity region ( −0.5 <
y0 < 0.5 ) is displayed in Fig. 9(b). A significant neg-
ative elliptic flow v2 at this energy region is consistent
with the excitation function of the elliptic flow studied
before [50]. An largest negative v2 about −12% in body-
body central collisions is observed which reflects the large
geometric and squeeze-out effects in the collisions. While
for tip-tip and Au+Au ones the maximum negative v2 are
obtained at mid-centrality. Since both high baryon, en-
ergy densities and large elliptic flow effects, which reflects
an early EoS of the hot dense compression nuclear matter
[17], are available in body-body central collisions. Thus
the body-body nucleon elliptic flow can also be taken as
a sensitive probe of nuclear EoS. The novel behaviors
of nucleon collective flow in tip-tip and body-body col-
lisions are mainly attributed to the large deformation of
the uranium nuclei.
V. SUMMARY
In summary, the CSR-ETF at Lanzhou provide a good
opportunity to systematically study the nuclear EoS at
the high net-baryon density region of nuclear matter
phase diagram. Due to the novel stopping effects in
largely deformed U+U collisions, the simulation based on
ART1.0 demonstrates that full stopping can be achieved
and also a bulk of hot, high densities nuclear matter with
large volume and long duration have been formed in both
tip-tip and body-body collisions. Large nucleon sideward
flow in tip-tip collisions and the significant negative nu-
cleon elliptic flow in body-body central collisions can pro-
vide a sensitive probe to extract nuclear EoS information.
Thus the extreme circumstance and the novel collective
flow in both tip-tip and body-body collisions can provide
a good condition and sensitive probe to study the nu-
clear EoS, respectively. More experimental observables
of U+U collision dynamics should be further studied, due
to the geometric effects.
VI. ACKNOWLEDGEMENT
This work is supported by National Natural Sci-
ence Foundation of China (10575101,10675111) and the
CAS/SAFEA International Partnership Program for
Creative Research Teams under the grant number of
CXTD-J2005-1. We wish to thank Bao-an Li, Feng Liu,
Qun Wang, Zhi-Gang Xiao and Nu Xu for their valuable
comments and suggestions.
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