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| # Community Pipeline Examples | |
| > **For more information about community pipelines, please have a look at [this issue](https://github.com/huggingface/diffusers/issues/841).** | |
| **Community pipeline** examples consist pipelines that have been added by the community. | |
| Please have a look at the following tables to get an overview of all community examples. Click on the **Code Example** to get a copy-and-paste ready code example that you can try out. | |
| If a community pipeline doesn't work as expected, please open an issue and ping the author on it. | |
| Please also check out our [Community Scripts](https://github.com/huggingface/diffusers/blob/main/examples/community/README_community_scripts.md) examples for tips and tricks that you can use with diffusers without having to run a community pipeline. | |
| | Example | Description | Code Example | Colab | Author | | |
| |:--------------------------------------------------------------------------------------------------------------------------------------|:---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------------|:-------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|--------------------------------------------------------------:| | |
| |Differential Diffusion|[Differential Diffusion](https://github.com/exx8/differential-diffusion) modifies an image according to a text prompt, and according to a map that specifies the amount of change in each region.|[Differential Diffusion](#differential-diffusion)|[](https://huggingface.co/spaces/exx8/differential-diffusion) [](https://colab.research.google.com/github/exx8/differential-diffusion/blob/main/examples/SD2.ipynb)|[Eran Levin](https://github.com/exx8) and [Ohad Fried](https://www.ohadf.com/)| | |
| | HD-Painter | [HD-Painter](https://github.com/Picsart-AI-Research/HD-Painter) enables prompt-faithfull and high resolution (up to 2k) image inpainting upon any diffusion-based image inpainting method. | [HD-Painter](#hd-painter) | [](https://huggingface.co/spaces/PAIR/HD-Painter) | [Manukyan Hayk](https://github.com/haikmanukyan) and [Sargsyan Andranik](https://github.com/AndranikSargsyan) | | |
| | Marigold Monocular Depth Estimation | A universal monocular depth estimator, utilizing Stable Diffusion, delivering sharp predictions in the wild. (See the [project page](https://marigoldmonodepth.github.io) and [full codebase](https://github.com/prs-eth/marigold) for more details.) | [Marigold Depth Estimation](#marigold-depth-estimation) | [](https://huggingface.co/spaces/toshas/marigold) [](https://colab.research.google.com/drive/12G8reD13DdpMie5ZQlaFNo2WCGeNUH-u?usp=sharing) | [Bingxin Ke](https://github.com/markkua) and [Anton Obukhov](https://github.com/toshas) | | |
| | LLM-grounded Diffusion (LMD+) | LMD greatly improves the prompt following ability of text-to-image generation models by introducing an LLM as a front-end prompt parser and layout planner. [Project page.](https://llm-grounded-diffusion.github.io/) [See our full codebase (also with diffusers).](https://github.com/TonyLianLong/LLM-groundedDiffusion) | [LLM-grounded Diffusion (LMD+)](#llm-grounded-diffusion) | [Huggingface Demo](https://huggingface.co/spaces/longlian/llm-grounded-diffusion) [](https://colab.research.google.com/drive/1SXzMSeAB-LJYISb2yrUOdypLz4OYWUKj) | [Long (Tony) Lian](https://tonylian.com/) | | |
| | CLIP Guided Stable Diffusion | Doing CLIP guidance for text to image generation with Stable Diffusion | [CLIP Guided Stable Diffusion](#clip-guided-stable-diffusion) | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/CLIP_Guided_Stable_diffusion_with_diffusers.ipynb) | [Suraj Patil](https://github.com/patil-suraj/) | | |
| | One Step U-Net (Dummy) | Example showcasing of how to use Community Pipelines (see <https://github.com/huggingface/diffusers/issues/841>) | [One Step U-Net](#one-step-unet) | - | [Patrick von Platen](https://github.com/patrickvonplaten/) | | |
| | Stable Diffusion Interpolation | Interpolate the latent space of Stable Diffusion between different prompts/seeds | [Stable Diffusion Interpolation](#stable-diffusion-interpolation) | - | [Nate Raw](https://github.com/nateraw/) | | |
| | Stable Diffusion Mega | **One** Stable Diffusion Pipeline with all functionalities of [Text2Image](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py), [Image2Image](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py) and [Inpainting](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_inpaint.py) | [Stable Diffusion Mega](#stable-diffusion-mega) | - | [Patrick von Platen](https://github.com/patrickvonplaten/) | | |
| | Long Prompt Weighting Stable Diffusion | **One** Stable Diffusion Pipeline without tokens length limit, and support parsing weighting in prompt. | [Long Prompt Weighting Stable Diffusion](#long-prompt-weighting-stable-diffusion) | - | [SkyTNT](https://github.com/SkyTNT) | | |
| | Speech to Image | Using automatic-speech-recognition to transcribe text and Stable Diffusion to generate images | [Speech to Image](#speech-to-image) | - | [Mikail Duzenli](https://github.com/MikailINTech) | |
| | Wild Card Stable Diffusion | Stable Diffusion Pipeline that supports prompts that contain wildcard terms (indicated by surrounding double underscores), with values instantiated randomly from a corresponding txt file or a dictionary of possible values | [Wildcard Stable Diffusion](#wildcard-stable-diffusion) | - | [Shyam Sudhakaran](https://github.com/shyamsn97) | | |
| | [Composable Stable Diffusion](https://energy-based-model.github.io/Compositional-Visual-Generation-with-Composable-Diffusion-Models/) | Stable Diffusion Pipeline that supports prompts that contain "|" in prompts (as an AND condition) and weights (separated by "|" as well) to positively / negatively weight prompts. | [Composable Stable Diffusion](#composable-stable-diffusion) | - | [Mark Rich](https://github.com/MarkRich) | | |
| | Seed Resizing Stable Diffusion | Stable Diffusion Pipeline that supports resizing an image and retaining the concepts of the 512 by 512 generation. | [Seed Resizing](#seed-resizing) | - | [Mark Rich](https://github.com/MarkRich) | | |
| | Imagic Stable Diffusion | Stable Diffusion Pipeline that enables writing a text prompt to edit an existing image | [Imagic Stable Diffusion](#imagic-stable-diffusion) | - | [Mark Rich](https://github.com/MarkRich) | | |
| | Multilingual Stable Diffusion | Stable Diffusion Pipeline that supports prompts in 50 different languages. | [Multilingual Stable Diffusion](#multilingual-stable-diffusion-pipeline) | - | [Juan Carlos Piñeros](https://github.com/juancopi81) | | |
| | GlueGen Stable Diffusion | Stable Diffusion Pipeline that supports prompts in different languages using GlueGen adapter. | [GlueGen Stable Diffusion](#gluegen-stable-diffusion-pipeline) | - | [Phạm Hồng Vinh](https://github.com/rootonchair) | | |
| | Image to Image Inpainting Stable Diffusion | Stable Diffusion Pipeline that enables the overlaying of two images and subsequent inpainting | [Image to Image Inpainting Stable Diffusion](#image-to-image-inpainting-stable-diffusion) | - | [Alex McKinney](https://github.com/vvvm23) | | |
| | Text Based Inpainting Stable Diffusion | Stable Diffusion Inpainting Pipeline that enables passing a text prompt to generate the mask for inpainting | [Text Based Inpainting Stable Diffusion](#image-to-image-inpainting-stable-diffusion) | - | [Dhruv Karan](https://github.com/unography) | | |
| | Bit Diffusion | Diffusion on discrete data | [Bit Diffusion](#bit-diffusion) | - | [Stuti R.](https://github.com/kingstut) | | |
| | K-Diffusion Stable Diffusion | Run Stable Diffusion with any of [K-Diffusion's samplers](https://github.com/crowsonkb/k-diffusion/blob/master/k_diffusion/sampling.py) | [Stable Diffusion with K Diffusion](#stable-diffusion-with-k-diffusion) | - | [Patrick von Platen](https://github.com/patrickvonplaten/) | | |
| | Checkpoint Merger Pipeline | Diffusion Pipeline that enables merging of saved model checkpoints | [Checkpoint Merger Pipeline](#checkpoint-merger-pipeline) | - | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) | | |
| | Stable Diffusion v1.1-1.4 Comparison | Run all 4 model checkpoints for Stable Diffusion and compare their results together | [Stable Diffusion Comparison](#stable-diffusion-comparisons) | - | [Suvaditya Mukherjee](https://github.com/suvadityamuk) | | |
| | MagicMix | Diffusion Pipeline for semantic mixing of an image and a text prompt | [MagicMix](#magic-mix) | - | [Partho Das](https://github.com/daspartho) | | |
| | Stable UnCLIP | Diffusion Pipeline for combining prior model (generate clip image embedding from text, UnCLIPPipeline `"kakaobrain/karlo-v1-alpha"`) and decoder pipeline (decode clip image embedding to image, StableDiffusionImageVariationPipeline `"lambdalabs/sd-image-variations-diffusers"` ). | [Stable UnCLIP](#stable-unclip) | - | [Ray Wang](https://wrong.wang) | | |
| | UnCLIP Text Interpolation Pipeline | Diffusion Pipeline that allows passing two prompts and produces images while interpolating between the text-embeddings of the two prompts | [UnCLIP Text Interpolation Pipeline](#unclip-text-interpolation-pipeline) | - | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) | | |
| | UnCLIP Image Interpolation Pipeline | Diffusion Pipeline that allows passing two images/image_embeddings and produces images while interpolating between their image-embeddings | [UnCLIP Image Interpolation Pipeline](#unclip-image-interpolation-pipeline) | - | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) | | |
| | DDIM Noise Comparative Analysis Pipeline | Investigating how the diffusion models learn visual concepts from each noise level (which is a contribution of [P2 weighting (CVPR 2022)](https://arxiv.org/abs/2204.00227)) | [DDIM Noise Comparative Analysis Pipeline](#ddim-noise-comparative-analysis-pipeline) | - | [Aengus (Duc-Anh)](https://github.com/aengusng8) | | |
| | CLIP Guided Img2Img Stable Diffusion Pipeline | Doing CLIP guidance for image to image generation with Stable Diffusion | [CLIP Guided Img2Img Stable Diffusion](#clip-guided-img2img-stable-diffusion) | - | [Nipun Jindal](https://github.com/nipunjindal/) | | |
| | TensorRT Stable Diffusion Text to Image Pipeline | Accelerates the Stable Diffusion Text2Image Pipeline using TensorRT | [TensorRT Stable Diffusion Text to Image Pipeline](#tensorrt-text2image-stable-diffusion-pipeline) | - | [Asfiya Baig](https://github.com/asfiyab-nvidia) | | |
| | EDICT Image Editing Pipeline | Diffusion pipeline for text-guided image editing | [EDICT Image Editing Pipeline](#edict-image-editing-pipeline) | - | [Joqsan Azocar](https://github.com/Joqsan) | | |
| | Stable Diffusion RePaint | Stable Diffusion pipeline using [RePaint](https://arxiv.org/abs/2201.0986) for inpainting. | [Stable Diffusion RePaint](#stable-diffusion-repaint ) | - | [Markus Pobitzer](https://github.com/Markus-Pobitzer) | | |
| | TensorRT Stable Diffusion Image to Image Pipeline | Accelerates the Stable Diffusion Image2Image Pipeline using TensorRT | [TensorRT Stable Diffusion Image to Image Pipeline](#tensorrt-image2image-stable-diffusion-pipeline) | - | [Asfiya Baig](https://github.com/asfiyab-nvidia) | | |
| | Stable Diffusion IPEX Pipeline | Accelerate Stable Diffusion inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [Stable Diffusion on IPEX](#stable-diffusion-on-ipex) | - | [Yingjie Han](https://github.com/yingjie-han/) | | |
| | CLIP Guided Images Mixing Stable Diffusion Pipeline | Сombine images using usual diffusion models. | [CLIP Guided Images Mixing Using Stable Diffusion](#clip-guided-images-mixing-with-stable-diffusion) | - | [Karachev Denis](https://github.com/TheDenk) | | |
| | TensorRT Stable Diffusion Inpainting Pipeline | Accelerates the Stable Diffusion Inpainting Pipeline using TensorRT | [TensorRT Stable Diffusion Inpainting Pipeline](#tensorrt-inpainting-stable-diffusion-pipeline) | - | [Asfiya Baig](https://github.com/asfiyab-nvidia) | | |
| | IADB Pipeline | Implementation of [Iterative α-(de)Blending: a Minimalist Deterministic Diffusion Model](https://arxiv.org/abs/2305.03486) | [IADB Pipeline](#iadb-pipeline) | - | [Thomas Chambon](https://github.com/tchambon) | |
| | Zero1to3 Pipeline | Implementation of [Zero-1-to-3: Zero-shot One Image to 3D Object](https://arxiv.org/abs/2303.11328) | [Zero1to3 Pipeline](#zero1to3-pipeline) | - | [Xin Kong](https://github.com/kxhit) | | |
| | Stable Diffusion XL Long Weighted Prompt Pipeline | A pipeline support unlimited length of prompt and negative prompt, use A1111 style of prompt weighting | [Stable Diffusion XL Long Weighted Prompt Pipeline](#stable-diffusion-xl-long-weighted-prompt-pipeline) | [](https://colab.research.google.com/drive/1LsqilswLR40XLLcp6XFOl5nKb_wOe26W?usp=sharing) | [Andrew Zhu](https://xhinker.medium.com/) | | |
| | FABRIC - Stable Diffusion with feedback Pipeline | pipeline supports feedback from liked and disliked images | [Stable Diffusion Fabric Pipeline](#stable-diffusion-fabric-pipeline) | - | [Shauray Singh](https://shauray8.github.io/about_shauray/) | | |
| | sketch inpaint - Inpainting with non-inpaint Stable Diffusion | sketch inpaint much like in automatic1111 | [Masked Im2Im Stable Diffusion Pipeline](#stable-diffusion-masked-im2im) | - | [Anatoly Belikov](https://github.com/noskill) | | |
| | prompt-to-prompt | change parts of a prompt and retain image structure (see [paper page](https://prompt-to-prompt.github.io/)) | [Prompt2Prompt Pipeline](#prompt2prompt-pipeline) | - | [Umer H. Adil](https://twitter.com/UmerHAdil) | | |
| | Latent Consistency Pipeline | Implementation of [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://arxiv.org/abs/2310.04378) | [Latent Consistency Pipeline](#latent-consistency-pipeline) | - | [Simian Luo](https://github.com/luosiallen) | | |
| | Latent Consistency Img2img Pipeline | Img2img pipeline for Latent Consistency Models | [Latent Consistency Img2Img Pipeline](#latent-consistency-img2img-pipeline) | - | [Logan Zoellner](https://github.com/nagolinc) | | |
| | Latent Consistency Interpolation Pipeline | Interpolate the latent space of Latent Consistency Models with multiple prompts | [Latent Consistency Interpolation Pipeline](#latent-consistency-interpolation-pipeline) | [](https://colab.research.google.com/drive/1pK3NrLWJSiJsBynLns1K1-IDTW9zbPvl?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) | | |
| | SDE Drag Pipeline | The pipeline supports drag editing of images using stochastic differential equations | [SDE Drag Pipeline](#sde-drag-pipeline) | - | [NieShen](https://github.com/NieShenRuc) [Fengqi Zhu](https://github.com/Monohydroxides) | | |
| | Regional Prompting Pipeline | Assign multiple prompts for different regions | [Regional Prompting Pipeline](#regional-prompting-pipeline) | - | [hako-mikan](https://github.com/hako-mikan) | | |
| | LDM3D-sr (LDM3D upscaler) | Upscale low resolution RGB and depth inputs to high resolution | [StableDiffusionUpscaleLDM3D Pipeline](https://github.com/estelleafl/diffusers/tree/ldm3d_upscaler_community/examples/community#stablediffusionupscaleldm3d-pipeline) | - | [Estelle Aflalo](https://github.com/estelleafl) | | |
| | AnimateDiff ControlNet Pipeline | Combines AnimateDiff with precise motion control using ControlNets | [AnimateDiff ControlNet Pipeline](#animatediff-controlnet-pipeline) | [](https://colab.research.google.com/drive/1SKboYeGjEQmQPWoFC0aLYpBlYdHXkvAu?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) and [Edoardo Botta](https://github.com/EdoardoBotta) | | |
| | DemoFusion Pipeline | Implementation of [DemoFusion: Democratising High-Resolution Image Generation With No $$$](https://arxiv.org/abs/2311.16973) | [DemoFusion Pipeline](#demofusion) | - | [Ruoyi Du](https://github.com/RuoyiDu) | | |
| | Instaflow Pipeline | Implementation of [InstaFlow! One-Step Stable Diffusion with Rectified Flow](https://arxiv.org/abs/2309.06380) | [Instaflow Pipeline](#instaflow-pipeline) | - | [Ayush Mangal](https://github.com/ayushtues) | | |
| | Null-Text Inversion Pipeline | Implement [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://arxiv.org/abs/2211.09794) as a pipeline. | [Null-Text Inversion](https://github.com/google/prompt-to-prompt/) | - | [Junsheng Luan](https://github.com/Junsheng121) | | |
| | Rerender A Video Pipeline | Implementation of [[SIGGRAPH Asia 2023] Rerender A Video: Zero-Shot Text-Guided Video-to-Video Translation](https://arxiv.org/abs/2306.07954) | [Rerender A Video Pipeline](#rerender-a-video) | - | [Yifan Zhou](https://github.com/SingleZombie) | | |
| | StyleAligned Pipeline | Implementation of [Style Aligned Image Generation via Shared Attention](https://arxiv.org/abs/2312.02133) | [StyleAligned Pipeline](#stylealigned-pipeline) | [](https://drive.google.com/file/d/15X2E0jFPTajUIjS0FzX50OaHsCbP2lQ0/view?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) | | |
| | AnimateDiff Image-To-Video Pipeline | Experimental Image-To-Video support for AnimateDiff (open to improvements) | [AnimateDiff Image To Video Pipeline](#animatediff-image-to-video-pipeline) | [](https://drive.google.com/file/d/1TvzCDPHhfFtdcJZe4RLloAwyoLKuttWK/view?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) | | |
| | IP Adapter FaceID Stable Diffusion | Stable Diffusion Pipeline that supports IP Adapter Face ID | [IP Adapter Face ID](#ip-adapter-face-id) | - | [Fabio Rigano](https://github.com/fabiorigano) | | |
| | InstantID Pipeline | Stable Diffusion XL Pipeline that supports InstantID | [InstantID Pipeline](#instantid-pipeline) | [](https://huggingface.co/spaces/InstantX/InstantID) | [Haofan Wang](https://github.com/haofanwang) | | |
| | UFOGen Scheduler | Scheduler for UFOGen Model (compatible with Stable Diffusion pipelines) | [UFOGen Scheduler](#ufogen-scheduler) | - | [dg845](https://github.com/dg845) | | |
| | Stable Diffusion XL IPEX Pipeline | Accelerate Stable Diffusion XL inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [Stable Diffusion XL on IPEX](#stable-diffusion-xl-on-ipex) | - | [Dan Li](https://github.com/ustcuna/) | | |
| | Stable Diffusion BoxDiff Pipeline | Training-free controlled generation with bounding boxes using [BoxDiff](https://github.com/showlab/BoxDiff) | [Stable Diffusion BoxDiff Pipeline](#stable-diffusion-boxdiff) | - | [Jingyang Zhang](https://github.com/zjysteven/) | | |
| | FRESCO V2V Pipeline | Implementation of [[CVPR 2024] FRESCO: Spatial-Temporal Correspondence for Zero-Shot Video Translation](https://arxiv.org/abs/2403.12962) | [FRESCO V2V Pipeline](#fresco) | - | [Yifan Zhou](https://github.com/SingleZombie) | | |
| To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly. | |
| ```py | |
| pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", custom_pipeline="filename_in_the_community_folder") | |
| ``` | |
| ## Example usages | |
| ### Differential Diffusion | |
| **Eran Levin, Ohad Fried** | |
| **Tel Aviv University, Reichman University** | |
| Diffusion models have revolutionized image generation and editing, producing state-of-the-art results in conditioned and unconditioned image synthesis. While current techniques enable user control over the degree of change in an image edit, the controllability is limited to global changes over an entire edited region. This paper introduces a novel framework that enables customization of the amount of change per pixel or per image region. Our framework can be integrated into any existing diffusion model, enhancing it with this capability. Such granular control on the quantity of change opens up a diverse array of new editing capabilities, such as control of the extent to which individual objects are modified, or the ability to introduce gradual spatial changes. Furthermore, we showcase the framework's effectiveness in soft-inpainting---the completion of portions of an image while subtly adjusting the surrounding areas to ensure seamless integration. Additionally, we introduce a new tool for exploring the effects of different change quantities. Our framework operates solely during inference, requiring no model training or fine-tuning. We demonstrate our method with the current open state-of-the-art models, and validate it via both quantitative and qualitative comparisons, and a user study. | |
|  | |
| You can find additional information about Differential Diffusion in the [paper](https://differential-diffusion.github.io/paper.pdf) or in the [project website](https://differential-diffusion.github.io/). | |
| #### Usage example | |
| ```python | |
| import torch | |
| from torchvision import transforms | |
| from diffusers import DPMSolverMultistepScheduler | |
| from diffusers.utils import load_image | |
| from examples.community.pipeline_stable_diffusion_xl_differential_img2img import ( | |
| StableDiffusionXLDifferentialImg2ImgPipeline, | |
| ) | |
| pipeline = StableDiffusionXLDifferentialImg2ImgPipeline.from_pretrained( | |
| "SG161222/RealVisXL_V4.0", torch_dtype=torch.float16, variant="fp16" | |
| ).to("cuda") | |
| pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config, use_karras_sigmas=True) | |
| def preprocess_image(image): | |
| image = image.convert("RGB") | |
| image = transforms.CenterCrop((image.size[1] // 64 * 64, image.size[0] // 64 * 64))(image) | |
| image = transforms.ToTensor()(image) | |
| image = image * 2 - 1 | |
| image = image.unsqueeze(0).to("cuda") | |
| return image | |
| def preprocess_map(map): | |
| map = map.convert("L") | |
| map = transforms.CenterCrop((map.size[1] // 64 * 64, map.size[0] // 64 * 64))(map) | |
| map = transforms.ToTensor()(map) | |
| map = map.to("cuda") | |
| return map | |
| image = preprocess_image( | |
| load_image( | |
| "https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/20240329211129_4024911930.png?download=true" | |
| ) | |
| ) | |
| mask = preprocess_map( | |
| load_image( | |
| "https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/gradient_mask.png?download=true" | |
| ) | |
| ) | |
| prompt = "a green pear" | |
| negative_prompt = "blurry" | |
| image = pipeline( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| guidance_scale=7.5, | |
| num_inference_steps=25, | |
| original_image=image, | |
| image=image, | |
| strength=1.0, | |
| map=mask, | |
| ).images[0] | |
| image.save("result.png") | |
| ``` | |
| ### HD-Painter | |
| Implementation of [HD-Painter: High-Resolution and Prompt-Faithful Text-Guided Image Inpainting with Diffusion Models](https://arxiv.org/abs/2312.14091). | |
|  | |
| The abstract from the paper is: | |
| Recent progress in text-guided image inpainting, based on the unprecedented success of text-to-image diffusion models, has led to exceptionally realistic and visually plausible results. | |
| However, there is still significant potential for improvement in current text-to-image inpainting models, particularly in better aligning the inpainted area with user prompts and performing high-resolution inpainting. | |
| Therefore, in this paper we introduce _HD-Painter_, a completely **training-free** approach that **accurately follows to prompts** and coherently **scales to high-resolution** image inpainting. | |
| To this end, we design the _Prompt-Aware Introverted Attention (PAIntA)_ layer enhancing self-attention scores by prompt information and resulting in better text alignment generations. | |
| To further improve the prompt coherence we introduce the _Reweighting Attention Score Guidance (RASG)_ mechanism seamlessly integrating a post-hoc sampling strategy into general form of DDIM to prevent out-of-distribution latent shifts. | |
| Moreover, HD-Painter allows extension to larger scales by introducing a specialized super-resolution technique customized for inpainting, enabling the completion of missing regions in images of up to 2K resolution. | |
| Our experiments demonstrate that HD-Painter surpasses existing state-of-the-art approaches qualitatively and quantitatively, achieving an impressive generation accuracy improvement of **61.4** vs **51.9**. | |
| We will make the codes publicly available. | |
| You can find additional information about Text2Video-Zero in the [paper](https://arxiv.org/abs/2312.14091) or the [original codebase](https://github.com/Picsart-AI-Research/HD-Painter). | |
| #### Usage example | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline, DDIMScheduler | |
| from diffusers.utils import load_image, make_image_grid | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-2-inpainting", | |
| custom_pipeline="hd_painter" | |
| ) | |
| pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config) | |
| prompt = "wooden boat" | |
| init_image = load_image("https://raw.githubusercontent.com/Picsart-AI-Research/HD-Painter/main/__assets__/samples/images/2.jpg") | |
| mask_image = load_image("https://raw.githubusercontent.com/Picsart-AI-Research/HD-Painter/main/__assets__/samples/masks/2.png") | |
| image = pipe (prompt, init_image, mask_image, use_rasg = True, use_painta = True, generator=torch.manual_seed(12345)).images[0] | |
| make_image_grid([init_image, mask_image, image], rows=1, cols=3) | |
| ``` | |
| ### Marigold Depth Estimation | |
| Marigold is a universal monocular depth estimator that delivers accurate and sharp predictions in the wild. Based on Stable Diffusion, it is trained exclusively with synthetic depth data and excels in zero-shot adaptation to real-world imagery. This pipeline is an official implementation of the inference process. More details can be found on our [project page](https://marigoldmonodepth.github.io) and [full codebase](https://github.com/prs-eth/marigold) (also implemented with diffusers). | |
|  | |
| This depth estimation pipeline processes a single input image through multiple diffusion denoising stages to estimate depth maps. These maps are subsequently merged to produce the final output. Below is an example code snippet, including optional arguments: | |
| ```python | |
| import numpy as np | |
| import torch | |
| from PIL import Image | |
| from diffusers import DiffusionPipeline | |
| from diffusers.utils import load_image | |
| # Original DDIM version (higher quality) | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "prs-eth/marigold-v1-0", | |
| custom_pipeline="marigold_depth_estimation" | |
| # torch_dtype=torch.float16, # (optional) Run with half-precision (16-bit float). | |
| # variant="fp16", # (optional) Use with `torch_dtype=torch.float16`, to directly load fp16 checkpoint | |
| ) | |
| # (New) LCM version (faster speed) | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "prs-eth/marigold-lcm-v1-0", | |
| custom_pipeline="marigold_depth_estimation" | |
| # torch_dtype=torch.float16, # (optional) Run with half-precision (16-bit float). | |
| # variant="fp16", # (optional) Use with `torch_dtype=torch.float16`, to directly load fp16 checkpoint | |
| ) | |
| pipe.to("cuda") | |
| img_path_or_url = "https://share.phys.ethz.ch/~pf/bingkedata/marigold/pipeline_example.jpg" | |
| image: Image.Image = load_image(img_path_or_url) | |
| pipeline_output = pipe( | |
| image, # Input image. | |
| # ----- recommended setting for DDIM version ----- | |
| # denoising_steps=10, # (optional) Number of denoising steps of each inference pass. Default: 10. | |
| # ensemble_size=10, # (optional) Number of inference passes in the ensemble. Default: 10. | |
| # ------------------------------------------------ | |
| # ----- recommended setting for LCM version ------ | |
| # denoising_steps=4, | |
| # ensemble_size=5, | |
| # ------------------------------------------------- | |
| # processing_res=768, # (optional) Maximum resolution of processing. If set to 0: will not resize at all. Defaults to 768. | |
| # match_input_res=True, # (optional) Resize depth prediction to match input resolution. | |
| # batch_size=0, # (optional) Inference batch size, no bigger than `num_ensemble`. If set to 0, the script will automatically decide the proper batch size. Defaults to 0. | |
| # seed=2024, # (optional) Random seed can be set to ensure additional reproducibility. Default: None (unseeded). Note: forcing --batch_size 1 helps to increase reproducibility. To ensure full reproducibility, deterministic mode needs to be used. | |
| # color_map="Spectral", # (optional) Colormap used to colorize the depth map. Defaults to "Spectral". Set to `None` to skip colormap generation. | |
| # show_progress_bar=True, # (optional) If true, will show progress bars of the inference progress. | |
| ) | |
| depth: np.ndarray = pipeline_output.depth_np # Predicted depth map | |
| depth_colored: Image.Image = pipeline_output.depth_colored # Colorized prediction | |
| # Save as uint16 PNG | |
| depth_uint16 = (depth * 65535.0).astype(np.uint16) | |
| Image.fromarray(depth_uint16).save("./depth_map.png", mode="I;16") | |
| # Save colorized depth map | |
| depth_colored.save("./depth_colored.png") | |
| ``` | |
| ### LLM-grounded Diffusion | |
| LMD and LMD+ greatly improves the prompt understanding ability of text-to-image generation models by introducing an LLM as a front-end prompt parser and layout planner. It improves spatial reasoning, the understanding of negation, attribute binding, generative numeracy, etc. in a unified manner without explicitly aiming for each. LMD is completely training-free (i.e., uses SD model off-the-shelf). LMD+ takes in additional adapters for better control. This is a reproduction of LMD+ model used in our work. [Project page.](https://llm-grounded-diffusion.github.io/) [See our full codebase (also with diffusers).](https://github.com/TonyLianLong/LLM-groundedDiffusion) | |
|  | |
|  | |
| This pipeline can be used with an LLM or on its own. We provide a parser that parses LLM outputs to the layouts. You can obtain the prompt to input to the LLM for layout generation [here](https://github.com/TonyLianLong/LLM-groundedDiffusion/blob/main/prompt.py). After feeding the prompt to an LLM (e.g., GPT-4 on ChatGPT website), you can feed the LLM response into our pipeline. | |
| The following code has been tested on 1x RTX 4090, but it should also support GPUs with lower GPU memory. | |
| #### Use this pipeline with an LLM | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "longlian/lmd_plus", | |
| custom_pipeline="llm_grounded_diffusion", | |
| custom_revision="main", | |
| variant="fp16", torch_dtype=torch.float16 | |
| ) | |
| pipe.enable_model_cpu_offload() | |
| # Generate directly from a text prompt and an LLM response | |
| prompt = "a waterfall and a modern high speed train in a beautiful forest with fall foliage" | |
| phrases, boxes, bg_prompt, neg_prompt = pipe.parse_llm_response(""" | |
| [('a waterfall', [71, 105, 148, 258]), ('a modern high speed train', [255, 223, 181, 149])] | |
| Background prompt: A beautiful forest with fall foliage | |
| Negative prompt: | |
| """) | |
| images = pipe( | |
| prompt=prompt, | |
| negative_prompt=neg_prompt, | |
| phrases=phrases, | |
| boxes=boxes, | |
| gligen_scheduled_sampling_beta=0.4, | |
| output_type="pil", | |
| num_inference_steps=50, | |
| lmd_guidance_kwargs={} | |
| ).images | |
| images[0].save("./lmd_plus_generation.jpg") | |
| ``` | |
| #### Use this pipeline on its own for layout generation | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "longlian/lmd_plus", | |
| custom_pipeline="llm_grounded_diffusion", | |
| variant="fp16", torch_dtype=torch.float16 | |
| ) | |
| pipe.enable_model_cpu_offload() | |
| # Generate an image described by the prompt and | |
| # insert objects described by text at the region defined by bounding boxes | |
| prompt = "a waterfall and a modern high speed train in a beautiful forest with fall foliage" | |
| boxes = [[0.1387, 0.2051, 0.4277, 0.7090], [0.4980, 0.4355, 0.8516, 0.7266]] | |
| phrases = ["a waterfall", "a modern high speed train"] | |
| images = pipe( | |
| prompt=prompt, | |
| phrases=phrases, | |
| boxes=boxes, | |
| gligen_scheduled_sampling_beta=0.4, | |
| output_type="pil", | |
| num_inference_steps=50, | |
| lmd_guidance_kwargs={} | |
| ).images | |
| images[0].save("./lmd_plus_generation.jpg") | |
| ``` | |
| ### CLIP Guided Stable Diffusion | |
| CLIP guided stable diffusion can help to generate more realistic images | |
| by guiding stable diffusion at every denoising step with an additional CLIP model. | |
| The following code requires roughly 12GB of GPU RAM. | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| from transformers import CLIPImageProcessor, CLIPModel | |
| import torch | |
| feature_extractor = CLIPImageProcessor.from_pretrained("laion/CLIP-ViT-B-32-laion2B-s34B-b79K") | |
| clip_model = CLIPModel.from_pretrained("laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16) | |
| guided_pipeline = DiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", | |
| custom_pipeline="clip_guided_stable_diffusion", | |
| clip_model=clip_model, | |
| feature_extractor=feature_extractor, | |
| torch_dtype=torch.float16, | |
| ) | |
| guided_pipeline.enable_attention_slicing() | |
| guided_pipeline = guided_pipeline.to("cuda") | |
| prompt = "fantasy book cover, full moon, fantasy forest landscape, golden vector elements, fantasy magic, dark light night, intricate, elegant, sharp focus, illustration, highly detailed, digital painting, concept art, matte, art by WLOP and Artgerm and Albert Bierstadt, masterpiece" | |
| generator = torch.Generator(device="cuda").manual_seed(0) | |
| images = [] | |
| for i in range(4): | |
| image = guided_pipeline( | |
| prompt, | |
| num_inference_steps=50, | |
| guidance_scale=7.5, | |
| clip_guidance_scale=100, | |
| num_cutouts=4, | |
| use_cutouts=False, | |
| generator=generator, | |
| ).images[0] | |
| images.append(image) | |
| # save images locally | |
| for i, img in enumerate(images): | |
| img.save(f"./clip_guided_sd/image_{i}.png") | |
| ``` | |
| The `images` list contains a list of PIL images that can be saved locally or displayed directly in a google colab. | |
| Generated images tend to be of higher qualtiy than natively using stable diffusion. E.g. the above script generates the following images: | |
| . | |
| ### One Step Unet | |
| The dummy "one-step-unet" can be run as follows: | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained("google/ddpm-cifar10-32", custom_pipeline="one_step_unet") | |
| pipe() | |
| ``` | |
| **Note**: This community pipeline is not useful as a feature, but rather just serves as an example of how community pipelines can be added (see <https://github.com/huggingface/diffusers/issues/841>). | |
| ### Stable Diffusion Interpolation | |
| The following code can be run on a GPU of at least 8GB VRAM and should take approximately 5 minutes. | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| revision='fp16', | |
| torch_dtype=torch.float16, | |
| safety_checker=None, # Very important for videos...lots of false positives while interpolating | |
| custom_pipeline="interpolate_stable_diffusion", | |
| ).to('cuda') | |
| pipe.enable_attention_slicing() | |
| frame_filepaths = pipe.walk( | |
| prompts=['a dog', 'a cat', 'a horse'], | |
| seeds=[42, 1337, 1234], | |
| num_interpolation_steps=16, | |
| output_dir='./dreams', | |
| batch_size=4, | |
| height=512, | |
| width=512, | |
| guidance_scale=8.5, | |
| num_inference_steps=50, | |
| ) | |
| ``` | |
| The output of the `walk(...)` function returns a list of images saved under the folder as defined in `output_dir`. You can use these images to create videos of stable diffusion. | |
| > **Please have a look at <https://github.com/nateraw/stable-diffusion-videos> for more in-detail information on how to create videos using stable diffusion as well as more feature-complete functionality.** | |
| ### Stable Diffusion Mega | |
| The Stable Diffusion Mega Pipeline lets you use the main use cases of the stable diffusion pipeline in a single class. | |
| ```python | |
| #!/usr/bin/env python3 | |
| from diffusers import DiffusionPipeline | |
| import PIL | |
| import requests | |
| from io import BytesIO | |
| import torch | |
| def download_image(url): | |
| response = requests.get(url) | |
| return PIL.Image.open(BytesIO(response.content)).convert("RGB") | |
| pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_mega", torch_dtype=torch.float16, revision="fp16") | |
| pipe.to("cuda") | |
| pipe.enable_attention_slicing() | |
| ### Text-to-Image | |
| images = pipe.text2img("An astronaut riding a horse").images | |
| ### Image-to-Image | |
| init_image = download_image("https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg") | |
| prompt = "A fantasy landscape, trending on artstation" | |
| images = pipe.img2img(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images | |
| ### Inpainting | |
| img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
| mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
| init_image = download_image(img_url).resize((512, 512)) | |
| mask_image = download_image(mask_url).resize((512, 512)) | |
| prompt = "a cat sitting on a bench" | |
| images = pipe.inpaint(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.75).images | |
| ``` | |
| As shown above this one pipeline can run all both "text-to-image", "image-to-image", and "inpainting" in one pipeline. | |
| ### Long Prompt Weighting Stable Diffusion | |
| Features of this custom pipeline: | |
| - Input a prompt without the 77 token length limit. | |
| - Includes tx2img, img2img. and inpainting pipelines. | |
| - Emphasize/weigh part of your prompt with parentheses as so: `a baby deer with (big eyes)` | |
| - De-emphasize part of your prompt as so: `a [baby] deer with big eyes` | |
| - Precisely weigh part of your prompt as so: `a baby deer with (big eyes:1.3)` | |
| Prompt weighting equivalents: | |
| - `a baby deer with` == `(a baby deer with:1.0)` | |
| - `(big eyes)` == `(big eyes:1.1)` | |
| - `((big eyes))` == `(big eyes:1.21)` | |
| - `[big eyes]` == `(big eyes:0.91)` | |
| You can run this custom pipeline as so: | |
| #### pytorch | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained( | |
| 'hakurei/waifu-diffusion', | |
| custom_pipeline="lpw_stable_diffusion", | |
| torch_dtype=torch.float16 | |
| ) | |
| pipe=pipe.to("cuda") | |
| prompt = "best_quality (1girl:1.3) bow bride brown_hair closed_mouth frilled_bow frilled_hair_tubes frills (full_body:1.3) fox_ear hair_bow hair_tubes happy hood japanese_clothes kimono long_sleeves red_bow smile solo tabi uchikake white_kimono wide_sleeves cherry_blossoms" | |
| neg_prompt = "lowres, bad_anatomy, error_body, error_hair, error_arm, error_hands, bad_hands, error_fingers, bad_fingers, missing_fingers, error_legs, bad_legs, multiple_legs, missing_legs, error_lighting, error_shadow, error_reflection, text, error, extra_digit, fewer_digits, cropped, worst_quality, low_quality, normal_quality, jpeg_artifacts, signature, watermark, username, blurry" | |
| pipe.text2img(prompt, negative_prompt=neg_prompt, width=512,height=512,max_embeddings_multiples=3).images[0] | |
| ``` | |
| #### onnxruntime | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained( | |
| 'CompVis/stable-diffusion-v1-4', | |
| custom_pipeline="lpw_stable_diffusion_onnx", | |
| revision="onnx", | |
| provider="CUDAExecutionProvider" | |
| ) | |
| prompt = "a photo of an astronaut riding a horse on mars, best quality" | |
| neg_prompt = "lowres, bad anatomy, error body, error hair, error arm, error hands, bad hands, error fingers, bad fingers, missing fingers, error legs, bad legs, multiple legs, missing legs, error lighting, error shadow, error reflection, text, error, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry" | |
| pipe.text2img(prompt,negative_prompt=neg_prompt, width=512, height=512, max_embeddings_multiples=3).images[0] | |
| ``` | |
| if you see `Token indices sequence length is longer than the specified maximum sequence length for this model ( *** > 77 ) . Running this sequence through the model will result in indexing errors`. Do not worry, it is normal. | |
| ### Speech to Image | |
| The following code can generate an image from an audio sample using pre-trained OpenAI whisper-small and Stable Diffusion. | |
| ```Python | |
| import torch | |
| import matplotlib.pyplot as plt | |
| from datasets import load_dataset | |
| from diffusers import DiffusionPipeline | |
| from transformers import ( | |
| WhisperForConditionalGeneration, | |
| WhisperProcessor, | |
| ) | |
| device = "cuda" if torch.cuda.is_available() else "cpu" | |
| ds = load_dataset("hf-internal-testing/librispeech_asr_dummy", "clean", split="validation") | |
| audio_sample = ds[3] | |
| text = audio_sample["text"].lower() | |
| speech_data = audio_sample["audio"]["array"] | |
| model = WhisperForConditionalGeneration.from_pretrained("openai/whisper-small").to(device) | |
| processor = WhisperProcessor.from_pretrained("openai/whisper-small") | |
| diffuser_pipeline = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="speech_to_image_diffusion", | |
| speech_model=model, | |
| speech_processor=processor, | |
| torch_dtype=torch.float16, | |
| ) | |
| diffuser_pipeline.enable_attention_slicing() | |
| diffuser_pipeline = diffuser_pipeline.to(device) | |
| output = diffuser_pipeline(speech_data) | |
| plt.imshow(output.images[0]) | |
| ``` | |
| This example produces the following image: | |
|  | |
| ### Wildcard Stable Diffusion | |
| Following the great examples from <https://github.com/jtkelm2/stable-diffusion-webui-1/blob/master/scripts/wildcards.py> and <https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Custom-Scripts#wildcards>, here's a minimal implementation that allows for users to add "wildcards", denoted by `__wildcard__` to prompts that are used as placeholders for randomly sampled values given by either a dictionary or a `.txt` file. For example: | |
| Say we have a prompt: | |
| ``` | |
| prompt = "__animal__ sitting on a __object__ wearing a __clothing__" | |
| ``` | |
| We can then define possible values to be sampled for `animal`, `object`, and `clothing`. These can either be from a `.txt` with the same name as the category. | |
| The possible values can also be defined / combined by using a dictionary like: `{"animal":["dog", "cat", mouse"]}`. | |
| The actual pipeline works just like `StableDiffusionPipeline`, except the `__call__` method takes in: | |
| `wildcard_files`: list of file paths for wild card replacement | |
| `wildcard_option_dict`: dict with key as `wildcard` and values as a list of possible replacements | |
| `num_prompt_samples`: number of prompts to sample, uniformly sampling wildcards | |
| A full example: | |
| create `animal.txt`, with contents like: | |
| ``` | |
| dog | |
| cat | |
| mouse | |
| ``` | |
| create `object.txt`, with contents like: | |
| ``` | |
| chair | |
| sofa | |
| bench | |
| ``` | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="wildcard_stable_diffusion", | |
| torch_dtype=torch.float16, | |
| ) | |
| prompt = "__animal__ sitting on a __object__ wearing a __clothing__" | |
| out = pipe( | |
| prompt, | |
| wildcard_option_dict={ | |
| "clothing":["hat", "shirt", "scarf", "beret"] | |
| }, | |
| wildcard_files=["object.txt", "animal.txt"], | |
| num_prompt_samples=1 | |
| ) | |
| ``` | |
| ### Composable Stable diffusion | |
| [Composable Stable Diffusion](https://energy-based-model.github.io/Compositional-Visual-Generation-with-Composable-Diffusion-Models/) proposes conjunction and negation (negative prompts) operators for compositional generation with conditional diffusion models. | |
| ```python | |
| import torch as th | |
| import numpy as np | |
| import torchvision.utils as tvu | |
| from diffusers import DiffusionPipeline | |
| import argparse | |
| parser = argparse.ArgumentParser() | |
| parser.add_argument("--prompt", type=str, default="mystical trees | A magical pond | dark", | |
| help="use '|' as the delimiter to compose separate sentences.") | |
| parser.add_argument("--steps", type=int, default=50) | |
| parser.add_argument("--scale", type=float, default=7.5) | |
| parser.add_argument("--weights", type=str, default="7.5 | 7.5 | -7.5") | |
| parser.add_argument("--seed", type=int, default=2) | |
| parser.add_argument("--model_path", type=str, default="CompVis/stable-diffusion-v1-4") | |
| parser.add_argument("--num_images", type=int, default=1) | |
| args = parser.parse_args() | |
| has_cuda = th.cuda.is_available() | |
| device = th.device('cpu' if not has_cuda else 'cuda') | |
| prompt = args.prompt | |
| scale = args.scale | |
| steps = args.steps | |
| pipe = DiffusionPipeline.from_pretrained( | |
| args.model_path, | |
| custom_pipeline="composable_stable_diffusion", | |
| ).to(device) | |
| pipe.safety_checker = None | |
| images = [] | |
| generator = th.Generator("cuda").manual_seed(args.seed) | |
| for i in range(args.num_images): | |
| image = pipe(prompt, guidance_scale=scale, num_inference_steps=steps, | |
| weights=args.weights, generator=generator).images[0] | |
| images.append(th.from_numpy(np.array(image)).permute(2, 0, 1) / 255.) | |
| grid = tvu.make_grid(th.stack(images, dim=0), nrow=4, padding=0) | |
| tvu.save_image(grid, f'{prompt}_{args.weights}' + '.png') | |
| ``` | |
| ### Imagic Stable Diffusion | |
| Allows you to edit an image using stable diffusion. | |
| ```python | |
| import requests | |
| from PIL import Image | |
| from io import BytesIO | |
| import torch | |
| import os | |
| from diffusers import DiffusionPipeline, DDIMScheduler | |
| has_cuda = torch.cuda.is_available() | |
| device = torch.device('cpu' if not has_cuda else 'cuda') | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| safety_checker=None, | |
| custom_pipeline="imagic_stable_diffusion", | |
| scheduler = DDIMScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False) | |
| ).to(device) | |
| generator = torch.Generator("cuda").manual_seed(0) | |
| seed = 0 | |
| prompt = "A photo of Barack Obama smiling with a big grin" | |
| url = 'https://www.dropbox.com/s/6tlwzr73jd1r9yk/obama.png?dl=1' | |
| response = requests.get(url) | |
| init_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| init_image = init_image.resize((512, 512)) | |
| res = pipe.train( | |
| prompt, | |
| image=init_image, | |
| generator=generator) | |
| res = pipe(alpha=1, guidance_scale=7.5, num_inference_steps=50) | |
| os.makedirs("imagic", exist_ok=True) | |
| image = res.images[0] | |
| image.save('./imagic/imagic_image_alpha_1.png') | |
| res = pipe(alpha=1.5, guidance_scale=7.5, num_inference_steps=50) | |
| image = res.images[0] | |
| image.save('./imagic/imagic_image_alpha_1_5.png') | |
| res = pipe(alpha=2, guidance_scale=7.5, num_inference_steps=50) | |
| image = res.images[0] | |
| image.save('./imagic/imagic_image_alpha_2.png') | |
| ``` | |
| ### Seed Resizing | |
| Test seed resizing. Originally generate an image in 512 by 512, then generate image with same seed at 512 by 592 using seed resizing. Finally, generate 512 by 592 using original stable diffusion pipeline. | |
| ```python | |
| import torch as th | |
| import numpy as np | |
| from diffusers import DiffusionPipeline | |
| has_cuda = th.cuda.is_available() | |
| device = th.device('cpu' if not has_cuda else 'cuda') | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="seed_resize_stable_diffusion" | |
| ).to(device) | |
| def dummy(images, **kwargs): | |
| return images, False | |
| pipe.safety_checker = dummy | |
| images = [] | |
| th.manual_seed(0) | |
| generator = th.Generator("cuda").manual_seed(0) | |
| seed = 0 | |
| prompt = "A painting of a futuristic cop" | |
| width = 512 | |
| height = 512 | |
| res = pipe( | |
| prompt, | |
| guidance_scale=7.5, | |
| num_inference_steps=50, | |
| height=height, | |
| width=width, | |
| generator=generator) | |
| image = res.images[0] | |
| image.save('./seed_resize/seed_resize_{w}_{h}_image.png'.format(w=width, h=height)) | |
| th.manual_seed(0) | |
| generator = th.Generator("cuda").manual_seed(0) | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="/home/mark/open_source/diffusers/examples/community/" | |
| ).to(device) | |
| width = 512 | |
| height = 592 | |
| res = pipe( | |
| prompt, | |
| guidance_scale=7.5, | |
| num_inference_steps=50, | |
| height=height, | |
| width=width, | |
| generator=generator) | |
| image = res.images[0] | |
| image.save('./seed_resize/seed_resize_{w}_{h}_image.png'.format(w=width, h=height)) | |
| pipe_compare = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="/home/mark/open_source/diffusers/examples/community/" | |
| ).to(device) | |
| res = pipe_compare( | |
| prompt, | |
| guidance_scale=7.5, | |
| num_inference_steps=50, | |
| height=height, | |
| width=width, | |
| generator=generator | |
| ) | |
| image = res.images[0] | |
| image.save('./seed_resize/seed_resize_{w}_{h}_image_compare.png'.format(w=width, h=height)) | |
| ``` | |
| ### Multilingual Stable Diffusion Pipeline | |
| The following code can generate an images from texts in different languages using the pre-trained [mBART-50 many-to-one multilingual machine translation model](https://huggingface.co/facebook/mbart-large-50-many-to-one-mmt) and Stable Diffusion. | |
| ```python | |
| from PIL import Image | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| from transformers import ( | |
| pipeline, | |
| MBart50TokenizerFast, | |
| MBartForConditionalGeneration, | |
| ) | |
| device = "cuda" if torch.cuda.is_available() else "cpu" | |
| device_dict = {"cuda": 0, "cpu": -1} | |
| # helper function taken from: https://huggingface.co/blog/stable_diffusion | |
| def image_grid(imgs, rows, cols): | |
| assert len(imgs) == rows*cols | |
| w, h = imgs[0].size | |
| grid = Image.new('RGB', size=(cols*w, rows*h)) | |
| grid_w, grid_h = grid.size | |
| for i, img in enumerate(imgs): | |
| grid.paste(img, box=(i%cols*w, i//cols*h)) | |
| return grid | |
| # Add language detection pipeline | |
| language_detection_model_ckpt = "papluca/xlm-roberta-base-language-detection" | |
| language_detection_pipeline = pipeline("text-classification", | |
| model=language_detection_model_ckpt, | |
| device=device_dict[device]) | |
| # Add model for language translation | |
| trans_tokenizer = MBart50TokenizerFast.from_pretrained("facebook/mbart-large-50-many-to-one-mmt") | |
| trans_model = MBartForConditionalGeneration.from_pretrained("facebook/mbart-large-50-many-to-one-mmt").to(device) | |
| diffuser_pipeline = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="multilingual_stable_diffusion", | |
| detection_pipeline=language_detection_pipeline, | |
| translation_model=trans_model, | |
| translation_tokenizer=trans_tokenizer, | |
| torch_dtype=torch.float16, | |
| ) | |
| diffuser_pipeline.enable_attention_slicing() | |
| diffuser_pipeline = diffuser_pipeline.to(device) | |
| prompt = ["a photograph of an astronaut riding a horse", | |
| "Una casa en la playa", | |
| "Ein Hund, der Orange isst", | |
| "Un restaurant parisien"] | |
| output = diffuser_pipeline(prompt) | |
| images = output.images | |
| grid = image_grid(images, rows=2, cols=2) | |
| ``` | |
| This example produces the following images: | |
|  | |
| ### GlueGen Stable Diffusion Pipeline | |
| GlueGen is a minimal adapter that allow alignment between any encoder (Text Encoder of different language, Multilingual Roberta, AudioClip) and CLIP text encoder used in standard Stable Diffusion model. This method allows easy language adaptation to available english Stable Diffusion checkpoints without the need of an image captioning dataset as well as long training hours. | |
| Make sure you downloaded `gluenet_French_clip_overnorm_over3_noln.ckpt` for French (there are also pre-trained weights for Chinese, Italian, Japanese, Spanish or train your own) at [GlueGen's official repo](https://github.com/salesforce/GlueGen/tree/main) | |
| ```python | |
| from PIL import Image | |
| import torch | |
| from transformers import AutoModel, AutoTokenizer | |
| from diffusers import DiffusionPipeline | |
| if __name__ == "__main__": | |
| device = "cuda" | |
| lm_model_id = "xlm-roberta-large" | |
| token_max_length = 77 | |
| text_encoder = AutoModel.from_pretrained(lm_model_id) | |
| tokenizer = AutoTokenizer.from_pretrained(lm_model_id, model_max_length=token_max_length, use_fast=False) | |
| tensor_norm = torch.Tensor([[43.8203],[28.3668],[27.9345],[28.0084],[28.2958],[28.2576],[28.3373],[28.2695],[28.4097],[28.2790],[28.2825],[28.2807],[28.2775],[28.2708],[28.2682],[28.2624],[28.2589],[28.2611],[28.2616],[28.2639],[28.2613],[28.2566],[28.2615],[28.2665],[28.2799],[28.2885],[28.2852],[28.2863],[28.2780],[28.2818],[28.2764],[28.2532],[28.2412],[28.2336],[28.2514],[28.2734],[28.2763],[28.2977],[28.2971],[28.2948],[28.2818],[28.2676],[28.2831],[28.2890],[28.2979],[28.2999],[28.3117],[28.3363],[28.3554],[28.3626],[28.3589],[28.3597],[28.3543],[28.3660],[28.3731],[28.3717],[28.3812],[28.3753],[28.3810],[28.3777],[28.3693],[28.3713],[28.3670],[28.3691],[28.3679],[28.3624],[28.3703],[28.3703],[28.3720],[28.3594],[28.3576],[28.3562],[28.3438],[28.3376],[28.3389],[28.3433],[28.3191]]) | |
| pipeline = DiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", | |
| text_encoder=text_encoder, | |
| tokenizer=tokenizer, | |
| custom_pipeline="gluegen" | |
| ).to(device) | |
| pipeline.load_language_adapter("gluenet_French_clip_overnorm_over3_noln.ckpt", num_token=token_max_length, dim=1024, dim_out=768, tensor_norm=tensor_norm) | |
| prompt = "une voiture sur la plage" | |
| generator = torch.Generator(device=device).manual_seed(42) | |
| image = pipeline(prompt, generator=generator).images[0] | |
| image.save("gluegen_output_fr.png") | |
| ``` | |
| Which will produce: | |
|  | |
| ### Image to Image Inpainting Stable Diffusion | |
| Similar to the standard stable diffusion inpainting example, except with the addition of an `inner_image` argument. | |
| `image`, `inner_image`, and `mask` should have the same dimensions. `inner_image` should have an alpha (transparency) channel. | |
| The aim is to overlay two images, then mask out the boundary between `image` and `inner_image` to allow stable diffusion to make the connection more seamless. | |
| For example, this could be used to place a logo on a shirt and make it blend seamlessly. | |
| ```python | |
| import PIL | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| image_path = "./path-to-image.png" | |
| inner_image_path = "./path-to-inner-image.png" | |
| mask_path = "./path-to-mask.png" | |
| init_image = PIL.Image.open(image_path).convert("RGB").resize((512, 512)) | |
| inner_image = PIL.Image.open(inner_image_path).convert("RGBA").resize((512, 512)) | |
| mask_image = PIL.Image.open(mask_path).convert("RGB").resize((512, 512)) | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-inpainting", | |
| custom_pipeline="img2img_inpainting", | |
| torch_dtype=torch.float16 | |
| ) | |
| pipe = pipe.to("cuda") | |
| prompt = "Your prompt here!" | |
| image = pipe(prompt=prompt, image=init_image, inner_image=inner_image, mask_image=mask_image).images[0] | |
| ``` | |
|  | |
| ### Text Based Inpainting Stable Diffusion | |
| Use a text prompt to generate the mask for the area to be inpainted. | |
| Currently uses the CLIPSeg model for mask generation, then calls the standard Stable Diffusion Inpainting pipeline to perform the inpainting. | |
| ```python | |
| from transformers import CLIPSegProcessor, CLIPSegForImageSegmentation | |
| from diffusers import DiffusionPipeline | |
| from PIL import Image | |
| import requests | |
| processor = CLIPSegProcessor.from_pretrained("CIDAS/clipseg-rd64-refined") | |
| model = CLIPSegForImageSegmentation.from_pretrained("CIDAS/clipseg-rd64-refined") | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-inpainting", | |
| custom_pipeline="text_inpainting", | |
| segmentation_model=model, | |
| segmentation_processor=processor | |
| ) | |
| pipe = pipe.to("cuda") | |
| url = "https://github.com/timojl/clipseg/blob/master/example_image.jpg?raw=true" | |
| image = Image.open(requests.get(url, stream=True).raw).resize((512, 512)) | |
| text = "a glass" # will mask out this text | |
| prompt = "a cup" # the masked out region will be replaced with this | |
| image = pipe(image=image, text=text, prompt=prompt).images[0] | |
| ``` | |
| ### Bit Diffusion | |
| Based <https://arxiv.org/abs/2208.04202>, this is used for diffusion on discrete data - eg, discreate image data, DNA sequence data. An unconditional discreate image can be generated like this: | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained("google/ddpm-cifar10-32", custom_pipeline="bit_diffusion") | |
| image = pipe().images[0] | |
| ``` | |
| ### Stable Diffusion with K Diffusion | |
| Make sure you have @crowsonkb's <https://github.com/crowsonkb/k-diffusion> installed: | |
| ```sh | |
| pip install k-diffusion | |
| ``` | |
| You can use the community pipeline as follows: | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="sd_text2img_k_diffusion") | |
| pipe = pipe.to("cuda") | |
| prompt = "an astronaut riding a horse on mars" | |
| pipe.set_scheduler("sample_heun") | |
| generator = torch.Generator(device="cuda").manual_seed(seed) | |
| image = pipe(prompt, generator=generator, num_inference_steps=20).images[0] | |
| image.save("./astronaut_heun_k_diffusion.png") | |
| ``` | |
| To make sure that K Diffusion and `diffusers` yield the same results: | |
| **Diffusers**: | |
| ```python | |
| from diffusers import DiffusionPipeline, EulerDiscreteScheduler | |
| seed = 33 | |
| pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4") | |
| pipe.scheduler = EulerDiscreteScheduler.from_config(pipe.scheduler.config) | |
| pipe = pipe.to("cuda") | |
| generator = torch.Generator(device="cuda").manual_seed(seed) | |
| image = pipe(prompt, generator=generator, num_inference_steps=50).images[0] | |
| ``` | |
|  | |
| **K Diffusion**: | |
| ```python | |
| from diffusers import DiffusionPipeline, EulerDiscreteScheduler | |
| seed = 33 | |
| pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="sd_text2img_k_diffusion") | |
| pipe.scheduler = EulerDiscreteScheduler.from_config(pipe.scheduler.config) | |
| pipe = pipe.to("cuda") | |
| pipe.set_scheduler("sample_euler") | |
| generator = torch.Generator(device="cuda").manual_seed(seed) | |
| image = pipe(prompt, generator=generator, num_inference_steps=50).images[0] | |
| ``` | |
|  | |
| ### Checkpoint Merger Pipeline | |
| Based on the AUTOMATIC1111/webui for checkpoint merging. This is a custom pipeline that merges upto 3 pretrained model checkpoints as long as they are in the HuggingFace model_index.json format. | |
| The checkpoint merging is currently memory intensive as it modifies the weights of a DiffusionPipeline object in place. Expect at least 13GB RAM Usage on Kaggle GPU kernels and | |
| on colab you might run out of the 12GB memory even while merging two checkpoints. | |
| Usage:- | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| #Return a CheckpointMergerPipeline class that allows you to merge checkpoints. | |
| #The checkpoint passed here is ignored. But still pass one of the checkpoints you plan to | |
| #merge for convenience | |
| pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="checkpoint_merger") | |
| #There are multiple possible scenarios: | |
| #The pipeline with the merged checkpoints is returned in all the scenarios | |
| #Compatible checkpoints a.k.a matched model_index.json files. Ignores the meta attributes in model_index.json during comparison.( attrs with _ as prefix ) | |
| merged_pipe = pipe.merge(["CompVis/stable-diffusion-v1-4","CompVis/stable-diffusion-v1-2"], interp = "sigmoid", alpha = 0.4) | |
| #Incompatible checkpoints in model_index.json but merge might be possible. Use force = True to ignore model_index.json compatibility | |
| merged_pipe_1 = pipe.merge(["CompVis/stable-diffusion-v1-4","hakurei/waifu-diffusion"], force = True, interp = "sigmoid", alpha = 0.4) | |
| #Three checkpoint merging. Only "add_difference" method actually works on all three checkpoints. Using any other options will ignore the 3rd checkpoint. | |
| merged_pipe_2 = pipe.merge(["CompVis/stable-diffusion-v1-4","hakurei/waifu-diffusion","prompthero/openjourney"], force = True, interp = "add_difference", alpha = 0.4) | |
| prompt = "An astronaut riding a horse on Mars" | |
| image = merged_pipe(prompt).images[0] | |
| ``` | |
| Some examples along with the merge details: | |
| 1. "CompVis/stable-diffusion-v1-4" + "hakurei/waifu-diffusion" ; Sigmoid interpolation; alpha = 0.8 | |
|  | |
| 2. "hakurei/waifu-diffusion" + "prompthero/openjourney" ; Inverse Sigmoid interpolation; alpha = 0.8 | |
|  | |
| 3. "CompVis/stable-diffusion-v1-4" + "hakurei/waifu-diffusion" + "prompthero/openjourney"; Add Difference interpolation; alpha = 0.5 | |
|  | |
| ### Stable Diffusion Comparisons | |
| This Community Pipeline enables the comparison between the 4 checkpoints that exist for Stable Diffusion. They can be found through the following links: | |
| 1. [Stable Diffusion v1.1](https://huggingface.co/CompVis/stable-diffusion-v1-1) | |
| 2. [Stable Diffusion v1.2](https://huggingface.co/CompVis/stable-diffusion-v1-2) | |
| 3. [Stable Diffusion v1.3](https://huggingface.co/CompVis/stable-diffusion-v1-3) | |
| 4. [Stable Diffusion v1.4](https://huggingface.co/CompVis/stable-diffusion-v1-4) | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import matplotlib.pyplot as plt | |
| pipe = DiffusionPipeline.from_pretrained('CompVis/stable-diffusion-v1-4', custom_pipeline='suvadityamuk/StableDiffusionComparison') | |
| pipe.enable_attention_slicing() | |
| pipe = pipe.to('cuda') | |
| prompt = "an astronaut riding a horse on mars" | |
| output = pipe(prompt) | |
| plt.subplots(2,2,1) | |
| plt.imshow(output.images[0]) | |
| plt.title('Stable Diffusion v1.1') | |
| plt.axis('off') | |
| plt.subplots(2,2,2) | |
| plt.imshow(output.images[1]) | |
| plt.title('Stable Diffusion v1.2') | |
| plt.axis('off') | |
| plt.subplots(2,2,3) | |
| plt.imshow(output.images[2]) | |
| plt.title('Stable Diffusion v1.3') | |
| plt.axis('off') | |
| plt.subplots(2,2,4) | |
| plt.imshow(output.images[3]) | |
| plt.title('Stable Diffusion v1.4') | |
| plt.axis('off') | |
| plt.show() | |
| ``` | |
| As a result, you can look at a grid of all 4 generated images being shown together, that captures a difference the advancement of the training between the 4 checkpoints. | |
| ### Magic Mix | |
| Implementation of the [MagicMix: Semantic Mixing with Diffusion Models](https://arxiv.org/abs/2210.16056) paper. This is a Diffusion Pipeline for semantic mixing of an image and a text prompt to create a new concept while preserving the spatial layout and geometry of the subject in the image. The pipeline takes an image that provides the layout semantics and a prompt that provides the content semantics for the mixing process. | |
| There are 3 parameters for the method- | |
| - `mix_factor`: It is the interpolation constant used in the layout generation phase. The greater the value of `mix_factor`, the greater the influence of the prompt on the layout generation process. | |
| - `kmax` and `kmin`: These determine the range for the layout and content generation process. A higher value of kmax results in loss of more information about the layout of the original image and a higher value of kmin results in more steps for content generation process. | |
| Here is an example usage- | |
| ```python | |
| from diffusers import DiffusionPipeline, DDIMScheduler | |
| from PIL import Image | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="magic_mix", | |
| scheduler = DDIMScheduler.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="scheduler"), | |
| ).to('cuda') | |
| img = Image.open('phone.jpg') | |
| mix_img = pipe( | |
| img, | |
| prompt = 'bed', | |
| kmin = 0.3, | |
| kmax = 0.5, | |
| mix_factor = 0.5, | |
| ) | |
| mix_img.save('phone_bed_mix.jpg') | |
| ``` | |
| The `mix_img` is a PIL image that can be saved locally or displayed directly in a google colab. Generated image is a mix of the layout semantics of the given image and the content semantics of the prompt. | |
| E.g. the above script generates the following image: | |
| `phone.jpg` | |
|  | |
| `phone_bed_mix.jpg` | |
|  | |
| For more example generations check out this [demo notebook](https://github.com/daspartho/MagicMix/blob/main/demo.ipynb). | |
| ### Stable UnCLIP | |
| UnCLIPPipeline("kakaobrain/karlo-v1-alpha") provide a prior model that can generate clip image embedding from text. | |
| StableDiffusionImageVariationPipeline("lambdalabs/sd-image-variations-diffusers") provide a decoder model than can generate images from clip image embedding. | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| device = torch.device("cpu" if not torch.cuda.is_available() else "cuda") | |
| pipeline = DiffusionPipeline.from_pretrained( | |
| "kakaobrain/karlo-v1-alpha", | |
| torch_dtype=torch.float16, | |
| custom_pipeline="stable_unclip", | |
| decoder_pipe_kwargs=dict( | |
| image_encoder=None, | |
| ), | |
| ) | |
| pipeline.to(device) | |
| prompt = "a shiba inu wearing a beret and black turtleneck" | |
| random_generator = torch.Generator(device=device).manual_seed(1000) | |
| output = pipeline( | |
| prompt=prompt, | |
| width=512, | |
| height=512, | |
| generator=random_generator, | |
| prior_guidance_scale=4, | |
| prior_num_inference_steps=25, | |
| decoder_guidance_scale=8, | |
| decoder_num_inference_steps=50, | |
| ) | |
| image = output.images[0] | |
| image.save("./shiba-inu.jpg") | |
| # debug | |
| # `pipeline.decoder_pipe` is a regular StableDiffusionImageVariationPipeline instance. | |
| # It is used to convert clip image embedding to latents, then fed into VAE decoder. | |
| print(pipeline.decoder_pipe.__class__) | |
| # <class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_image_variation.StableDiffusionImageVariationPipeline'> | |
| # this pipeline only use prior module in "kakaobrain/karlo-v1-alpha" | |
| # It is used to convert clip text embedding to clip image embedding. | |
| print(pipeline) | |
| # StableUnCLIPPipeline { | |
| # "_class_name": "StableUnCLIPPipeline", | |
| # "_diffusers_version": "0.12.0.dev0", | |
| # "prior": [ | |
| # "diffusers", | |
| # "PriorTransformer" | |
| # ], | |
| # "prior_scheduler": [ | |
| # "diffusers", | |
| # "UnCLIPScheduler" | |
| # ], | |
| # "text_encoder": [ | |
| # "transformers", | |
| # "CLIPTextModelWithProjection" | |
| # ], | |
| # "tokenizer": [ | |
| # "transformers", | |
| # "CLIPTokenizer" | |
| # ] | |
| # } | |
| # pipeline.prior_scheduler is the scheduler used for prior in UnCLIP. | |
| print(pipeline.prior_scheduler) | |
| # UnCLIPScheduler { | |
| # "_class_name": "UnCLIPScheduler", | |
| # "_diffusers_version": "0.12.0.dev0", | |
| # "clip_sample": true, | |
| # "clip_sample_range": 5.0, | |
| # "num_train_timesteps": 1000, | |
| # "prediction_type": "sample", | |
| # "variance_type": "fixed_small_log" | |
| # } | |
| ``` | |
| `shiba-inu.jpg` | |
|  | |
| ### UnCLIP Text Interpolation Pipeline | |
| This Diffusion Pipeline takes two prompts and interpolates between the two input prompts using spherical interpolation ( slerp ). The input prompts are converted to text embeddings by the pipeline's text_encoder and the interpolation is done on the resulting text_embeddings over the number of steps specified. Defaults to 5 steps. | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| device = torch.device("cpu" if not torch.cuda.is_available() else "cuda") | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "kakaobrain/karlo-v1-alpha", | |
| torch_dtype=torch.float16, | |
| custom_pipeline="unclip_text_interpolation" | |
| ) | |
| pipe.to(device) | |
| start_prompt = "A photograph of an adult lion" | |
| end_prompt = "A photograph of a lion cub" | |
| #For best results keep the prompts close in length to each other. Of course, feel free to try out with differing lengths. | |
| generator = torch.Generator(device=device).manual_seed(42) | |
| output = pipe(start_prompt, end_prompt, steps = 6, generator = generator, enable_sequential_cpu_offload=False) | |
| for i,image in enumerate(output.images): | |
| img.save('result%s.jpg' % i) | |
| ``` | |
| The resulting images in order:- | |
|  | |
|  | |
|  | |
|  | |
|  | |
|  | |
| ### UnCLIP Image Interpolation Pipeline | |
| This Diffusion Pipeline takes two images or an image_embeddings tensor of size 2 and interpolates between their embeddings using spherical interpolation ( slerp ). The input images/image_embeddings are converted to image embeddings by the pipeline's image_encoder and the interpolation is done on the resulting image_embeddings over the number of steps specified. Defaults to 5 steps. | |
| ```python | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| from PIL import Image | |
| device = torch.device("cpu" if not torch.cuda.is_available() else "cuda") | |
| dtype = torch.float16 if torch.cuda.is_available() else torch.bfloat16 | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "kakaobrain/karlo-v1-alpha-image-variations", | |
| torch_dtype=dtype, | |
| custom_pipeline="unclip_image_interpolation" | |
| ) | |
| pipe.to(device) | |
| images = [Image.open('./starry_night.jpg'), Image.open('./flowers.jpg')] | |
| #For best results keep the prompts close in length to each other. Of course, feel free to try out with differing lengths. | |
| generator = torch.Generator(device=device).manual_seed(42) | |
| output = pipe(image = images ,steps = 6, generator = generator) | |
| for i,image in enumerate(output.images): | |
| image.save('starry_to_flowers_%s.jpg' % i) | |
| ``` | |
| The original images:- | |
|  | |
|  | |
| The resulting images in order:- | |
|  | |
|  | |
|  | |
|  | |
|  | |
|  | |
| ### DDIM Noise Comparative Analysis Pipeline | |
| #### **Research question: What visual concepts do the diffusion models learn from each noise level during training?** | |
| The [P2 weighting (CVPR 2022)](https://arxiv.org/abs/2204.00227) paper proposed an approach to answer the above question, which is their second contribution. | |
| The approach consists of the following steps: | |
| 1. The input is an image x0. | |
| 2. Perturb it to xt using a diffusion process q(xt|x0). | |
| - `strength` is a value between 0.0 and 1.0, that controls the amount of noise that is added to the input image. Values that approach 1.0 allow for lots of variations but will also produce images that are not semantically consistent with the input. | |
| 3. Reconstruct the image with the learned denoising process pθ(ˆx0|xt). | |
| 4. Compare x0 and ˆx0 among various t to show how each step contributes to the sample. | |
| The authors used [openai/guided-diffusion](https://github.com/openai/guided-diffusion) model to denoise images in FFHQ dataset. This pipeline extends their second contribution by investigating DDIM on any input image. | |
| ```python | |
| import torch | |
| from PIL import Image | |
| import numpy as np | |
| image_path = "path/to/your/image" # images from CelebA-HQ might be better | |
| image_pil = Image.open(image_path) | |
| image_name = image_path.split("/")[-1].split(".")[0] | |
| device = torch.device("cpu" if not torch.cuda.is_available() else "cuda") | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "google/ddpm-ema-celebahq-256", | |
| custom_pipeline="ddim_noise_comparative_analysis", | |
| ) | |
| pipe = pipe.to(device) | |
| for strength in np.linspace(0.1, 1, 25): | |
| denoised_image, latent_timestep = pipe( | |
| image_pil, strength=strength, return_dict=False | |
| ) | |
| denoised_image = denoised_image[0] | |
| denoised_image.save( | |
| f"noise_comparative_analysis_{image_name}_{latent_timestep}.png" | |
| ) | |
| ``` | |
| Here is the result of this pipeline (which is DDIM) on CelebA-HQ dataset. | |
|  | |
| ### CLIP Guided Img2Img Stable Diffusion | |
| CLIP guided Img2Img stable diffusion can help to generate more realistic images with an initial image | |
| by guiding stable diffusion at every denoising step with an additional CLIP model. | |
| The following code requires roughly 12GB of GPU RAM. | |
| ```python | |
| from io import BytesIO | |
| import requests | |
| import torch | |
| from diffusers import DiffusionPipeline | |
| from PIL import Image | |
| from transformers import CLIPFeatureExtractor, CLIPModel | |
| feature_extractor = CLIPFeatureExtractor.from_pretrained( | |
| "laion/CLIP-ViT-B-32-laion2B-s34B-b79K" | |
| ) | |
| clip_model = CLIPModel.from_pretrained( | |
| "laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16 | |
| ) | |
| guided_pipeline = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| # custom_pipeline="clip_guided_stable_diffusion", | |
| custom_pipeline="/home/njindal/diffusers/examples/community/clip_guided_stable_diffusion.py", | |
| clip_model=clip_model, | |
| feature_extractor=feature_extractor, | |
| torch_dtype=torch.float16, | |
| ) | |
| guided_pipeline.enable_attention_slicing() | |
| guided_pipeline = guided_pipeline.to("cuda") | |
| prompt = "fantasy book cover, full moon, fantasy forest landscape, golden vector elements, fantasy magic, dark light night, intricate, elegant, sharp focus, illustration, highly detailed, digital painting, concept art, matte, art by WLOP and Artgerm and Albert Bierstadt, masterpiece" | |
| url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg" | |
| response = requests.get(url) | |
| init_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| image = guided_pipeline( | |
| prompt=prompt, | |
| num_inference_steps=30, | |
| image=init_image, | |
| strength=0.75, | |
| guidance_scale=7.5, | |
| clip_guidance_scale=100, | |
| num_cutouts=4, | |
| use_cutouts=False, | |
| ).images[0] | |
| display(image) | |
| ``` | |
| Init Image | |
|  | |
| Output Image | |
|  | |
| ### TensorRT Text2Image Stable Diffusion Pipeline | |
| The TensorRT Pipeline can be used to accelerate the Text2Image Stable Diffusion Inference run. | |
| NOTE: The ONNX conversions and TensorRT engine build may take up to 30 minutes. | |
| ```python | |
| import torch | |
| from diffusers import DDIMScheduler | |
| from diffusers.pipelines.stable_diffusion import StableDiffusionPipeline | |
| # Use the DDIMScheduler scheduler here instead | |
| scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1", | |
| subfolder="scheduler") | |
| pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1", | |
| custom_pipeline="stable_diffusion_tensorrt_txt2img", | |
| revision='fp16', | |
| torch_dtype=torch.float16, | |
| scheduler=scheduler,) | |
| # re-use cached folder to save ONNX models and TensorRT Engines | |
| pipe.set_cached_folder("stabilityai/stable-diffusion-2-1", revision='fp16',) | |
| pipe = pipe.to("cuda") | |
| prompt = "a beautiful photograph of Mt. Fuji during cherry blossom" | |
| image = pipe(prompt).images[0] | |
| image.save('tensorrt_mt_fuji.png') | |
| ``` | |
| ### EDICT Image Editing Pipeline | |
| This pipeline implements the text-guided image editing approach from the paper [EDICT: Exact Diffusion Inversion via Coupled Transformations](https://arxiv.org/abs/2211.12446). You have to pass: | |
| - (`PIL`) `image` you want to edit. | |
| - `base_prompt`: the text prompt describing the current image (before editing). | |
| - `target_prompt`: the text prompt describing with the edits. | |
| ```python | |
| from diffusers import DiffusionPipeline, DDIMScheduler | |
| from transformers import CLIPTextModel | |
| import torch, PIL, requests | |
| from io import BytesIO | |
| from IPython.display import display | |
| def center_crop_and_resize(im): | |
| width, height = im.size | |
| d = min(width, height) | |
| left = (width - d) / 2 | |
| upper = (height - d) / 2 | |
| right = (width + d) / 2 | |
| lower = (height + d) / 2 | |
| return im.crop((left, upper, right, lower)).resize((512, 512)) | |
| torch_dtype = torch.float16 | |
| device = torch.device('cuda' if torch.cuda.is_available() else 'cpu') | |
| # scheduler and text_encoder param values as in the paper | |
| scheduler = DDIMScheduler( | |
| num_train_timesteps=1000, | |
| beta_start=0.00085, | |
| beta_end=0.012, | |
| beta_schedule="scaled_linear", | |
| set_alpha_to_one=False, | |
| clip_sample=False, | |
| ) | |
| text_encoder = CLIPTextModel.from_pretrained( | |
| pretrained_model_name_or_path="openai/clip-vit-large-patch14", | |
| torch_dtype=torch_dtype, | |
| ) | |
| # initialize pipeline | |
| pipeline = DiffusionPipeline.from_pretrained( | |
| pretrained_model_name_or_path="CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="edict_pipeline", | |
| revision="fp16", | |
| scheduler=scheduler, | |
| text_encoder=text_encoder, | |
| leapfrog_steps=True, | |
| torch_dtype=torch_dtype, | |
| ).to(device) | |
| # download image | |
| image_url = "https://huggingface.co/datasets/Joqsan/images/resolve/main/imagenet_dog_1.jpeg" | |
| response = requests.get(image_url) | |
| image = PIL.Image.open(BytesIO(response.content)) | |
| # preprocess it | |
| cropped_image = center_crop_and_resize(image) | |
| # define the prompts | |
| base_prompt = "A dog" | |
| target_prompt = "A golden retriever" | |
| # run the pipeline | |
| result_image = pipeline( | |
| base_prompt=base_prompt, | |
| target_prompt=target_prompt, | |
| image=cropped_image, | |
| ) | |
| display(result_image) | |
| ``` | |
| Init Image | |
|  | |
| Output Image | |
|  | |
| ### Stable Diffusion RePaint | |
| This pipeline uses the [RePaint](https://arxiv.org/abs/2201.09865) logic on the latent space of stable diffusion. It can | |
| be used similarly to other image inpainting pipelines but does not rely on a specific inpainting model. This means you can use | |
| models that are not specifically created for inpainting. | |
| Make sure to use the ```RePaintScheduler``` as shown in the example below. | |
| Disclaimer: The mask gets transferred into latent space, this may lead to unexpected changes on the edge of the masked part. | |
| The inference time is a lot slower. | |
| ```py | |
| import PIL | |
| import requests | |
| import torch | |
| from io import BytesIO | |
| from diffusers import StableDiffusionPipeline, RePaintScheduler | |
| def download_image(url): | |
| response = requests.get(url) | |
| return PIL.Image.open(BytesIO(response.content)).convert("RGB") | |
| img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
| mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
| init_image = download_image(img_url).resize((512, 512)) | |
| mask_image = download_image(mask_url).resize((512, 512)) | |
| mask_image = PIL.ImageOps.invert(mask_image) | |
| pipe = StableDiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16, custom_pipeline="stable_diffusion_repaint", | |
| ) | |
| pipe.scheduler = RePaintScheduler.from_config(pipe.scheduler.config) | |
| pipe = pipe.to("cuda") | |
| prompt = "Face of a yellow cat, high resolution, sitting on a park bench" | |
| image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0] | |
| ``` | |
| ### TensorRT Image2Image Stable Diffusion Pipeline | |
| The TensorRT Pipeline can be used to accelerate the Image2Image Stable Diffusion Inference run. | |
| NOTE: The ONNX conversions and TensorRT engine build may take up to 30 minutes. | |
| ```python | |
| import requests | |
| from io import BytesIO | |
| from PIL import Image | |
| import torch | |
| from diffusers import DDIMScheduler | |
| from diffusers.pipelines.stable_diffusion import StableDiffusionImg2ImgPipeline | |
| # Use the DDIMScheduler scheduler here instead | |
| scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1", | |
| subfolder="scheduler") | |
| pipe = StableDiffusionImg2ImgPipeline.from_pretrained("stabilityai/stable-diffusion-2-1", | |
| custom_pipeline="stable_diffusion_tensorrt_img2img", | |
| revision='fp16', | |
| torch_dtype=torch.float16, | |
| scheduler=scheduler,) | |
| # re-use cached folder to save ONNX models and TensorRT Engines | |
| pipe.set_cached_folder("stabilityai/stable-diffusion-2-1", revision='fp16',) | |
| pipe = pipe.to("cuda") | |
| url = "https://pajoca.com/wp-content/uploads/2022/09/tekito-yamakawa-1.png" | |
| response = requests.get(url) | |
| input_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| prompt = "photorealistic new zealand hills" | |
| image = pipe(prompt, image=input_image, strength=0.75,).images[0] | |
| image.save('tensorrt_img2img_new_zealand_hills.png') | |
| ``` | |
| ### Stable Diffusion BoxDiff | |
| BoxDiff is a training-free method for controlled generation with bounding box coordinates. It shoud work with any Stable Diffusion model. Below shows an example with `stable-diffusion-2-1-base`. | |
| ```py | |
| import torch | |
| from PIL import Image, ImageDraw | |
| from copy import deepcopy | |
| from examples.community.pipeline_stable_diffusion_boxdiff import StableDiffusionBoxDiffPipeline | |
| def draw_box_with_text(img, boxes, names): | |
| colors = ["red", "olive", "blue", "green", "orange", "brown", "cyan", "purple"] | |
| img_new = deepcopy(img) | |
| draw = ImageDraw.Draw(img_new) | |
| W, H = img.size | |
| for bid, box in enumerate(boxes): | |
| draw.rectangle([box[0] * W, box[1] * H, box[2] * W, box[3] * H], outline=colors[bid % len(colors)], width=4) | |
| draw.text((box[0] * W, box[1] * H), names[bid], fill=colors[bid % len(colors)]) | |
| return img_new | |
| pipe = StableDiffusionBoxDiffPipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-2-1-base", | |
| torch_dtype=torch.float16, | |
| ) | |
| pipe.to("cuda") | |
| # example 1 | |
| prompt = "as the aurora lights up the sky, a herd of reindeer leisurely wanders on the grassy meadow, admiring the breathtaking view, a serene lake quietly reflects the magnificent display, and in the distance, a snow-capped mountain stands majestically, fantasy, 8k, highly detailed" | |
| phrases = [ | |
| "aurora", | |
| "reindeer", | |
| "meadow", | |
| "lake", | |
| "mountain" | |
| ] | |
| boxes = [[1,3,512,202], [75,344,421,495], [1,327,508,507], [2,217,507,341], [1,135,509,242]] | |
| # example 2 | |
| # prompt = "A rabbit wearing sunglasses looks very proud" | |
| # phrases = ["rabbit", "sunglasses"] | |
| # boxes = [[67,87,366,512], [66,130,364,262]] | |
| boxes = [[x / 512 for x in box] for box in boxes] | |
| images = pipe( | |
| prompt, | |
| boxdiff_phrases=phrases, | |
| boxdiff_boxes=boxes, | |
| boxdiff_kwargs={ | |
| "attention_res": 16, | |
| "normalize_eot": True | |
| }, | |
| num_inference_steps=50, | |
| guidance_scale=7.5, | |
| generator=torch.manual_seed(42), | |
| safety_checker=None | |
| ).images | |
| draw_box_with_text(images[0], boxes, phrases).save("output.png") | |
| ``` | |
| ### Stable Diffusion Reference | |
| This pipeline uses the Reference Control. Refer to the [sd-webui-controlnet discussion: Reference-only Control](https://github.com/Mikubill/sd-webui-controlnet/discussions/1236)[sd-webui-controlnet discussion: Reference-adain Control](https://github.com/Mikubill/sd-webui-controlnet/discussions/1280). | |
| Based on [this issue](https://github.com/huggingface/diffusers/issues/3566), | |
| - `EulerAncestralDiscreteScheduler` got poor results. | |
| ```py | |
| import torch | |
| from diffusers import UniPCMultistepScheduler | |
| from diffusers.utils import load_image | |
| input_image = load_image("https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png") | |
| pipe = StableDiffusionReferencePipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", | |
| safety_checker=None, | |
| torch_dtype=torch.float16 | |
| ).to('cuda:0') | |
| pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) | |
| result_img = pipe(ref_image=input_image, | |
| prompt="1girl", | |
| num_inference_steps=20, | |
| reference_attn=True, | |
| reference_adain=True).images[0] | |
| ``` | |
| Reference Image | |
|  | |
| Output Image of `reference_attn=True` and `reference_adain=False` | |
|  | |
| Output Image of `reference_attn=False` and `reference_adain=True` | |
|  | |
| Output Image of `reference_attn=True` and `reference_adain=True` | |
|  | |
| ### Stable Diffusion ControlNet Reference | |
| This pipeline uses the Reference Control with ControlNet. Refer to the [sd-webui-controlnet discussion: Reference-only Control](https://github.com/Mikubill/sd-webui-controlnet/discussions/1236)[sd-webui-controlnet discussion: Reference-adain Control](https://github.com/Mikubill/sd-webui-controlnet/discussions/1280). | |
| Based on [this issue](https://github.com/huggingface/diffusers/issues/3566), | |
| - `EulerAncestralDiscreteScheduler` got poor results. | |
| - `guess_mode=True` works well for ControlNet v1.1 | |
| ```py | |
| import cv2 | |
| import torch | |
| import numpy as np | |
| from PIL import Image | |
| from diffusers import UniPCMultistepScheduler | |
| from diffusers.utils import load_image | |
| input_image = load_image("https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png") | |
| # get canny image | |
| image = cv2.Canny(np.array(input_image), 100, 200) | |
| image = image[:, :, None] | |
| image = np.concatenate([image, image, image], axis=2) | |
| canny_image = Image.fromarray(image) | |
| controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16) | |
| pipe = StableDiffusionControlNetReferencePipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", | |
| controlnet=controlnet, | |
| safety_checker=None, | |
| torch_dtype=torch.float16 | |
| ).to('cuda:0') | |
| pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) | |
| result_img = pipe(ref_image=input_image, | |
| prompt="1girl", | |
| image=canny_image, | |
| num_inference_steps=20, | |
| reference_attn=True, | |
| reference_adain=True).images[0] | |
| ``` | |
| Reference Image | |
|  | |
| Output Image | |
|  | |
| ### Stable Diffusion on IPEX | |
| This diffusion pipeline aims to accelarate the inference of Stable-Diffusion on Intel Xeon CPUs with BF16/FP32 precision using [IPEX](https://github.com/intel/intel-extension-for-pytorch). | |
| To use this pipeline, you need to: | |
| 1. Install [IPEX](https://github.com/intel/intel-extension-for-pytorch) | |
| **Note:** For each PyTorch release, there is a corresponding release of the IPEX. Here is the mapping relationship. It is recommended to install Pytorch/IPEX2.0 to get the best performance. | |
| |PyTorch Version|IPEX Version| | |
| |--|--| | |
| |[v2.0.\*](https://github.com/pytorch/pytorch/tree/v2.0.1 "v2.0.1")|[v2.0.\*](https://github.com/intel/intel-extension-for-pytorch/tree/v2.0.100+cpu)| | |
| |[v1.13.\*](https://github.com/pytorch/pytorch/tree/v1.13.0 "v1.13.0")|[v1.13.\*](https://github.com/intel/intel-extension-for-pytorch/tree/v1.13.100+cpu)| | |
| You can simply use pip to install IPEX with the latest version. | |
| ```sh | |
| python -m pip install intel_extension_for_pytorch | |
| ``` | |
| **Note:** To install a specific version, run with the following command: | |
| ```sh | |
| python -m pip install intel_extension_for_pytorch==<version_name> -f https://developer.intel.com/ipex-whl-stable-cpu | |
| ``` | |
| 2. After pipeline initialization, `prepare_for_ipex()` should be called to enable IPEX accelaration. Supported inference datatypes are Float32 and BFloat16. | |
| **Note:** The setting of generated image height/width for `prepare_for_ipex()` should be same as the setting of pipeline inference. | |
| ```python | |
| pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_ipex") | |
| # For Float32 | |
| pipe.prepare_for_ipex(prompt, dtype=torch.float32, height=512, width=512) #value of image height/width should be consistent with the pipeline inference | |
| # For BFloat16 | |
| pipe.prepare_for_ipex(prompt, dtype=torch.bfloat16, height=512, width=512) #value of image height/width should be consistent with the pipeline inference | |
| ``` | |
| Then you can use the ipex pipeline in a similar way to the default stable diffusion pipeline. | |
| ```python | |
| # For Float32 | |
| image = pipe(prompt, num_inference_steps=20, height=512, width=512).images[0] #value of image height/width should be consistent with 'prepare_for_ipex()' | |
| # For BFloat16 | |
| with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16): | |
| image = pipe(prompt, num_inference_steps=20, height=512, width=512).images[0] #value of image height/width should be consistent with 'prepare_for_ipex()' | |
| ``` | |
| The following code compares the performance of the original stable diffusion pipeline with the ipex-optimized pipeline. | |
| ```python | |
| import torch | |
| import intel_extension_for_pytorch as ipex | |
| from diffusers import StableDiffusionPipeline | |
| import time | |
| prompt = "sailing ship in storm by Rembrandt" | |
| model_id = "runwayml/stable-diffusion-v1-5" | |
| # Helper function for time evaluation | |
| def elapsed_time(pipeline, nb_pass=3, num_inference_steps=20): | |
| # warmup | |
| for _ in range(2): | |
| images = pipeline(prompt, num_inference_steps=num_inference_steps, height=512, width=512).images | |
| #time evaluation | |
| start = time.time() | |
| for _ in range(nb_pass): | |
| pipeline(prompt, num_inference_steps=num_inference_steps, height=512, width=512) | |
| end = time.time() | |
| return (end - start) / nb_pass | |
| ############## bf16 inference performance ############### | |
| # 1. IPEX Pipeline initialization | |
| pipe = DiffusionPipeline.from_pretrained(model_id, custom_pipeline="stable_diffusion_ipex") | |
| pipe.prepare_for_ipex(prompt, dtype=torch.bfloat16, height=512, width=512) | |
| # 2. Original Pipeline initialization | |
| pipe2 = StableDiffusionPipeline.from_pretrained(model_id) | |
| # 3. Compare performance between Original Pipeline and IPEX Pipeline | |
| with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16): | |
| latency = elapsed_time(pipe) | |
| print("Latency of StableDiffusionIPEXPipeline--bf16", latency) | |
| latency = elapsed_time(pipe2) | |
| print("Latency of StableDiffusionPipeline--bf16",latency) | |
| ############## fp32 inference performance ############### | |
| # 1. IPEX Pipeline initialization | |
| pipe3 = DiffusionPipeline.from_pretrained(model_id, custom_pipeline="stable_diffusion_ipex") | |
| pipe3.prepare_for_ipex(prompt, dtype=torch.float32, height=512, width=512) | |
| # 2. Original Pipeline initialization | |
| pipe4 = StableDiffusionPipeline.from_pretrained(model_id) | |
| # 3. Compare performance between Original Pipeline and IPEX Pipeline | |
| latency = elapsed_time(pipe3) | |
| print("Latency of StableDiffusionIPEXPipeline--fp32", latency) | |
| latency = elapsed_time(pipe4) | |
| print("Latency of StableDiffusionPipeline--fp32",latency) | |
| ``` | |
| ### Stable Diffusion XL on IPEX | |
| This diffusion pipeline aims to accelarate the inference of Stable-Diffusion XL on Intel Xeon CPUs with BF16/FP32 precision using [IPEX](https://github.com/intel/intel-extension-for-pytorch). | |
| To use this pipeline, you need to: | |
| 1. Install [IPEX](https://github.com/intel/intel-extension-for-pytorch) | |
| **Note:** For each PyTorch release, there is a corresponding release of IPEX. Here is the mapping relationship. It is recommended to install Pytorch/IPEX2.0 to get the best performance. | |
| |PyTorch Version|IPEX Version| | |
| |--|--| | |
| |[v2.0.\*](https://github.com/pytorch/pytorch/tree/v2.0.1 "v2.0.1")|[v2.0.\*](https://github.com/intel/intel-extension-for-pytorch/tree/v2.0.100+cpu)| | |
| |[v1.13.\*](https://github.com/pytorch/pytorch/tree/v1.13.0 "v1.13.0")|[v1.13.\*](https://github.com/intel/intel-extension-for-pytorch/tree/v1.13.100+cpu)| | |
| You can simply use pip to install IPEX with the latest version. | |
| ```sh | |
| python -m pip install intel_extension_for_pytorch | |
| ``` | |
| **Note:** To install a specific version, run with the following command: | |
| ```sh | |
| python -m pip install intel_extension_for_pytorch==<version_name> -f https://developer.intel.com/ipex-whl-stable-cpu | |
| ``` | |
| 2. After pipeline initialization, `prepare_for_ipex()` should be called to enable IPEX accelaration. Supported inference datatypes are Float32 and BFloat16. | |
| **Note:** The values of `height` and `width` used during preparation with `prepare_for_ipex()` should be the same when running inference with the prepared pipeline. | |
| ```python | |
| pipe = StableDiffusionXLPipelineIpex.from_pretrained("stabilityai/sdxl-turbo", low_cpu_mem_usage=True, use_safetensors=True) | |
| # value of image height/width should be consistent with the pipeline inference | |
| # For Float32 | |
| pipe.prepare_for_ipex(torch.float32, prompt, height=512, width=512) | |
| # For BFloat16 | |
| pipe.prepare_for_ipex(torch.bfloat16, prompt, height=512, width=512) | |
| ``` | |
| Then you can use the ipex pipeline in a similar way to the default stable diffusion xl pipeline. | |
| ```python | |
| # value of image height/width should be consistent with 'prepare_for_ipex()' | |
| # For Float32 | |
| image = pipe(prompt, num_inference_steps=num_inference_steps, height=512, width=512, guidance_scale=guidance_scale).images[0] | |
| # For BFloat16 | |
| with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16): | |
| image = pipe(prompt, num_inference_steps=num_inference_steps, height=512, width=512, guidance_scale=guidance_scale).images[0] | |
| ``` | |
| The following code compares the performance of the original stable diffusion xl pipeline with the ipex-optimized pipeline. | |
| By using this optimized pipeline, we can get about 1.4-2 times performance boost with BFloat16 on fourth generation of Intel Xeon CPUs, | |
| code-named Sapphire Rapids. | |
| ```python | |
| import torch | |
| from diffusers import StableDiffusionXLPipeline | |
| from pipeline_stable_diffusion_xl_ipex import StableDiffusionXLPipelineIpex | |
| import time | |
| prompt = "sailing ship in storm by Rembrandt" | |
| model_id = "stabilityai/sdxl-turbo" | |
| steps = 4 | |
| # Helper function for time evaluation | |
| def elapsed_time(pipeline, nb_pass=3, num_inference_steps=1): | |
| # warmup | |
| for _ in range(2): | |
| images = pipeline(prompt, num_inference_steps=num_inference_steps, height=512, width=512, guidance_scale=0.0).images | |
| #time evaluation | |
| start = time.time() | |
| for _ in range(nb_pass): | |
| pipeline(prompt, num_inference_steps=num_inference_steps, height=512, width=512, guidance_scale=0.0) | |
| end = time.time() | |
| return (end - start) / nb_pass | |
| ############## bf16 inference performance ############### | |
| # 1. IPEX Pipeline initialization | |
| pipe = StableDiffusionXLPipelineIpex.from_pretrained(model_id, low_cpu_mem_usage=True, use_safetensors=True) | |
| pipe.prepare_for_ipex(torch.bfloat16, prompt, height=512, width=512) | |
| # 2. Original Pipeline initialization | |
| pipe2 = StableDiffusionXLPipeline.from_pretrained(model_id, low_cpu_mem_usage=True, use_safetensors=True) | |
| # 3. Compare performance between Original Pipeline and IPEX Pipeline | |
| with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16): | |
| latency = elapsed_time(pipe, num_inference_steps=steps) | |
| print("Latency of StableDiffusionXLPipelineIpex--bf16", latency, "s for total", steps, "steps") | |
| latency = elapsed_time(pipe2, num_inference_steps=steps) | |
| print("Latency of StableDiffusionXLPipeline--bf16", latency, "s for total", steps, "steps") | |
| ############## fp32 inference performance ############### | |
| # 1. IPEX Pipeline initialization | |
| pipe3 = StableDiffusionXLPipelineIpex.from_pretrained(model_id, low_cpu_mem_usage=True, use_safetensors=True) | |
| pipe3.prepare_for_ipex(torch.float32, prompt, height=512, width=512) | |
| # 2. Original Pipeline initialization | |
| pipe4 = StableDiffusionXLPipeline.from_pretrained(model_id, low_cpu_mem_usage=True, use_safetensors=True) | |
| # 3. Compare performance between Original Pipeline and IPEX Pipeline | |
| latency = elapsed_time(pipe3, num_inference_steps=steps) | |
| print("Latency of StableDiffusionXLPipelineIpex--fp32", latency, "s for total", steps, "steps") | |
| latency = elapsed_time(pipe4, num_inference_steps=steps) | |
| print("Latency of StableDiffusionXLPipeline--fp32",latency, "s for total", steps, "steps") | |
| ``` | |
| ### CLIP Guided Images Mixing With Stable Diffusion | |
|  | |
| CLIP guided stable diffusion images mixing pipeline allows to combine two images using standard diffusion models. | |
| This approach is using (optional) CoCa model to avoid writing image description. | |
| [More code examples](https://github.com/TheDenk/images_mixing) | |
| ### Stable Diffusion XL Long Weighted Prompt Pipeline | |
| This SDXL pipeline support unlimited length prompt and negative prompt, compatible with A1111 prompt weighted style. | |
| You can provide both `prompt` and `prompt_2`. If only one prompt is provided, `prompt_2` will be a copy of the provided `prompt`. Here is a sample code to use this pipeline. | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| from diffusers.utils import load_image | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-xl-base-1.0" | |
| , torch_dtype = torch.float16 | |
| , use_safetensors = True | |
| , variant = "fp16" | |
| , custom_pipeline = "lpw_stable_diffusion_xl", | |
| ) | |
| prompt = "photo of a cute (white) cat running on the grass" * 20 | |
| prompt2 = "chasing (birds:1.5)" * 20 | |
| prompt = f"{prompt},{prompt2}" | |
| neg_prompt = "blur, low quality, carton, animate" | |
| pipe.to("cuda") | |
| # text2img | |
| t2i_images = pipe( | |
| prompt=prompt, | |
| negative_prompt=neg_prompt, | |
| ).images # alternatively, you can call the .text2img() function | |
| # img2img | |
| input_image = load_image("/path/to/local/image.png") # or URL to your input image | |
| i2i_images = pipe.img2img( | |
| prompt=prompt, | |
| negative_prompt=neg_prompt, | |
| image=input_image, | |
| strength=0.8, # higher strength will result in more variation compared to original image | |
| ).images | |
| # inpaint | |
| input_mask = load_image("/path/to/local/mask.png") # or URL to your input inpainting mask | |
| inpaint_images = pipe.inpaint( | |
| prompt="photo of a cute (black) cat running on the grass" * 20, | |
| negative_prompt=neg_prompt, | |
| image=input_image, | |
| mask=input_mask, | |
| strength=0.6, # higher strength will result in more variation compared to original image | |
| ).images | |
| pipe.to("cpu") | |
| torch.cuda.empty_cache() | |
| from IPython.display import display # assuming you are using this code in a notebook | |
| display(t2i_images[0]) | |
| display(i2i_images[0]) | |
| display(inpaint_images[0]) | |
| ``` | |
| In the above code, the `prompt2` is appended to the `prompt`, which is more than 77 tokens. "birds" are showing up in the result. | |
|  | |
| For more results, checkout [PR #6114](https://github.com/huggingface/diffusers/pull/6114). | |
| ### Example Images Mixing (with CoCa) | |
| ```python | |
| import requests | |
| from io import BytesIO | |
| import PIL | |
| import torch | |
| import open_clip | |
| from open_clip import SimpleTokenizer | |
| from diffusers import DiffusionPipeline | |
| from transformers import CLIPFeatureExtractor, CLIPModel | |
| def download_image(url): | |
| response = requests.get(url) | |
| return PIL.Image.open(BytesIO(response.content)).convert("RGB") | |
| # Loading additional models | |
| feature_extractor = CLIPFeatureExtractor.from_pretrained( | |
| "laion/CLIP-ViT-B-32-laion2B-s34B-b79K" | |
| ) | |
| clip_model = CLIPModel.from_pretrained( | |
| "laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16 | |
| ) | |
| coca_model = open_clip.create_model('coca_ViT-L-14', pretrained='laion2B-s13B-b90k').to('cuda') | |
| coca_model.dtype = torch.float16 | |
| coca_transform = open_clip.image_transform( | |
| coca_model.visual.image_size, | |
| is_train = False, | |
| mean = getattr(coca_model.visual, 'image_mean', None), | |
| std = getattr(coca_model.visual, 'image_std', None), | |
| ) | |
| coca_tokenizer = SimpleTokenizer() | |
| # Pipeline creating | |
| mixing_pipeline = DiffusionPipeline.from_pretrained( | |
| "CompVis/stable-diffusion-v1-4", | |
| custom_pipeline="clip_guided_images_mixing_stable_diffusion", | |
| clip_model=clip_model, | |
| feature_extractor=feature_extractor, | |
| coca_model=coca_model, | |
| coca_tokenizer=coca_tokenizer, | |
| coca_transform=coca_transform, | |
| torch_dtype=torch.float16, | |
| ) | |
| mixing_pipeline.enable_attention_slicing() | |
| mixing_pipeline = mixing_pipeline.to("cuda") | |
| # Pipeline running | |
| generator = torch.Generator(device="cuda").manual_seed(17) | |
| def download_image(url): | |
| response = requests.get(url) | |
| return PIL.Image.open(BytesIO(response.content)).convert("RGB") | |
| content_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir.jpg") | |
| style_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/gigachad.jpg") | |
| pipe_images = mixing_pipeline( | |
| num_inference_steps=50, | |
| content_image=content_image, | |
| style_image=style_image, | |
| noise_strength=0.65, | |
| slerp_latent_style_strength=0.9, | |
| slerp_prompt_style_strength=0.1, | |
| slerp_clip_image_style_strength=0.1, | |
| guidance_scale=9.0, | |
| batch_size=1, | |
| clip_guidance_scale=100, | |
| generator=generator, | |
| ).images | |
| ``` | |
|  | |
| ### Stable Diffusion Mixture Tiling | |
| This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details. | |
| ```python | |
| from diffusers import LMSDiscreteScheduler, DiffusionPipeline | |
| # Creater scheduler and model (similar to StableDiffusionPipeline) | |
| scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) | |
| pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler, custom_pipeline="mixture_tiling") | |
| pipeline.to("cuda") | |
| # Mixture of Diffusers generation | |
| image = pipeline( | |
| prompt=[[ | |
| "A charming house in the countryside, by jakub rozalski, sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece", | |
| "A dirt road in the countryside crossing pastures, by jakub rozalski, sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece", | |
| "An old and rusty giant robot lying on a dirt road, by jakub rozalski, dark sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece" | |
| ]], | |
| tile_height=640, | |
| tile_width=640, | |
| tile_row_overlap=0, | |
| tile_col_overlap=256, | |
| guidance_scale=8, | |
| seed=7178915308, | |
| num_inference_steps=50, | |
| )["images"][0] | |
| ``` | |
|  | |
| ### TensorRT Inpainting Stable Diffusion Pipeline | |
| The TensorRT Pipeline can be used to accelerate the Inpainting Stable Diffusion Inference run. | |
| NOTE: The ONNX conversions and TensorRT engine build may take up to 30 minutes. | |
| ```python | |
| import requests | |
| from io import BytesIO | |
| from PIL import Image | |
| import torch | |
| from diffusers import PNDMScheduler | |
| from diffusers.pipelines.stable_diffusion import StableDiffusionInpaintPipeline | |
| # Use the PNDMScheduler scheduler here instead | |
| scheduler = PNDMScheduler.from_pretrained("stabilityai/stable-diffusion-2-inpainting", subfolder="scheduler") | |
| pipe = StableDiffusionInpaintPipeline.from_pretrained("stabilityai/stable-diffusion-2-inpainting", | |
| custom_pipeline="stable_diffusion_tensorrt_inpaint", | |
| revision='fp16', | |
| torch_dtype=torch.float16, | |
| scheduler=scheduler, | |
| ) | |
| # re-use cached folder to save ONNX models and TensorRT Engines | |
| pipe.set_cached_folder("stabilityai/stable-diffusion-2-inpainting", revision='fp16',) | |
| pipe = pipe.to("cuda") | |
| url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
| response = requests.get(url) | |
| input_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
| response = requests.get(mask_url) | |
| mask_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| prompt = "a mecha robot sitting on a bench" | |
| image = pipe(prompt, image=input_image, mask_image=mask_image, strength=0.75,).images[0] | |
| image.save('tensorrt_inpaint_mecha_robot.png') | |
| ``` | |
| ### Stable Diffusion Mixture Canvas | |
| This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details. | |
| ```python | |
| from PIL import Image | |
| from diffusers import LMSDiscreteScheduler, DiffusionPipeline | |
| from diffusers.pipelines.pipeline_utils import Image2ImageRegion, Text2ImageRegion, preprocess_image | |
| # Load and preprocess guide image | |
| iic_image = preprocess_image(Image.open("input_image.png").convert("RGB")) | |
| # Creater scheduler and model (similar to StableDiffusionPipeline) | |
| scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) | |
| pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler).to("cuda:0", custom_pipeline="mixture_canvas") | |
| pipeline.to("cuda") | |
| # Mixture of Diffusers generation | |
| output = pipeline( | |
| canvas_height=800, | |
| canvas_width=352, | |
| regions=[ | |
| Text2ImageRegion(0, 800, 0, 352, guidance_scale=8, | |
| prompt=f"best quality, masterpiece, WLOP, sakimichan, art contest winner on pixiv, 8K, intricate details, wet effects, rain drops, ethereal, mysterious, futuristic, UHD, HDR, cinematic lighting, in a beautiful forest, rainy day, award winning, trending on artstation, beautiful confident cheerful young woman, wearing a futuristic sleeveless dress, ultra beautiful detailed eyes, hyper-detailed face, complex, perfect, model, textured, chiaroscuro, professional make-up, realistic, figure in frame, "), | |
| Image2ImageRegion(352-800, 352, 0, 352, reference_image=iic_image, strength=1.0), | |
| ], | |
| num_inference_steps=100, | |
| seed=5525475061, | |
| )["images"][0] | |
| ``` | |
|  | |
|  | |
| ### IADB pipeline | |
| This pipeline is the implementation of the [α-(de)Blending: a Minimalist Deterministic Diffusion Model](https://arxiv.org/abs/2305.03486) paper. | |
| It is a simple and minimalist diffusion model. | |
| The following code shows how to use the IADB pipeline to generate images using a pretrained celebahq-256 model. | |
| ```python | |
| pipeline_iadb = DiffusionPipeline.from_pretrained("thomasc4/iadb-celebahq-256", custom_pipeline='iadb') | |
| pipeline_iadb = pipeline_iadb.to('cuda') | |
| output = pipeline_iadb(batch_size=4,num_inference_steps=128) | |
| for i in range(len(output[0])): | |
| plt.imshow(output[0][i]) | |
| plt.show() | |
| ``` | |
| Sampling with the IADB formulation is easy, and can be done in a few lines (the pipeline already implements it): | |
| ```python | |
| def sample_iadb(model, x0, nb_step): | |
| x_alpha = x0 | |
| for t in range(nb_step): | |
| alpha = (t/nb_step) | |
| alpha_next =((t+1)/nb_step) | |
| d = model(x_alpha, torch.tensor(alpha, device=x_alpha.device))['sample'] | |
| x_alpha = x_alpha + (alpha_next-alpha)*d | |
| return x_alpha | |
| ``` | |
| The training loop is also straightforward: | |
| ```python | |
| # Training loop | |
| while True: | |
| x0 = sample_noise() | |
| x1 = sample_dataset() | |
| alpha = torch.rand(batch_size) | |
| # Blend | |
| x_alpha = (1-alpha) * x0 + alpha * x1 | |
| # Loss | |
| loss = torch.sum((D(x_alpha, alpha)- (x1-x0))**2) | |
| optimizer.zero_grad() | |
| loss.backward() | |
| optimizer.step() | |
| ``` | |
| ### Zero1to3 pipeline | |
| This pipeline is the implementation of the [Zero-1-to-3: Zero-shot One Image to 3D Object](https://arxiv.org/abs/2303.11328) paper. | |
| The original pytorch-lightning [repo](https://github.com/cvlab-columbia/zero123) and a diffusers [repo](https://github.com/kxhit/zero123-hf). | |
| The following code shows how to use the Zero1to3 pipeline to generate novel view synthesis images using a pretrained stable diffusion model. | |
| ```python | |
| import os | |
| import torch | |
| from pipeline_zero1to3 import Zero1to3StableDiffusionPipeline | |
| from diffusers.utils import load_image | |
| model_id = "kxic/zero123-165000" # zero123-105000, zero123-165000, zero123-xl | |
| pipe = Zero1to3StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16) | |
| pipe.enable_xformers_memory_efficient_attention() | |
| pipe.enable_vae_tiling() | |
| pipe.enable_attention_slicing() | |
| pipe = pipe.to("cuda") | |
| num_images_per_prompt = 4 | |
| # test inference pipeline | |
| # x y z, Polar angle (vertical rotation in degrees) Azimuth angle (horizontal rotation in degrees) Zoom (relative distance from center) | |
| query_pose1 = [-75.0, 100.0, 0.0] | |
| query_pose2 = [-20.0, 125.0, 0.0] | |
| query_pose3 = [-55.0, 90.0, 0.0] | |
| # load image | |
| # H, W = (256, 256) # H, W = (512, 512) # zero123 training is 256,256 | |
| # for batch input | |
| input_image1 = load_image("./demo/4_blackarm.png") #load_image("https://cvlab-zero123-live.hf.space/file=/home/user/app/configs/4_blackarm.png") | |
| input_image2 = load_image("./demo/8_motor.png") #load_image("https://cvlab-zero123-live.hf.space/file=/home/user/app/configs/8_motor.png") | |
| input_image3 = load_image("./demo/7_london.png") #load_image("https://cvlab-zero123-live.hf.space/file=/home/user/app/configs/7_london.png") | |
| input_images = [input_image1, input_image2, input_image3] | |
| query_poses = [query_pose1, query_pose2, query_pose3] | |
| # # for single input | |
| # H, W = (256, 256) | |
| # input_images = [input_image2.resize((H, W), PIL.Image.NEAREST)] | |
| # query_poses = [query_pose2] | |
| # better do preprocessing | |
| from gradio_new import preprocess_image, create_carvekit_interface | |
| import numpy as np | |
| import PIL.Image as Image | |
| pre_images = [] | |
| models = dict() | |
| print('Instantiating Carvekit HiInterface...') | |
| models['carvekit'] = create_carvekit_interface() | |
| if not isinstance(input_images, list): | |
| input_images = [input_images] | |
| for raw_im in input_images: | |
| input_im = preprocess_image(models, raw_im, True) | |
| H, W = input_im.shape[:2] | |
| pre_images.append(Image.fromarray((input_im * 255.0).astype(np.uint8))) | |
| input_images = pre_images | |
| # infer pipeline, in original zero123 num_inference_steps=76 | |
| images = pipe(input_imgs=input_images, prompt_imgs=input_images, poses=query_poses, height=H, width=W, | |
| guidance_scale=3.0, num_images_per_prompt=num_images_per_prompt, num_inference_steps=50).images | |
| # save imgs | |
| log_dir = "logs" | |
| os.makedirs(log_dir, exist_ok=True) | |
| bs = len(input_images) | |
| i = 0 | |
| for obj in range(bs): | |
| for idx in range(num_images_per_prompt): | |
| images[i].save(os.path.join(log_dir,f"obj{obj}_{idx}.jpg")) | |
| i += 1 | |
| ``` | |
| ### Stable Diffusion XL Reference | |
| This pipeline uses the Reference . Refer to the [stable_diffusion_reference](https://github.com/huggingface/diffusers/blob/main/examples/community/README.md#stable-diffusion-reference). | |
| ```py | |
| import torch | |
| from PIL import Image | |
| from diffusers.utils import load_image | |
| from diffusers import DiffusionPipeline | |
| from diffusers.schedulers import UniPCMultistepScheduler | |
| input_image = load_image("https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png") | |
| # pipe = DiffusionPipeline.from_pretrained( | |
| # "stabilityai/stable-diffusion-xl-base-1.0", | |
| # custom_pipeline="stable_diffusion_xl_reference", | |
| # torch_dtype=torch.float16, | |
| # use_safetensors=True, | |
| # variant="fp16").to('cuda:0') | |
| pipe = StableDiffusionXLReferencePipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-xl-base-1.0", | |
| torch_dtype=torch.float16, | |
| use_safetensors=True, | |
| variant="fp16").to('cuda:0') | |
| pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) | |
| result_img = pipe(ref_image=input_image, | |
| prompt="1girl", | |
| num_inference_steps=20, | |
| reference_attn=True, | |
| reference_adain=True).images[0] | |
| ``` | |
| Reference Image | |
|  | |
| Output Image | |
| `prompt: 1 girl` | |
| `reference_attn=True, reference_adain=True, num_inference_steps=20` | |
|  | |
| Reference Image | |
|  | |
| Output Image | |
| `prompt: A dog` | |
| `reference_attn=True, reference_adain=False, num_inference_steps=20` | |
|  | |
| Reference Image | |
|  | |
| Output Image | |
| `prompt: An astronaut riding a lion` | |
| `reference_attn=True, reference_adain=True, num_inference_steps=20` | |
|  | |
| ### Stable diffusion fabric pipeline | |
| FABRIC approach applicable to a wide range of popular diffusion models, which exploits | |
| the self-attention layer present in the most widely used architectures to condition | |
| the diffusion process on a set of feedback images. | |
| ```python | |
| import requests | |
| import torch | |
| from PIL import Image | |
| from io import BytesIO | |
| from diffusers import DiffusionPipeline | |
| # load the pipeline | |
| # make sure you're logged in with `huggingface-cli login` | |
| model_id_or_path = "runwayml/stable-diffusion-v1-5" | |
| #can also be used with dreamlike-art/dreamlike-photoreal-2.0 | |
| pipe = DiffusionPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16, custom_pipeline="pipeline_fabric").to("cuda") | |
| # let's specify a prompt | |
| prompt = "An astronaut riding an elephant" | |
| negative_prompt = "lowres, cropped" | |
| # call the pipeline | |
| image = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| num_inference_steps=20, | |
| generator=torch.manual_seed(12) | |
| ).images[0] | |
| image.save("horse_to_elephant.jpg") | |
| # let's try another example with feedback | |
| url = "https://raw.githubusercontent.com/ChenWu98/cycle-diffusion/main/data/dalle2/A%20black%20colored%20car.png" | |
| response = requests.get(url) | |
| init_image = Image.open(BytesIO(response.content)).convert("RGB") | |
| prompt = "photo, A blue colored car, fish eye" | |
| liked = [init_image] | |
| ## same goes with disliked | |
| # call the pipeline | |
| torch.manual_seed(0) | |
| image = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| liked = liked, | |
| num_inference_steps=20, | |
| ).images[0] | |
| image.save("black_to_blue.png") | |
| ``` | |
| *With enough feedbacks you can create very similar high quality images.* | |
| The original codebase can be found at [sd-fabric/fabric](https://github.com/sd-fabric/fabric), and available checkpoints are [dreamlike-art/dreamlike-photoreal-2.0](https://huggingface.co/dreamlike-art/dreamlike-photoreal-2.0), [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5), and [stabilityai/stable-diffusion-2-1](https://huggingface.co/stabilityai/stable-diffusion-2-1) (may give unexpected results). | |
| Let's have a look at the images (_512X512_) | |
| | Without Feedback | With Feedback (1st image) | | |
| |---------------------|---------------------| | |
| |  |  | | |
| ### Masked Im2Im Stable Diffusion Pipeline | |
| This pipeline reimplements sketch inpaint feature from A1111 for non-inpaint models. The following code reads two images, original and one with mask painted over it. It computes mask as a difference of two images and does the inpainting in the area defined by the mask. | |
| ```python | |
| img = PIL.Image.open("./mech.png") | |
| # read image with mask painted over | |
| img_paint = PIL.Image.open("./mech_painted.png") | |
| neq = numpy.any(numpy.array(img) != numpy.array(img_paint), axis=-1) | |
| mask = neq / neq.max() | |
| pipeline = MaskedStableDiffusionImg2ImgPipeline.from_pretrained("frankjoshua/icbinpICantBelieveIts_v8") | |
| # works best with EulerAncestralDiscreteScheduler | |
| pipeline.scheduler = EulerAncestralDiscreteScheduler.from_config(pipeline.scheduler.config) | |
| generator = torch.Generator(device="cpu").manual_seed(4) | |
| prompt = "a man wearing a mask" | |
| result = pipeline(prompt=prompt, image=img_paint, mask=mask, strength=0.75, | |
| generator=generator) | |
| result.images[0].save("result.png") | |
| ``` | |
| original image mech.png | |
| <img src=<https://github.com/noskill/diffusers/assets/733626/10ad972d-d655-43cb-8de1-039e3d79e849> width="25%" > | |
| image with mask mech_painted.png | |
| <img src=<https://github.com/noskill/diffusers/assets/733626/c334466a-67fe-4377-9ff7-f46021b9c224> width="25%" > | |
| result: | |
| <img src=<https://github.com/noskill/diffusers/assets/733626/23a0a71d-51db-471e-926a-107ac62512a8> width="25%" > | |
| ### Prompt2Prompt Pipeline | |
| Prompt2Prompt allows the following edits: | |
| - ReplaceEdit (change words in prompt) | |
| - ReplaceEdit with local blend (change words in prompt, keep image part unrelated to changes constant) | |
| - RefineEdit (add words to prompt) | |
| - RefineEdit with local blend (add words to prompt, keep image part unrelated to changes constant) | |
| - ReweightEdit (modulate importance of words) | |
| Here's a full example for `ReplaceEdit``: | |
| ```python | |
| import torch | |
| import numpy as np | |
| import matplotlib.pyplot as plt | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="pipeline_prompt2prompt").to("cuda") | |
| prompts = ["A turtle playing with a ball", | |
| "A monkey playing with a ball"] | |
| cross_attention_kwargs = { | |
| "edit_type": "replace", | |
| "cross_replace_steps": 0.4, | |
| "self_replace_steps": 0.4 | |
| } | |
| outputs = pipe(prompt=prompts, height=512, width=512, num_inference_steps=50, cross_attention_kwargs=cross_attention_kwargs) | |
| ``` | |
| And abbreviated examples for the other edits: | |
| `ReplaceEdit with local blend` | |
| ```python | |
| prompts = ["A turtle playing with a ball", | |
| "A monkey playing with a ball"] | |
| cross_attention_kwargs = { | |
| "edit_type": "replace", | |
| "cross_replace_steps": 0.4, | |
| "self_replace_steps": 0.4, | |
| "local_blend_words": ["turtle", "monkey"] | |
| } | |
| ``` | |
| `RefineEdit` | |
| ```python | |
| prompts = ["A turtle", | |
| "A turtle in a forest"] | |
| cross_attention_kwargs = { | |
| "edit_type": "refine", | |
| "cross_replace_steps": 0.4, | |
| "self_replace_steps": 0.4, | |
| } | |
| ``` | |
| `RefineEdit with local blend` | |
| ```python | |
| prompts = ["A turtle", | |
| "A turtle in a forest"] | |
| cross_attention_kwargs = { | |
| "edit_type": "refine", | |
| "cross_replace_steps": 0.4, | |
| "self_replace_steps": 0.4, | |
| "local_blend_words": ["in", "a" , "forest"] | |
| } | |
| ``` | |
| `ReweightEdit` | |
| ```python | |
| prompts = ["A smiling turtle"] * 2 | |
| edit_kcross_attention_kwargswargs = { | |
| "edit_type": "reweight", | |
| "cross_replace_steps": 0.4, | |
| "self_replace_steps": 0.4, | |
| "equalizer_words": ["smiling"], | |
| "equalizer_strengths": [5] | |
| } | |
| ``` | |
| Side note: See [this GitHub gist](https://gist.github.com/UmerHA/b65bb5fb9626c9c73f3ade2869e36164) if you want to visualize the attention maps. | |
| ### Latent Consistency Pipeline | |
| Latent Consistency Models was proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://arxiv.org/abs/2310.04378) by _Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, Hang Zhao_ from Tsinghua University. | |
| The abstract of the paper reads as follows: | |
| *Latent Diffusion models (LDMs) have achieved remarkable results in synthesizing high-resolution images. However, the iterative sampling process is computationally intensive and leads to slow generation. Inspired by Consistency Models (song et al.), we propose Latent Consistency Models (LCMs), enabling swift inference with minimal steps on any pre-trained LDMs, including Stable Diffusion (rombach et al). Viewing the guided reverse diffusion process as solving an augmented probability flow ODE (PF-ODE), LCMs are designed to directly predict the solution of such ODE in latent space, mitigating the need for numerous iterations and allowing rapid, high-fidelity sampling. Efficiently distilled from pre-trained classifier-free guided diffusion models, a high-quality 768 x 768 2~4-step LCM takes only 32 A100 GPU hours for training. Furthermore, we introduce Latent Consistency Fine-tuning (LCF), a novel method that is tailored for fine-tuning LCMs on customized image datasets. Evaluation on the LAION-5B-Aesthetics dataset demonstrates that LCMs achieve state-of-the-art text-to-image generation performance with few-step inference. Project Page: [this https URL](https://latent-consistency-models.github.io/)* | |
| The model can be used with `diffusers` as follows: | |
| - *1. Load the model from the community pipeline.* | |
| ```py | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained("SimianLuo/LCM_Dreamshaper_v7", custom_pipeline="latent_consistency_txt2img", custom_revision="main") | |
| # To save GPU memory, torch.float16 can be used, but it may compromise image quality. | |
| pipe.to(torch_device="cuda", torch_dtype=torch.float32) | |
| ``` | |
| - 2. Run inference with as little as 4 steps: | |
| ```py | |
| prompt = "Self-portrait oil painting, a beautiful cyborg with golden hair, 8k" | |
| # Can be set to 1~50 steps. LCM support fast inference even <= 4 steps. Recommend: 1~8 steps. | |
| num_inference_steps = 4 | |
| images = pipe(prompt=prompt, num_inference_steps=num_inference_steps, guidance_scale=8.0, lcm_origin_steps=50, output_type="pil").images | |
| ``` | |
| For any questions or feedback, feel free to reach out to [Simian Luo](https://github.com/luosiallen). | |
| You can also try this pipeline directly in the [🚀 official spaces](https://huggingface.co/spaces/SimianLuo/Latent_Consistency_Model). | |
| ### Latent Consistency Img2img Pipeline | |
| This pipeline extends the Latent Consistency Pipeline to allow it to take an input image. | |
| ```py | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained("SimianLuo/LCM_Dreamshaper_v7", custom_pipeline="latent_consistency_img2img") | |
| # To save GPU memory, torch.float16 can be used, but it may compromise image quality. | |
| pipe.to(torch_device="cuda", torch_dtype=torch.float32) | |
| ``` | |
| - 2. Run inference with as little as 4 steps: | |
| ```py | |
| prompt = "Self-portrait oil painting, a beautiful cyborg with golden hair, 8k" | |
| input_image=Image.open("myimg.png") | |
| strength = 0.5 #strength =0 (no change) strength=1 (completely overwrite image) | |
| # Can be set to 1~50 steps. LCM support fast inference even <= 4 steps. Recommend: 1~8 steps. | |
| num_inference_steps = 4 | |
| images = pipe(prompt=prompt, image=input_image, strength=strength, num_inference_steps=num_inference_steps, guidance_scale=8.0, lcm_origin_steps=50, output_type="pil").images | |
| ``` | |
| ### Latent Consistency Interpolation Pipeline | |
| This pipeline extends the Latent Consistency Pipeline to allow for interpolation of the latent space between multiple prompts. It is similar to the [Stable Diffusion Interpolate](https://github.com/huggingface/diffusers/blob/main/examples/community/interpolate_stable_diffusion.py) and [unCLIP Interpolate](https://github.com/huggingface/diffusers/blob/main/examples/community/unclip_text_interpolation.py) community pipelines. | |
| ```py | |
| import torch | |
| import numpy as np | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained("SimianLuo/LCM_Dreamshaper_v7", custom_pipeline="latent_consistency_interpolate") | |
| # To save GPU memory, torch.float16 can be used, but it may compromise image quality. | |
| pipe.to(torch_device="cuda", torch_dtype=torch.float32) | |
| prompts = [ | |
| "Self-portrait oil painting, a beautiful cyborg with golden hair, Margot Robbie, 8k", | |
| "Self-portrait oil painting, an extremely strong man, body builder, Huge Jackman, 8k", | |
| "An astronaut floating in space, renaissance art, realistic, high quality, 8k", | |
| "Oil painting of a cat, cute, dream-like", | |
| "Hugging face emoji, cute, realistic" | |
| ] | |
| num_inference_steps = 4 | |
| num_interpolation_steps = 60 | |
| seed = 1337 | |
| torch.manual_seed(seed) | |
| np.random.seed(seed) | |
| images = pipe( | |
| prompt=prompts, | |
| height=512, | |
| width=512, | |
| num_inference_steps=num_inference_steps, | |
| num_interpolation_steps=num_interpolation_steps, | |
| guidance_scale=8.0, | |
| embedding_interpolation_type="lerp", | |
| latent_interpolation_type="slerp", | |
| process_batch_size=4, # Make it higher or lower based on your GPU memory | |
| generator=torch.Generator(seed), | |
| ) | |
| assert len(images) == (len(prompts) - 1) * num_interpolation_steps | |
| ``` | |
| ### StableDiffusionUpscaleLDM3D Pipeline | |
| [LDM3D-VR](https://arxiv.org/pdf/2311.03226.pdf) is an extended version of LDM3D. | |
| The abstract from the paper is: | |
| *Latent diffusion models have proven to be state-of-the-art in the creation and manipulation of visual outputs. However, as far as we know, the generation of depth maps jointly with RGB is still limited. We introduce LDM3D-VR, a suite of diffusion models targeting virtual reality development that includes LDM3D-pano and LDM3D-SR. These models enable the generation of panoramic RGBD based on textual prompts and the upscaling of low-resolution inputs to high-resolution RGBD, respectively. Our models are fine-tuned from existing pretrained models on datasets containing panoramic/high-resolution RGB images, depth maps and captions. Both models are evaluated in comparison to existing related methods* | |
| Two checkpoints are available for use: | |
| - [ldm3d-pano](https://huggingface.co/Intel/ldm3d-pano). This checkpoint enables the generation of panoramic images and requires the StableDiffusionLDM3DPipeline pipeline to be used. | |
| - [ldm3d-sr](https://huggingface.co/Intel/ldm3d-sr). This checkpoint enables the upscaling of RGB and depth images. Can be used in cascade after the original LDM3D pipeline using the StableDiffusionUpscaleLDM3DPipeline pipeline. | |
| '''py | |
| from PIL import Image | |
| import os | |
| import torch | |
| from diffusers import StableDiffusionLDM3DPipeline, DiffusionPipeline | |
| # Generate a rgb/depth output from LDM3D | |
| pipe_ldm3d = StableDiffusionLDM3DPipeline.from_pretrained("Intel/ldm3d-4c") | |
| pipe_ldm3d.to("cuda") | |
| prompt =f"A picture of some lemons on a table" | |
| output = pipe_ldm3d(prompt) | |
| rgb_image, depth_image = output.rgb, output.depth | |
| rgb_image[0].save(f"lemons_ldm3d_rgb.jpg") | |
| depth_image[0].save(f"lemons_ldm3d_depth.png") | |
| # Upscale the previous output to a resolution of (1024, 1024) | |
| pipe_ldm3d_upscale = DiffusionPipeline.from_pretrained("Intel/ldm3d-sr", custom_pipeline="pipeline_stable_diffusion_upscale_ldm3d") | |
| pipe_ldm3d_upscale.to("cuda") | |
| low_res_img = Image.open(f"lemons_ldm3d_rgb.jpg").convert("RGB") | |
| low_res_depth = Image.open(f"lemons_ldm3d_depth.png").convert("L") | |
| outputs = pipe_ldm3d_upscale(prompt="high quality high resolution uhd 4k image", rgb=low_res_img, depth=low_res_depth, num_inference_steps=50, target_res=[1024, 1024]) | |
| upscaled_rgb, upscaled_depth =outputs.rgb[0], outputs.depth[0] | |
| upscaled_rgb.save(f"upscaled_lemons_rgb.png") | |
| upscaled_depth.save(f"upscaled_lemons_depth.png") | |
| ''' | |
| ### ControlNet + T2I Adapter Pipeline | |
| This pipelines combines both ControlNet and T2IAdapter into a single pipeline, where the forward pass is executed once. | |
| It receives `control_image` and `adapter_image`, as well as `controlnet_conditioning_scale` and `adapter_conditioning_scale`, for the ControlNet and Adapter modules, respectively. Whenever `adapter_conditioning_scale = 0` or `controlnet_conditioning_scale = 0`, it will act as a full ControlNet module or as a full T2IAdapter module, respectively. | |
| ```py | |
| import cv2 | |
| import numpy as np | |
| import torch | |
| from controlnet_aux.midas import MidasDetector | |
| from PIL import Image | |
| from diffusers import AutoencoderKL, ControlNetModel, MultiAdapter, T2IAdapter | |
| from diffusers.pipelines.controlnet.multicontrolnet import MultiControlNetModel | |
| from diffusers.utils import load_image | |
| from examples.community.pipeline_stable_diffusion_xl_controlnet_adapter import ( | |
| StableDiffusionXLControlNetAdapterPipeline, | |
| ) | |
| controlnet_depth = ControlNetModel.from_pretrained( | |
| "diffusers/controlnet-depth-sdxl-1.0", | |
| torch_dtype=torch.float16, | |
| variant="fp16", | |
| use_safetensors=True | |
| ) | |
| adapter_depth = T2IAdapter.from_pretrained( | |
| "TencentARC/t2i-adapter-depth-midas-sdxl-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True | |
| ) | |
| vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16, use_safetensors=True) | |
| pipe = StableDiffusionXLControlNetAdapterPipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-xl-base-1.0", | |
| controlnet=controlnet_depth, | |
| adapter=adapter_depth, | |
| vae=vae, | |
| variant="fp16", | |
| use_safetensors=True, | |
| torch_dtype=torch.float16, | |
| ) | |
| pipe = pipe.to("cuda") | |
| pipe.enable_xformers_memory_efficient_attention() | |
| # pipe.enable_freeu(s1=0.6, s2=0.4, b1=1.1, b2=1.2) | |
| midas_depth = MidasDetector.from_pretrained( | |
| "valhalla/t2iadapter-aux-models", filename="dpt_large_384.pt", model_type="dpt_large" | |
| ).to("cuda") | |
| prompt = "a tiger sitting on a park bench" | |
| img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
| image = load_image(img_url).resize((1024, 1024)) | |
| depth_image = midas_depth( | |
| image, detect_resolution=512, image_resolution=1024 | |
| ) | |
| strength = 0.5 | |
| images = pipe( | |
| prompt, | |
| control_image=depth_image, | |
| adapter_image=depth_image, | |
| num_inference_steps=30, | |
| controlnet_conditioning_scale=strength, | |
| adapter_conditioning_scale=strength, | |
| ).images | |
| images[0].save("controlnet_and_adapter.png") | |
| ``` | |
| ### ControlNet + T2I Adapter + Inpainting Pipeline | |
| ```py | |
| import cv2 | |
| import numpy as np | |
| import torch | |
| from controlnet_aux.midas import MidasDetector | |
| from PIL import Image | |
| from diffusers import AutoencoderKL, ControlNetModel, MultiAdapter, T2IAdapter | |
| from diffusers.pipelines.controlnet.multicontrolnet import MultiControlNetModel | |
| from diffusers.utils import load_image | |
| from examples.community.pipeline_stable_diffusion_xl_controlnet_adapter_inpaint import ( | |
| StableDiffusionXLControlNetAdapterInpaintPipeline, | |
| ) | |
| controlnet_depth = ControlNetModel.from_pretrained( | |
| "diffusers/controlnet-depth-sdxl-1.0", | |
| torch_dtype=torch.float16, | |
| variant="fp16", | |
| use_safetensors=True | |
| ) | |
| adapter_depth = T2IAdapter.from_pretrained( | |
| "TencentARC/t2i-adapter-depth-midas-sdxl-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True | |
| ) | |
| vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16, use_safetensors=True) | |
| pipe = StableDiffusionXLControlNetAdapterInpaintPipeline.from_pretrained( | |
| "diffusers/stable-diffusion-xl-1.0-inpainting-0.1", | |
| controlnet=controlnet_depth, | |
| adapter=adapter_depth, | |
| vae=vae, | |
| variant="fp16", | |
| use_safetensors=True, | |
| torch_dtype=torch.float16, | |
| ) | |
| pipe = pipe.to("cuda") | |
| pipe.enable_xformers_memory_efficient_attention() | |
| # pipe.enable_freeu(s1=0.6, s2=0.4, b1=1.1, b2=1.2) | |
| midas_depth = MidasDetector.from_pretrained( | |
| "valhalla/t2iadapter-aux-models", filename="dpt_large_384.pt", model_type="dpt_large" | |
| ).to("cuda") | |
| prompt = "a tiger sitting on a park bench" | |
| img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
| mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
| image = load_image(img_url).resize((1024, 1024)) | |
| mask_image = load_image(mask_url).resize((1024, 1024)) | |
| depth_image = midas_depth( | |
| image, detect_resolution=512, image_resolution=1024 | |
| ) | |
| strength = 0.4 | |
| images = pipe( | |
| prompt, | |
| image=image, | |
| mask_image=mask_image, | |
| control_image=depth_image, | |
| adapter_image=depth_image, | |
| num_inference_steps=30, | |
| controlnet_conditioning_scale=strength, | |
| adapter_conditioning_scale=strength, | |
| strength=0.7, | |
| ).images | |
| images[0].save("controlnet_and_adapter_inpaint.png") | |
| ``` | |
| ### Regional Prompting Pipeline | |
| This pipeline is a port of the [Regional Prompter extension](https://github.com/hako-mikan/sd-webui-regional-prompter) for [Stable Diffusion web UI](https://github.com/AUTOMATIC1111/stable-diffusion-webui) to diffusers. | |
| This code implements a pipeline for the Stable Diffusion model, enabling the division of the canvas into multiple regions, with different prompts applicable to each region. Users can specify regions in two ways: using `Cols` and `Rows` modes for grid-like divisions, or the `Prompt` mode for regions calculated based on prompts. | |
|  | |
| ### Usage | |
| ### Sample Code | |
| ```py | |
| from examples.community.regional_prompting_stable_diffusion import RegionalPromptingStableDiffusionPipeline | |
| pipe = RegionalPromptingStableDiffusionPipeline.from_single_file(model_path, vae=vae) | |
| rp_args = { | |
| "mode":"rows", | |
| "div": "1;1;1" | |
| } | |
| prompt =""" | |
| green hair twintail BREAK | |
| red blouse BREAK | |
| blue skirt | |
| """ | |
| images = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| guidance_scale=7.5, | |
| height = 768, | |
| width = 512, | |
| num_inference_steps =20, | |
| num_images_per_prompt = 1, | |
| rp_args = rp_args | |
| ).images | |
| time = time.strftime(r"%Y%m%d%H%M%S") | |
| i = 1 | |
| for image in images: | |
| i += 1 | |
| fileName = f'img-{time}-{i+1}.png' | |
| image.save(fileName) | |
| ``` | |
| ### Cols, Rows mode | |
| In the Cols, Rows mode, you can split the screen vertically and horizontally and assign prompts to each region. The split ratio can be specified by 'div', and you can set the division ratio like '3;3;2' or '0.1;0.5'. Furthermore, as will be described later, you can also subdivide the split Cols, Rows to specify more complex regions. | |
| In this image, the image is divided into three parts, and a separate prompt is applied to each. The prompts are divided by 'BREAK', and each is applied to the respective region. | |
|  | |
| ``` | |
| green hair twintail BREAK | |
| red blouse BREAK | |
| blue skirt | |
| ``` | |
| ### 2-Dimentional division | |
| The prompt consists of instructions separated by the term `BREAK` and is assigned to different regions of a two-dimensional space. The image is initially split in the main splitting direction, which in this case is rows, due to the presence of a single semicolon`;`, dividing the space into an upper and a lower section. Additional sub-splitting is then applied, indicated by commas. The upper row is split into ratios of `2:1:1`, while the lower row is split into a ratio of `4:6`. Rows themselves are split in a `1:2` ratio. According to the reference image, the blue sky is designated as the first region, green hair as the second, the bookshelf as the third, and so on, in a sequence based on their position from the top left. The terrarium is placed on the desk in the fourth region, and the orange dress and sofa are in the fifth region, conforming to their respective splits. | |
| ``` | |
| rp_args = { | |
| "mode":"rows", | |
| "div": "1,2,1,1;2,4,6" | |
| } | |
| prompt =""" | |
| blue sky BREAK | |
| green hair BREAK | |
| book shelf BREAK | |
| terrarium on desk BREAK | |
| orange dress and sofa | |
| """ | |
| ``` | |
|  | |
| ### Prompt Mode | |
| There are limitations to methods of specifying regions in advance. This is because specifying regions can be a hindrance when designating complex shapes or dynamic compositions. In the region specified by the prompt, the regions is determined after the image generation has begun. This allows us to accommodate compositions and complex regions. | |
| For further infomagen, see [here](https://github.com/hako-mikan/sd-webui-regional-prompter/blob/main/prompt_en.md). | |
| ### syntax | |
| ``` | |
| baseprompt target1 target2 BREAK | |
| effect1, target1 BREAK | |
| effect2 ,target2 | |
| ``` | |
| First, write the base prompt. In the base prompt, write the words (target1, target2) for which you want to create a mask. Next, separate them with BREAK. Next, write the prompt corresponding to target1. Then enter a comma and write target1. The order of the targets in the base prompt and the order of the BREAK-separated targets can be back to back. | |
| ``` | |
| target2 baseprompt target1 BREAK | |
| effect1, target1 BREAK | |
| effect2 ,target2 | |
| ``` | |
| is also effective. | |
| ### Sample | |
| In this example, masks are calculated for shirt, tie, skirt, and color prompts are specified only for those regions. | |
| ``` | |
| rp_args = { | |
| "mode":"prompt-ex", | |
| "save_mask":True, | |
| "th": "0.4,0.6,0.6", | |
| } | |
| prompt =""" | |
| a girl in street with shirt, tie, skirt BREAK | |
| red, shirt BREAK | |
| green, tie BREAK | |
| blue , skirt | |
| """ | |
| ``` | |
|  | |
| ### threshold | |
| The threshold used to determine the mask created by the prompt. This can be set as many times as there are masks, as the range varies widely depending on the target prompt. If multiple regions are used, enter them separated by commas. For example, hair tends to be ambiguous and requires a small value, while face tends to be large and requires a small value. These should be ordered by BREAK. | |
| ``` | |
| a lady ,hair, face BREAK | |
| red, hair BREAK | |
| tanned ,face | |
| ``` | |
| `threshold : 0.4,0.6` | |
| If only one input is given for multiple regions, they are all assumed to be the same value. | |
| ### Prompt and Prompt-EX | |
| The difference is that in Prompt, duplicate regions are added, whereas in Prompt-EX, duplicate regions are overwritten sequentially. Since they are processed in order, setting a TARGET with a large regions first makes it easier for the effect of small regions to remain unmuffled. | |
| ### Accuracy | |
| In the case of a 512 x 512 image, Attention mode reduces the size of the region to about 8 x 8 pixels deep in the U-Net, so that small regions get mixed up; Latent mode calculates 64*64, so that the region is exact. | |
| ``` | |
| girl hair twintail frills,ribbons, dress, face BREAK | |
| girl, ,face | |
| ``` | |
| ### Mask | |
| When an image is generated, the generated mask is displayed. It is generated at the same size as the image, but is actually used at a much smaller size. | |
| ### Use common prompt | |
| You can attach the prompt up to ADDCOMM to all prompts by separating it first with ADDCOMM. This is useful when you want to include elements common to all regions. For example, when generating pictures of three people with different appearances, it's necessary to include the instruction of 'three people' in all regions. It's also useful when inserting quality tags and other things."For example, if you write as follows: | |
| ``` | |
| best quality, 3persons in garden, ADDCOMM | |
| a girl white dress BREAK | |
| a boy blue shirt BREAK | |
| an old man red suit | |
| ``` | |
| If common is enabled, this prompt is converted to the following: | |
| ``` | |
| best quality, 3persons in garden, a girl white dress BREAK | |
| best quality, 3persons in garden, a boy blue shirt BREAK | |
| best quality, 3persons in garden, an old man red suit | |
| ``` | |
| ### Negative prompt | |
| Negative prompts are equally effective across all regions, but it is possible to set region-specific prompts for negative prompts as well. The number of BREAKs must be the same as the number of prompts. If the number of prompts does not match, the negative prompts will be used without being divided into regions. | |
| ### Parameters | |
| To activate Regional Prompter, it is necessary to enter settings in rp_args. The items that can be set are as follows. rp_args is a dictionary type. | |
| ### Input Parameters | |
| Parameters are specified through the `rp_arg`(dictionary type). | |
| ``` | |
| rp_args = { | |
| "mode":"rows", | |
| "div": "1;1;1" | |
| } | |
| pipe(prompt =prompt, rp_args = rp_args) | |
| ``` | |
| ### Required Parameters | |
| - `mode`: Specifies the method for defining regions. Choose from `Cols`, `Rows`, `Prompt` or `Prompt-Ex`. This parameter is case-insensitive. | |
| - `divide`: Used in `Cols` and `Rows` modes. Details on how to specify this are provided under the respective `Cols` and `Rows` sections. | |
| - `th`: Used in `Prompt` mode. The method of specification is detailed under the `Prompt` section. | |
| ### Optional Parameters | |
| - `save_mask`: In `Prompt` mode, choose whether to output the generated mask along with the image. The default is `False`. | |
| The Pipeline supports `compel` syntax. Input prompts using the `compel` structure will be automatically applied and processed. | |
| ### Diffusion Posterior Sampling Pipeline | |
| - Reference paper | |
| ``` | |
| @article{chung2022diffusion, | |
| title={Diffusion posterior sampling for general noisy inverse problems}, | |
| author={Chung, Hyungjin and Kim, Jeongsol and Mccann, Michael T and Klasky, Marc L and Ye, Jong Chul}, | |
| journal={arXiv preprint arXiv:2209.14687}, | |
| year={2022} | |
| } | |
| ``` | |
| - This pipeline allows zero-shot conditional sampling from the posterior distribution $p(x|y)$, given observation on $y$, unconditional generative model $p(x)$ and differentiable operator $y=f(x)$. | |
| - For example, $f(.)$ can be downsample operator, then $y$ is a downsampled image, and the pipeline becomes a super-resolution pipeline. | |
| - To use this pipeline, you need to know your operator $f(.)$ and corrupted image $y$, and pass them during the call. For example, as in the main function of dps_pipeline.py, you need to first define the Gaussian blurring operator $f(.)$. The operator should be a callable nn.Module, with all the parameter gradient disabled: | |
| ```python | |
| import torch.nn.functional as F | |
| import scipy | |
| from torch import nn | |
| # define the Gaussian blurring operator first | |
| class GaussialBlurOperator(nn.Module): | |
| def __init__(self, kernel_size, intensity): | |
| super().__init__() | |
| class Blurkernel(nn.Module): | |
| def __init__(self, blur_type='gaussian', kernel_size=31, std=3.0): | |
| super().__init__() | |
| self.blur_type = blur_type | |
| self.kernel_size = kernel_size | |
| self.std = std | |
| self.seq = nn.Sequential( | |
| nn.ReflectionPad2d(self.kernel_size//2), | |
| nn.Conv2d(3, 3, self.kernel_size, stride=1, padding=0, bias=False, groups=3) | |
| ) | |
| self.weights_init() | |
| def forward(self, x): | |
| return self.seq(x) | |
| def weights_init(self): | |
| if self.blur_type == "gaussian": | |
| n = np.zeros((self.kernel_size, self.kernel_size)) | |
| n[self.kernel_size // 2,self.kernel_size // 2] = 1 | |
| k = scipy.ndimage.gaussian_filter(n, sigma=self.std) | |
| k = torch.from_numpy(k) | |
| self.k = k | |
| for name, f in self.named_parameters(): | |
| f.data.copy_(k) | |
| elif self.blur_type == "motion": | |
| k = Kernel(size=(self.kernel_size, self.kernel_size), intensity=self.std).kernelMatrix | |
| k = torch.from_numpy(k) | |
| self.k = k | |
| for name, f in self.named_parameters(): | |
| f.data.copy_(k) | |
| def update_weights(self, k): | |
| if not torch.is_tensor(k): | |
| k = torch.from_numpy(k) | |
| for name, f in self.named_parameters(): | |
| f.data.copy_(k) | |
| def get_kernel(self): | |
| return self.k | |
| self.kernel_size = kernel_size | |
| self.conv = Blurkernel(blur_type='gaussian', | |
| kernel_size=kernel_size, | |
| std=intensity) | |
| self.kernel = self.conv.get_kernel() | |
| self.conv.update_weights(self.kernel.type(torch.float32)) | |
| for param in self.parameters(): | |
| param.requires_grad=False | |
| def forward(self, data, **kwargs): | |
| return self.conv(data) | |
| def transpose(self, data, **kwargs): | |
| return data | |
| def get_kernel(self): | |
| return self.kernel.view(1, 1, self.kernel_size, self.kernel_size) | |
| ``` | |
| - Next, you should obtain the corrupted image $y$ by the operator. In this example, we generate $y$ from the source image $x$. However in practice, having the operator $f(.)$ and corrupted image $y$ is enough: | |
| ```python | |
| # set up source image | |
| src = Image.open('sample.png') | |
| # read image into [1,3,H,W] | |
| src = torch.from_numpy(np.array(src, dtype=np.float32)).permute(2,0,1)[None] | |
| # normalize image to [-1,1] | |
| src = (src / 127.5) - 1.0 | |
| src = src.to("cuda") | |
| # set up operator and measurement | |
| operator = GaussialBlurOperator(kernel_size=61, intensity=3.0).to("cuda") | |
| measurement = operator(src) | |
| # save the source and corrupted images | |
| save_image((src+1.0)/2.0, "dps_src.png") | |
| save_image((measurement+1.0)/2.0, "dps_mea.png") | |
| ``` | |
| - We provide an example pair of saved source and corrupted images, using the Gaussian blur operator above | |
| - Source image: | |
| -  | |
| - Gaussian blurred image: | |
| -  | |
| - You can download those image to run the example on your own. | |
| - Next, we need to define a loss function used for diffusion posterior sample. For most of the cases, the RMSE is fine: | |
| ```python | |
| def RMSELoss(yhat, y): | |
| return torch.sqrt(torch.sum((yhat-y)**2)) | |
| ``` | |
| - And next, as any other diffusion models, we need the score estimator and scheduler. As we are working with $256x256$ face images, we use ddmp-celebahq-256: | |
| ```python | |
| # set up scheduler | |
| scheduler = DDPMScheduler.from_pretrained("google/ddpm-celebahq-256") | |
| scheduler.set_timesteps(1000) | |
| # set up model | |
| model = UNet2DModel.from_pretrained("google/ddpm-celebahq-256").to("cuda") | |
| ``` | |
| - And finally, run the pipeline: | |
| ```python | |
| # finally, the pipeline | |
| dpspipe = DPSPipeline(model, scheduler) | |
| image = dpspipe( | |
| measurement = measurement, | |
| operator = operator, | |
| loss_fn = RMSELoss, | |
| zeta = 1.0, | |
| ).images[0] | |
| image.save("dps_generated_image.png") | |
| ``` | |
| - The zeta is a hyperparameter that is in range of $[0,1]$. It need to be tuned for best effect. By setting zeta=1, you should be able to have the reconstructed result: | |
| - Reconstructed image: | |
| -  | |
| - The reconstruction is perceptually similar to the source image, but different in details. | |
| - In dps_pipeline.py, we also provide a super-resolution example, which should produce: | |
| - Downsampled image: | |
| -  | |
| - Reconstructed image: | |
| -  | |
| ### AnimateDiff ControlNet Pipeline | |
| This pipeline combines AnimateDiff and ControlNet. Enjoy precise motion control for your videos! Refer to [this](https://github.com/huggingface/diffusers/issues/5866) issue for more details. | |
| ```py | |
| import torch | |
| from diffusers import AutoencoderKL, ControlNetModel, MotionAdapter | |
| from diffusers.pipelines import DiffusionPipeline | |
| from diffusers.schedulers import DPMSolverMultistepScheduler | |
| from PIL import Image | |
| motion_id = "guoyww/animatediff-motion-adapter-v1-5-2" | |
| adapter = MotionAdapter.from_pretrained(motion_id) | |
| controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_openpose", torch_dtype=torch.float16) | |
| vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16) | |
| model_id = "SG161222/Realistic_Vision_V5.1_noVAE" | |
| pipe = DiffusionPipeline.from_pretrained( | |
| model_id, | |
| motion_adapter=adapter, | |
| controlnet=controlnet, | |
| vae=vae, | |
| custom_pipeline="pipeline_animatediff_controlnet", | |
| ).to(device="cuda", dtype=torch.float16) | |
| pipe.scheduler = DPMSolverMultistepScheduler.from_pretrained( | |
| model_id, subfolder="scheduler", beta_schedule="linear", clip_sample=False, timestep_spacing="linspace", steps_offset=1 | |
| ) | |
| pipe.enable_vae_slicing() | |
| conditioning_frames = [] | |
| for i in range(1, 16 + 1): | |
| conditioning_frames.append(Image.open(f"frame_{i}.png")) | |
| prompt = "astronaut in space, dancing" | |
| negative_prompt = "bad quality, worst quality, jpeg artifacts, ugly" | |
| result = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| width=512, | |
| height=768, | |
| conditioning_frames=conditioning_frames, | |
| num_inference_steps=20, | |
| ).frames[0] | |
| from diffusers.utils import export_to_gif | |
| export_to_gif(result.frames[0], "result.gif") | |
| ``` | |
| <table> | |
| <tr><td colspan="2" align=center><b>Conditioning Frames</b></td></tr> | |
| <tr align=center> | |
| <td align=center><img src="https://user-images.githubusercontent.com/7365912/265043418-23291941-864d-495a-8ba8-d02e05756396.gif" alt="input-frames"></td> | |
| </tr> | |
| <tr><td colspan="2" align=center><b>AnimateDiff model: SG161222/Realistic_Vision_V5.1_noVAE</b></td></tr> | |
| <tr> | |
| <td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/baf301e2-d03c-4129-bd84-203a1de2b2be" alt="gif-1"></td> | |
| <td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/9f923475-ecaf-452b-92c8-4e42171182d8" alt="gif-2"></td> | |
| </tr> | |
| <tr><td colspan="2" align=center><b>AnimateDiff model: CardosAnime</b></td></tr> | |
| <tr> | |
| <td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/b2c41028-38a0-45d6-86ed-fec7446b87f7" alt="gif-1"></td> | |
| <td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/eb7d2952-72e4-44fa-b664-077c79b4fc70" alt="gif-2"></td> | |
| </tr> | |
| </table> | |
| You can also use multiple controlnets at once! | |
| ```python | |
| import torch | |
| from diffusers import AutoencoderKL, ControlNetModel, MotionAdapter | |
| from diffusers.pipelines import DiffusionPipeline | |
| from diffusers.schedulers import DPMSolverMultistepScheduler | |
| from PIL import Image | |
| motion_id = "guoyww/animatediff-motion-adapter-v1-5-2" | |
| adapter = MotionAdapter.from_pretrained(motion_id) | |
| controlnet1 = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_openpose", torch_dtype=torch.float16) | |
| controlnet2 = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16) | |
| vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16) | |
| model_id = "SG161222/Realistic_Vision_V5.1_noVAE" | |
| pipe = DiffusionPipeline.from_pretrained( | |
| model_id, | |
| motion_adapter=adapter, | |
| controlnet=[controlnet1, controlnet2], | |
| vae=vae, | |
| custom_pipeline="pipeline_animatediff_controlnet", | |
| ).to(device="cuda", dtype=torch.float16) | |
| pipe.scheduler = DPMSolverMultistepScheduler.from_pretrained( | |
| model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", steps_offset=1, beta_schedule="linear", | |
| ) | |
| pipe.enable_vae_slicing() | |
| def load_video(file_path: str): | |
| images = [] | |
| if file_path.startswith(('http://', 'https://')): | |
| # If the file_path is a URL | |
| response = requests.get(file_path) | |
| response.raise_for_status() | |
| content = BytesIO(response.content) | |
| vid = imageio.get_reader(content) | |
| else: | |
| # Assuming it's a local file path | |
| vid = imageio.get_reader(file_path) | |
| for frame in vid: | |
| pil_image = Image.fromarray(frame) | |
| images.append(pil_image) | |
| return images | |
| video = load_video("dance.gif") | |
| # You need to install it using `pip install controlnet_aux` | |
| from controlnet_aux.processor import Processor | |
| p1 = Processor("openpose_full") | |
| cn1 = [p1(frame) for frame in video] | |
| p2 = Processor("canny") | |
| cn2 = [p2(frame) for frame in video] | |
| prompt = "astronaut in space, dancing" | |
| negative_prompt = "bad quality, worst quality, jpeg artifacts, ugly" | |
| result = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| width=512, | |
| height=768, | |
| conditioning_frames=[cn1, cn2], | |
| num_inference_steps=20, | |
| ) | |
| from diffusers.utils import export_to_gif | |
| export_to_gif(result.frames[0], "result.gif") | |
| ``` | |
| ### DemoFusion | |
| This pipeline is the official implementation of [DemoFusion: Democratising High-Resolution Image Generation With No $$$](https://arxiv.org/abs/2311.16973). | |
| The original repo can be found at [repo](https://github.com/PRIS-CV/DemoFusion). | |
| - `view_batch_size` (`int`, defaults to 16): | |
| The batch size for multiple denoising paths. Typically, a larger batch size can result in higher efficiency but comes with increased GPU memory requirements. | |
| - `stride` (`int`, defaults to 64): | |
| The stride of moving local patches. A smaller stride is better for alleviating seam issues, but it also introduces additional computational overhead and inference time. | |
| - `cosine_scale_1` (`float`, defaults to 3): | |
| Control the strength of skip-residual. For specific impacts, please refer to Appendix C in the DemoFusion paper. | |
| - `cosine_scale_2` (`float`, defaults to 1): | |
| Control the strength of dilated sampling. For specific impacts, please refer to Appendix C in the DemoFusion paper. | |
| - `cosine_scale_3` (`float`, defaults to 1): | |
| Control the strength of the Gaussian filter. For specific impacts, please refer to Appendix C in the DemoFusion paper. | |
| - `sigma` (`float`, defaults to 1): | |
| The standard value of the Gaussian filter. Larger sigma promotes the global guidance of dilated sampling, but has the potential of over-smoothing. | |
| - `multi_decoder` (`bool`, defaults to True): | |
| Determine whether to use a tiled decoder. Generally, when the resolution exceeds 3072x3072, a tiled decoder becomes necessary. | |
| - `show_image` (`bool`, defaults to False): | |
| Determine whether to show intermediate results during generation. | |
| ```py | |
| from diffusers import DiffusionPipeline | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "stabilityai/stable-diffusion-xl-base-1.0", | |
| custom_pipeline="pipeline_demofusion_sdxl", | |
| custom_revision="main", | |
| torch_dtype=torch.float16, | |
| ) | |
| pipe = pipe.to("cuda") | |
| prompt = "Envision a portrait of an elderly woman, her face a canvas of time, framed by a headscarf with muted tones of rust and cream. Her eyes, blue like faded denim. Her attire, simple yet dignified." | |
| negative_prompt = "blurry, ugly, duplicate, poorly drawn, deformed, mosaic" | |
| images = pipe( | |
| prompt, | |
| negative_prompt=negative_prompt, | |
| height=3072, | |
| width=3072, | |
| view_batch_size=16, | |
| stride=64, | |
| num_inference_steps=50, | |
| guidance_scale=7.5, | |
| cosine_scale_1=3, | |
| cosine_scale_2=1, | |
| cosine_scale_3=1, | |
| sigma=0.8, | |
| multi_decoder=True, | |
| show_image=True | |
| ) | |
| ``` | |
| You can display and save the generated images as: | |
| ```py | |
| def image_grid(imgs, save_path=None): | |
| w = 0 | |
| for i, img in enumerate(imgs): | |
| h_, w_ = imgs[i].size | |
| w += w_ | |
| h = h_ | |
| grid = Image.new('RGB', size=(w, h)) | |
| grid_w, grid_h = grid.size | |
| w = 0 | |
| for i, img in enumerate(imgs): | |
| h_, w_ = imgs[i].size | |
| grid.paste(img, box=(w, h - h_)) | |
| if save_path != None: | |
| img.save(save_path + "/img_{}.jpg".format((i + 1) * 1024)) | |
| w += w_ | |
| return grid | |
| image_grid(images, save_path="./outputs/") | |
| ``` | |
|  | |
| ### SDE Drag pipeline | |
| This pipeline provides drag-and-drop image editing using stochastic differential equations. It enables image editing by inputting prompt, image, mask_image, source_points, and target_points. | |
|  | |
| See [paper](https://arxiv.org/abs/2311.01410), [paper page](https://ml-gsai.github.io/SDE-Drag-demo/), [original repo](https://github.com/ML-GSAI/SDE-Drag) for more infomation. | |
| ```py | |
| import PIL | |
| import torch | |
| from diffusers import DDIMScheduler, DiffusionPipeline | |
| # Load the pipeline | |
| model_path = "runwayml/stable-diffusion-v1-5" | |
| scheduler = DDIMScheduler.from_pretrained(model_path, subfolder="scheduler") | |
| pipe = DiffusionPipeline.from_pretrained(model_path, scheduler=scheduler, custom_pipeline="sde_drag") | |
| pipe.to('cuda') | |
| # To save GPU memory, torch.float16 can be used, but it may compromise image quality. | |
| # If not training LoRA, please avoid using torch.float16 | |
| # pipe.to(torch.float16) | |
| # Provide prompt, image, mask image, and the starting and target points for drag editing. | |
| prompt = "prompt of the image" | |
| image = PIL.Image.open('/path/to/image') | |
| mask_image = PIL.Image.open('/path/to/mask_image') | |
| source_points = [[123, 456]] | |
| target_points = [[234, 567]] | |
| # train_lora is optional, and in most cases, using train_lora can better preserve consistency with the original image. | |
| pipe.train_lora(prompt, image) | |
| output = pipe(prompt, image, mask_image, source_points, target_points) | |
| output_image = PIL.Image.fromarray(output) | |
| output_image.save("./output.png") | |
| ``` | |
| ### Instaflow Pipeline | |
| InstaFlow is an ultra-fast, one-step image generator that achieves image quality close to Stable Diffusion, significantly reducing the demand of computational resources. This efficiency is made possible through a recent [Rectified Flow](https://github.com/gnobitab/RectifiedFlow) technique, which trains probability flows with straight trajectories, hence inherently requiring only a single step for fast inference. | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained("XCLIU/instaflow_0_9B_from_sd_1_5", torch_dtype=torch.float16, custom_pipeline="instaflow_one_step") | |
| pipe.to("cuda") ### if GPU is not available, comment this line | |
| prompt = "A hyper-realistic photo of a cute cat." | |
| images = pipe(prompt=prompt, | |
| num_inference_steps=1, | |
| guidance_scale=0.0).images | |
| images[0].save("./image.png") | |
| ``` | |
|  | |
| You can also combine it with LORA out of the box, like <https://huggingface.co/artificialguybr/logo-redmond-1-5v-logo-lora-for-liberteredmond-sd-1-5>, to unlock cool use cases in single step! | |
| ```python | |
| from diffusers import DiffusionPipeline | |
| import torch | |
| pipe = DiffusionPipeline.from_pretrained("XCLIU/instaflow_0_9B_from_sd_1_5", torch_dtype=torch.float16, custom_pipeline="instaflow_one_step") | |
| pipe.to("cuda") ### if GPU is not available, comment this line | |
| pipe.load_lora_weights("artificialguybr/logo-redmond-1-5v-logo-lora-for-liberteredmond-sd-1-5") | |
| prompt = "logo, A logo for a fitness app, dynamic running figure, energetic colors (red, orange) ),LogoRedAF ," | |
| images = pipe(prompt=prompt, | |
| num_inference_steps=1, | |
| guidance_scale=0.0).images | |
| images[0].save("./image.png") | |
| ``` | |
|  | |
| ### Null-Text Inversion pipeline | |
| This pipeline provides null-text inversion for editing real images. It enables null-text optimization, and DDIM reconstruction via w, w/o null-text optimization. No prompt-to-prompt code is implemented as there is a Prompt2PromptPipeline. | |
| - Reference paper | |
| ```@article{hertz2022prompt, | |
| title={Prompt-to-prompt image editing with cross attention control}, | |
| author={Hertz, Amir and Mokady, Ron and Tenenbaum, Jay and Aberman, Kfir and Pritch, Yael and Cohen-Or, Daniel}, | |
| booktitle={arXiv preprint arXiv:2208.01626}, | |
| year={2022} | |
| ```} | |
| ```py | |
| from diffusers.schedulers import DDIMScheduler | |
| from examples.community.pipeline_null_text_inversion import NullTextPipeline | |
| import torch | |
| # Load the pipeline | |
| device = "cuda" | |
| # Provide invert_prompt and the image for null-text optimization. | |
| invert_prompt = "A lying cat" | |
| input_image = "siamese.jpg" | |
| steps = 50 | |
| # Provide prompt used for generation. Same if reconstruction | |
| prompt = "A lying cat" | |
| # or different if editing. | |
| prompt = "A lying dog" | |
| #Float32 is essential to a well optimization | |
| model_path = "runwayml/stable-diffusion-v1-5" | |
| scheduler = DDIMScheduler(num_train_timesteps=1000, beta_start=0.00085, beta_end=0.0120, beta_schedule="scaled_linear") | |
| pipeline = NullTextPipeline.from_pretrained(model_path, scheduler = scheduler, torch_dtype=torch.float32).to(device) | |
| #Saves the inverted_latent to save time | |
| inverted_latent, uncond = pipeline.invert(input_image, invert_prompt, num_inner_steps=10, early_stop_epsilon= 1e-5, num_inference_steps = steps) | |
| pipeline(prompt, uncond, inverted_latent, guidance_scale=7.5, num_inference_steps=steps).images[0].save(input_image+".output.jpg") | |
| ``` | |
| ### Rerender A Video | |
| This is the Diffusers implementation of zero-shot video-to-video translation pipeline [Rerender A Video](https://github.com/williamyang1991/Rerender_A_Video) (without Ebsynth postprocessing). To run the code, please install gmflow. Then modify the path in `gmflow_dir`. After that, you can run the pipeline with: | |
| ```py | |
| import sys | |
| gmflow_dir = "/path/to/gmflow" | |
| sys.path.insert(0, gmflow_dir) | |
| from diffusers import ControlNetModel, AutoencoderKL, DDIMScheduler | |
| from diffusers.utils import export_to_video | |
| import numpy as np | |
| import torch | |
| import cv2 | |
| from PIL import Image | |
| def video_to_frame(video_path: str, interval: int): | |
| vidcap = cv2.VideoCapture(video_path) | |
| success = True | |
| count = 0 | |
| res = [] | |
| while success: | |
| count += 1 | |
| success, image = vidcap.read() | |
| if count % interval != 1: | |
| continue | |
| if image is not None: | |
| image = cv2.cvtColor(image, cv2.COLOR_BGR2RGB) | |
| res.append(image) | |
| vidcap.release() | |
| return res | |
| input_video_path = 'path/to/video' | |
| input_interval = 10 | |
| frames = video_to_frame( | |
| input_video_path, input_interval) | |
| control_frames = [] | |
| # get canny image | |
| for frame in frames: | |
| np_image = cv2.Canny(frame, 50, 100) | |
| np_image = np_image[:, :, None] | |
| np_image = np.concatenate([np_image, np_image, np_image], axis=2) | |
| canny_image = Image.fromarray(np_image) | |
| control_frames.append(canny_image) | |
| # You can use any ControlNet here | |
| controlnet = ControlNetModel.from_pretrained( | |
| "lllyasviel/sd-controlnet-canny").to('cuda') | |
| # You can use any fintuned SD here | |
| pipe = DiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", controlnet=controlnet, custom_pipeline='rerender_a_video').to('cuda') | |
| # Optional: you can download vae-ft-mse-840000-ema-pruned.ckpt to enhance the results | |
| # pipe.vae = AutoencoderKL.from_single_file( | |
| # "path/to/vae-ft-mse-840000-ema-pruned.ckpt").to('cuda') | |
| pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config) | |
| generator = torch.manual_seed(0) | |
| frames = [Image.fromarray(frame) for frame in frames] | |
| output_frames = pipe( | |
| "a beautiful woman in CG style, best quality, extremely detailed", | |
| frames, | |
| control_frames, | |
| num_inference_steps=20, | |
| strength=0.75, | |
| controlnet_conditioning_scale=0.7, | |
| generator=generator, | |
| warp_start=0.0, | |
| warp_end=0.1, | |
| mask_start=0.5, | |
| mask_end=0.8, | |
| mask_strength=0.5, | |
| negative_prompt='longbody, lowres, bad anatomy, bad hands, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality' | |
| ).frames[0] | |
| export_to_video( | |
| output_frames, "/path/to/video.mp4", 5) | |
| ``` | |
| ### StyleAligned Pipeline | |
| This pipeline is the implementation of [Style Aligned Image Generation via Shared Attention](https://arxiv.org/abs/2312.02133). You can find more results [here](https://github.com/huggingface/diffusers/pull/6489#issuecomment-1881209354). | |
| > Large-scale Text-to-Image (T2I) models have rapidly gained prominence across creative fields, generating visually compelling outputs from textual prompts. However, controlling these models to ensure consistent style remains challenging, with existing methods necessitating fine-tuning and manual intervention to disentangle content and style. In this paper, we introduce StyleAligned, a novel technique designed to establish style alignment among a series of generated images. By employing minimal `attention sharing' during the diffusion process, our method maintains style consistency across images within T2I models. This approach allows for the creation of style-consistent images using a reference style through a straightforward inversion operation. Our method's evaluation across diverse styles and text prompts demonstrates high-quality synthesis and fidelity, underscoring its efficacy in achieving consistent style across various inputs. | |
| ```python | |
| from typing import List | |
| import torch | |
| from diffusers.pipelines.pipeline_utils import DiffusionPipeline | |
| from PIL import Image | |
| model_id = "a-r-r-o-w/dreamshaper-xl-turbo" | |
| pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16, variant="fp16", custom_pipeline="pipeline_sdxl_style_aligned") | |
| pipe = pipe.to("cuda") | |
| # Enable memory saving techniques | |
| pipe.enable_vae_slicing() | |
| pipe.enable_vae_tiling() | |
| prompt = [ | |
| "a toy train. macro photo. 3d game asset", | |
| "a toy airplane. macro photo. 3d game asset", | |
| "a toy bicycle. macro photo. 3d game asset", | |
| "a toy car. macro photo. 3d game asset", | |
| ] | |
| negative_prompt = "low quality, worst quality, " | |
| # Enable StyleAligned | |
| pipe.enable_style_aligned( | |
| share_group_norm=False, | |
| share_layer_norm=False, | |
| share_attention=True, | |
| adain_queries=True, | |
| adain_keys=True, | |
| adain_values=False, | |
| full_attention_share=False, | |
| shared_score_scale=1.0, | |
| shared_score_shift=0.0, | |
| only_self_level=0.0, | |
| ) | |
| # Run inference | |
| images = pipe( | |
| prompt=prompt, | |
| negative_prompt=negative_prompt, | |
| guidance_scale=2, | |
| height=1024, | |
| width=1024, | |
| num_inference_steps=10, | |
| generator=torch.Generator().manual_seed(42), | |
| ).images | |
| # Disable StyleAligned if you do not wish to use it anymore | |
| pipe.disable_style_aligned() | |
| ``` | |
| ### AnimateDiff Image-To-Video Pipeline | |
| This pipeline adds experimental support for the image-to-video task using AnimateDiff. Refer to [this](https://github.com/huggingface/diffusers/pull/6328) PR for more examples and results. | |
| This pipeline relies on a "hack" discovered by the community that allows the generation of videos given an input image with AnimateDiff. It works by creating a copy of the image `num_frames` times and progressively adding more noise to the image based on the strength and latent interpolation method. | |
| ```py | |
| import torch | |
| from diffusers import MotionAdapter, DiffusionPipeline, DDIMScheduler | |
| from diffusers.utils import export_to_gif, load_image | |
| model_id = "SG161222/Realistic_Vision_V5.1_noVAE" | |
| adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2") | |
| pipe = DiffusionPipeline.from_pretrained(model_id, motion_adapter=adapter, custom_pipeline="pipeline_animatediff_img2video").to("cuda") | |
| pipe.scheduler = DDIMScheduler.from_pretrained(model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", beta_schedule="linear", steps_offset=1) | |
| image = load_image("snail.png") | |
| output = pipe( | |
| image=image, | |
| prompt="A snail moving on the ground", | |
| strength=0.8, | |
| latent_interpolation_method="slerp", # can be lerp, slerp, or your own callback | |
| ) | |
| frames = output.frames[0] | |
| export_to_gif(frames, "animation.gif") | |
| ``` | |
| ### IP Adapter Face ID | |
| IP Adapter FaceID is an experimental IP Adapter model that uses image embeddings generated by `insightface`, so no image encoder needs to be loaded. | |
| You need to install `insightface` and all its requirements to use this model. | |
| You must pass the image embedding tensor as `image_embeds` to the StableDiffusionPipeline instead of `ip_adapter_image`. | |
| You can find more results [here](https://github.com/huggingface/diffusers/pull/6276). | |
| ```py | |
| import diffusers | |
| import torch | |
| from diffusers.utils import load_image | |
| import cv2 | |
| import numpy as np | |
| from diffusers import DiffusionPipeline, AutoencoderKL, DDIMScheduler | |
| from insightface.app import FaceAnalysis | |
| noise_scheduler = DDIMScheduler( | |
| num_train_timesteps=1000, | |
| beta_start=0.00085, | |
| beta_end=0.012, | |
| beta_schedule="scaled_linear", | |
| clip_sample=False, | |
| set_alpha_to_one=False, | |
| steps_offset=1, | |
| ) | |
| vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse").to(dtype=torch.float16) | |
| pipeline = DiffusionPipeline.from_pretrained( | |
| "SG161222/Realistic_Vision_V4.0_noVAE", | |
| torch_dtype=torch.float16, | |
| scheduler=noise_scheduler, | |
| vae=vae, | |
| custom_pipeline="ip_adapter_face_id" | |
| ) | |
| pipeline.load_ip_adapter_face_id("h94/IP-Adapter-FaceID", "ip-adapter-faceid_sd15.bin") | |
| pipeline.to("cuda") | |
| generator = torch.Generator(device="cpu").manual_seed(42) | |
| num_images=2 | |
| image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png") | |
| app = FaceAnalysis(name="buffalo_l", providers=['CUDAExecutionProvider', 'CPUExecutionProvider']) | |
| app.prepare(ctx_id=0, det_size=(640, 640)) | |
| image = cv2.cvtColor(np.asarray(image), cv2.COLOR_BGR2RGB) | |
| faces = app.get(image) | |
| image = torch.from_numpy(faces[0].normed_embedding).unsqueeze(0) | |
| images = pipeline( | |
| prompt="A photo of a girl wearing a black dress, holding red roses in hand, upper body, behind is the Eiffel Tower", | |
| image_embeds=image, | |
| negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality", | |
| num_inference_steps=20, num_images_per_prompt=num_images, width=512, height=704, | |
| generator=generator | |
| ).images | |
| for i in range(num_images): | |
| images[i].save(f"c{i}.png") | |
| ``` | |
| ### InstantID Pipeline | |
| InstantID is a new state-of-the-art tuning-free method to achieve ID-Preserving generation with only single image, supporting various downstream tasks. For any usgae question, please refer to the [official implementation](https://github.com/InstantID/InstantID). | |
| ```py | |
| # !pip install opencv-python transformers accelerate insightface | |
| import diffusers | |
| from diffusers.utils import load_image | |
| from diffusers.models import ControlNetModel | |
| import cv2 | |
| import torch | |
| import numpy as np | |
| from PIL import Image | |
| from insightface.app import FaceAnalysis | |
| from pipeline_stable_diffusion_xl_instantid import StableDiffusionXLInstantIDPipeline, draw_kps | |
| # prepare 'antelopev2' under ./models | |
| # https://github.com/deepinsight/insightface/issues/1896#issuecomment-1023867304 | |
| app = FaceAnalysis(name='antelopev2', root='./', providers=['CUDAExecutionProvider', 'CPUExecutionProvider']) | |
| app.prepare(ctx_id=0, det_size=(640, 640)) | |
| # prepare models under ./checkpoints | |
| # https://huggingface.co/InstantX/InstantID | |
| from huggingface_hub import hf_hub_download | |
| hf_hub_download(repo_id="InstantX/InstantID", filename="ControlNetModel/config.json", local_dir="./checkpoints") | |
| hf_hub_download(repo_id="InstantX/InstantID", filename="ControlNetModel/diffusion_pytorch_model.safetensors", local_dir="./checkpoints") | |
| hf_hub_download(repo_id="InstantX/InstantID", filename="ip-adapter.bin", local_dir="./checkpoints") | |
| face_adapter = f'./checkpoints/ip-adapter.bin' | |
| controlnet_path = f'./checkpoints/ControlNetModel' | |
| # load IdentityNet | |
| controlnet = ControlNetModel.from_pretrained(controlnet_path, torch_dtype=torch.float16) | |
| base_model = 'wangqixun/YamerMIX_v8' | |
| pipe = StableDiffusionXLInstantIDPipeline.from_pretrained( | |
| base_model, | |
| controlnet=controlnet, | |
| torch_dtype=torch.float16 | |
| ) | |
| pipe.cuda() | |
| # load adapter | |
| pipe.load_ip_adapter_instantid(face_adapter) | |
| # load an image | |
| face_image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png") | |
| # prepare face emb | |
| face_info = app.get(cv2.cvtColor(np.array(face_image), cv2.COLOR_RGB2BGR)) | |
| face_info = sorted(face_info, key=lambda x:(x['bbox'][2]-x['bbox'][0])*x['bbox'][3]-x['bbox'][1])[-1] # only use the maximum face | |
| face_emb = face_info['embedding'] | |
| face_kps = draw_kps(face_image, face_info['kps']) | |
| # prompt | |
| prompt = "film noir style, ink sketch|vector, male man, highly detailed, sharp focus, ultra sharpness, monochrome, high contrast, dramatic shadows, 1940s style, mysterious, cinematic" | |
| negative_prompt = "ugly, deformed, noisy, blurry, low contrast, realism, photorealistic, vibrant, colorful" | |
| # generate image | |
| pipe.set_ip_adapter_scale(0.8) | |
| image = pipe( | |
| prompt, | |
| image_embeds=face_emb, | |
| image=face_kps, | |
| controlnet_conditioning_scale=0.8, | |
| ).images[0] | |
| ``` | |
| ### UFOGen Scheduler | |
| [UFOGen](https://arxiv.org/abs/2311.09257) is a generative model designed for fast one-step text-to-image generation, trained via adversarial training starting from an initial pretrained diffusion model such as Stable Diffusion. `scheduling_ufogen.py` implements a onestep and multistep sampling algorithm for UFOGen models compatible with pipelines like `StableDiffusionPipeline`. A usage example is as follows: | |
| ```py | |
| import torch | |
| from diffusers import StableDiffusionPipeline | |
| from scheduling_ufogen import UFOGenScheduler | |
| # NOTE: currently, I am not aware of any publicly available UFOGen model checkpoints trained from SD v1.5. | |
| ufogen_model_id_or_path = "/path/to/ufogen/model" | |
| pipe = StableDiffusionPipeline( | |
| ufogen_model_id_or_path, | |
| torch_dtype=torch.float16, | |
| ) | |
| # You can initialize a UFOGenScheduler as follows: | |
| pipe.scheduler = UFOGenScheduler.from_config(pipe.scheduler.config) | |
| prompt = "Three cats having dinner at a table at new years eve, cinematic shot, 8k." | |
| # Onestep sampling | |
| onestep_image = pipe(prompt, num_inference_steps=1).images[0] | |
| # Multistep sampling | |
| multistep_image = pipe(prompt, num_inference_steps=4).images[0] | |
| ``` | |
| ### FRESCO | |
| This is the Diffusers implementation of zero-shot video-to-video translation pipeline [FRESCO](https://github.com/williamyang1991/FRESCO) (without Ebsynth postprocessing and background smooth). To run the code, please install gmflow. Then modify the path in `gmflow_dir`. After that, you can run the pipeline with: | |
| ```py | |
| from PIL import Image | |
| import cv2 | |
| import torch | |
| import numpy as np | |
| from diffusers import ControlNetModel,DDIMScheduler, DiffusionPipeline | |
| import sys | |
| gmflow_dir = "/path/to/gmflow" | |
| sys.path.insert(0, gmflow_dir) | |
| def video_to_frame(video_path: str, interval: int): | |
| vidcap = cv2.VideoCapture(video_path) | |
| success = True | |
| count = 0 | |
| res = [] | |
| while success: | |
| count += 1 | |
| success, image = vidcap.read() | |
| if count % interval != 1: | |
| continue | |
| if image is not None: | |
| image = cv2.cvtColor(image, cv2.COLOR_BGR2RGB) | |
| res.append(image) | |
| if len(res) >= 8: | |
| break | |
| vidcap.release() | |
| return res | |
| input_video_path = 'https://github.com/williamyang1991/FRESCO/raw/main/data/car-turn.mp4' | |
| output_video_path = 'car.gif' | |
| # You can use any fintuned SD here | |
| model_path = 'SG161222/Realistic_Vision_V2.0' | |
| prompt = 'a red car turns in the winter' | |
| a_prompt = ', RAW photo, subject, (high detailed skin:1.2), 8k uhd, dslr, soft lighting, high quality, film grain, Fujifilm XT3, ' | |
| n_prompt = '(deformed iris, deformed pupils, semi-realistic, cgi, 3d, render, sketch, cartoon, drawing, anime, mutated hands and fingers:1.4), (deformed, distorted, disfigured:1.3), poorly drawn, bad anatomy, wrong anatomy, extra limb, missing limb, floating limbs, disconnected limbs, mutation, mutated, ugly, disgusting, amputation' | |
| input_interval = 5 | |
| frames = video_to_frame( | |
| input_video_path, input_interval) | |
| control_frames = [] | |
| # get canny image | |
| for frame in frames: | |
| image = cv2.Canny(frame, 50, 100) | |
| np_image = np.array(image) | |
| np_image = np_image[:, :, None] | |
| np_image = np.concatenate([np_image, np_image, np_image], axis=2) | |
| canny_image = Image.fromarray(np_image) | |
| control_frames.append(canny_image) | |
| # You can use any ControlNet here | |
| controlnet = ControlNetModel.from_pretrained( | |
| "lllyasviel/sd-controlnet-canny").to('cuda') | |
| pipe = DiffusionPipeline.from_pretrained( | |
| model_path, controlnet=controlnet, custom_pipeline='fresco_v2v').to('cuda') | |
| pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config) | |
| generator = torch.manual_seed(0) | |
| frames = [Image.fromarray(frame) for frame in frames] | |
| output_frames = pipe( | |
| prompt + a_prompt, | |
| frames, | |
| control_frames, | |
| num_inference_steps=20, | |
| strength=0.75, | |
| controlnet_conditioning_scale=0.7, | |
| generator=generator, | |
| negative_prompt=n_prompt | |
| ).images | |
| output_frames[0].save(output_video_path, save_all=True, | |
| append_images=output_frames[1:], duration=100, loop=0) | |
| ``` | |
| # Perturbed-Attention Guidance | |
| [Project](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) / [arXiv](https://arxiv.org/abs/2403.17377) / [GitHub](https://github.com/KU-CVLAB/Perturbed-Attention-Guidance) | |
| This implementation is based on [Diffusers](https://huggingface.co/docs/diffusers/index). StableDiffusionPAGPipeline is a modification of StableDiffusionPipeline to support Perturbed-Attention Guidance (PAG). | |
| ## Example Usage | |
| ```py | |
| import os | |
| import torch | |
| from accelerate.utils import set_seed | |
| from diffusers import StableDiffusionPipeline | |
| from diffusers.utils import load_image, make_image_grid | |
| from diffusers.utils.torch_utils import randn_tensor | |
| pipe = StableDiffusionPipeline.from_pretrained( | |
| "runwayml/stable-diffusion-v1-5", | |
| custom_pipeline="hyoungwoncho/sd_perturbed_attention_guidance", | |
| torch_dtype=torch.float16 | |
| ) | |
| device="cuda" | |
| pipe = pipe.to(device) | |
| pag_scale = 5.0 | |
| pag_applied_layers_index = ['m0'] | |
| batch_size = 4 | |
| seed=10 | |
| base_dir = "./results/" | |
| grid_dir = base_dir + "/pag" + str(pag_scale) + "/" | |
| if not os.path.exists(grid_dir): | |
| os.makedirs(grid_dir) | |
| set_seed(seed) | |
| latent_input = randn_tensor(shape=(batch_size,4,64,64),generator=None, device=device, dtype=torch.float16) | |
| output_baseline = pipe( | |
| "", | |
| width=512, | |
| height=512, | |
| num_inference_steps=50, | |
| guidance_scale=0.0, | |
| pag_scale=0.0, | |
| pag_applied_layers_index=pag_applied_layers_index, | |
| num_images_per_prompt=batch_size, | |
| latents=latent_input | |
| ).images | |
| output_pag = pipe( | |
| "", | |
| width=512, | |
| height=512, | |
| num_inference_steps=50, | |
| guidance_scale=0.0, | |
| pag_scale=5.0, | |
| pag_applied_layers_index=pag_applied_layers_index, | |
| num_images_per_prompt=batch_size, | |
| latents=latent_input | |
| ).images | |
| grid_image = make_image_grid(output_baseline + output_pag, rows=2, cols=batch_size) | |
| grid_image.save(grid_dir + "sample.png") | |
| ``` | |
| ## PAG Parameters | |
| pag_scale : gudiance scale of PAG (ex: 5.0) | |
| pag_applied_layers_index : index of the layer to apply perturbation (ex: ['m0']) | |