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--- |
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license: openrail++ |
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tags: |
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- stable-diffusion |
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- text-to-image |
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--- |
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# SD-XL 1.0-refiner Model Card |
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![row01](01.png) |
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## Model |
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![pipeline](pipeline.png) |
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[SDXL](https://arxiv.org/abs/2307.01952) consists of a mixture-of-experts pipeline for latent diffusion: |
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In a first step, the base model (available here: https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) is used to generate (noisy) latents, |
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which are then further processed with a refinement model specialized for the final denoising steps. |
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Note that the base model can be used as a standalone module. |
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Alternatively, we can use a two-stage pipeline as follows: |
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First, the base model is used to generate latents of the desired output size. |
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In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") |
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to the latents generated in the first step, using the same prompt. This technique is slightly slower than the first one, as it requires more function evaluations. |
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Source code is available at https://github.com/Stability-AI/generative-models . |
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### Model Description |
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- **Developed by:** Stability AI |
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- **Model type:** Diffusion-based text-to-image generative model |
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- **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/blob/main/LICENSE.md) |
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- **Model Description:** This is a model that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses two fixed, pretrained text encoders ([OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip) and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main)). |
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- **Resources for more information:** Check out our [GitHub Repository](https://github.com/Stability-AI/generative-models) and the [SDXL report on arXiv](https://arxiv.org/abs/2307.01952). |
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### Model Sources |
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For research purposes, we recommned our `generative-models` Github repository (https://github.com/Stability-AI/generative-models), which implements the most popoular diffusion frameworks (both training and inference) and for which new functionalities like distillation will be added over time. |
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[Clipdrop](https://clipdrop.co/stable-diffusion) provides free SDXL inference. |
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- **Repository:** https://github.com/Stability-AI/generative-models |
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- **Demo:** https://clipdrop.co/stable-diffusion |
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## Evaluation |
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![comparison](comparison.png) |
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The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0.9 and Stable Diffusion 1.5 and 2.1. |
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The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. |
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### 🧨 Diffusers |
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Make sure to upgrade diffusers to >= 0.18.0: |
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``` |
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pip install diffusers --upgrade |
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``` |
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In addition make sure to install `transformers`, `safetensors`, `accelerate` as well as the invisible watermark: |
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``` |
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pip install invisible_watermark transformers accelerate safetensors |
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``` |
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You can use the model then as follows |
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```py |
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from diffusers import DiffusionPipeline |
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import torch |
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pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, use_safetensors=True, variant="fp16") |
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pipe.to("cuda") |
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# if using torch < 2.0 |
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# pipe.enable_xformers_memory_efficient_attention() |
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prompt = "An astronaut riding a green horse" |
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images = pipe(prompt=prompt).images[0] |
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``` |
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When using `torch >= 2.0`, you can improve the inference speed by 20-30% with torch.compile. Simple wrap the unet with torch compile before running the pipeline: |
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```py |
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pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True) |
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``` |
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If you are limited by GPU VRAM, you can enable *cpu offloading* by calling `pipe.enable_model_cpu_offload` |
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instead of `.to("cuda")`: |
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```diff |
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- pipe.to("cuda") |
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+ pipe.enable_model_cpu_offload() |
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``` |
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## Uses |
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### Direct Use |
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The model is intended for research purposes only. Possible research areas and tasks include |
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- Generation of artworks and use in design and other artistic processes. |
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- Applications in educational or creative tools. |
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- Research on generative models. |
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- Safe deployment of models which have the potential to generate harmful content. |
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- Probing and understanding the limitations and biases of generative models. |
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Excluded uses are described below. |
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### Out-of-Scope Use |
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The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model. |
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## Limitations and Bias |
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### Limitations |
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- The model does not achieve perfect photorealism |
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- The model cannot render legible text |
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- The model struggles with more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere” |
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- Faces and people in general may not be generated properly. |
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- The autoencoding part of the model is lossy. |
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### Bias |
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While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases. |