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In a vast disc of gas and dust particles circling a young star, scientists have found evidence of a hypothesized but never-seen dust trap that may solve the mystery of how planets form. We know planets that orbit stars are abundant throughout our galaxy, and likely throughout the universe as well, but until recently, scientists weren't exactly sure how those planets came to be. The working theory is that they grew over time as tiny bits of dust collided and stuck together -- eventually forming comets, rocky planets and the cores of gaseous planets over millions of years. But there is a problem with that theory: Once these tiny bits of dust grow to the size of pebbles or boulders, they are likely to either smash into one another and break apart, or spiral toward their central star where their growth is inhibited. Theoretical astronomers had hypothesized that the flat discs of dust and gas that often surround young stars might occasionally contain dust traps -- an area in the disc where the gas is more dense and can create a barrier that keeps more substantial bits of dust from falling toward the star. And then, quite by accident, a team of scientists found evidence of one in a disc around the massive young star Oph IRS 48, about 400 light-years from Earth. The team, led by Nienke van der Marel, a doctoral student at Leiden Observatory in the Netherlands, was using ALMA, an array of radio telescopes in Chile, to observe just the gas in the disc. But ALMA also gave them data about the dust in the ring, "for free," Van Der Marel said. When the researchers looked at the dust data, they were confused: They had expected to see dust particles distributed evenly throughout the disc, but the images they saw showed the dust clumped on about one-third of the disc in the shape of a cashew. "The first time we saw the image of the dust, we thought there must be something wrong with the data," Van Der Marel told the Los Angeles Times. "But we had a really high, clear signal so it was clear it wasn't a mistake. Then we started looking into possibilities that could explain the separation between the gas and the dust." It turns out that they had just made the first observation of a dust trap. This dust trap was caused by a large gaseous planet or perhaps a small star that is also circling the central star, Van Der Marel said. She and her colleagues had observed that there was a hole in the gas disc, likely filled by one of these two types of bodies (but they are not sure which one). The gas around this hole is more dense than the rest of the disc, and can keep the dust particles that get stuck behind it from falling toward the central star. Van Der Marel said this particular dust trap is not likely to create planets because of its location in the disc, but it could create comets as large as .6 miles across. Van Der Marel describes it as "comet factory." Next, Van Der Marel said she planned to use ALMA to look for dust traps that are closer to their stars, where planet formation is more likely. A study describing the dust trap was published in Science this week. See below for video of how scientists envision the dust trap works.

http://www.latimes.com/science/sciencenow/la-sci-sn-comet-factory-dust-trap-20130606-story.html

About 13,000 years ago, the Earth was plunged into what is called the 'Big Freeze' — or more formally known as the 'Younger Dryas stadial' — where the planet's climate cooled significantly, ushering in a new glacial period that lasted for about 1,300 years. It has been well-established that the Big Freeze was caused when a lake of melt-water sitting on the Laurentide Ice Sheet — a 2-3 km thick sheet of ice that covered most of the land-mass of what is now Canada and parts of the U.S. Midwest — broke through an ice-damn and rushed into the north Atlantic. This massive influx of frigid fresh water into the ocean disrupted the global circulation of heat and salt-content in the oceans, and quickly altered the Earth's climate. "This episode was the last time the Earth underwent a major cooling, so understanding exactly what caused it is very important for understanding how our modern-day climate might change in the future," says Alan Condron, a physical oceanographer with the University of Massachusetts Amherst's Climate System Research Center, according to Science Daily. There has been some debate over the years as to the path of this fresh water, though. The most commonly used hypothesis, first proposed by Wallace Broecker of Columbia University in 1989, was that the water flowed down the St. Lawrence River into the north Atlantic. Others suggested that it took a route down the Mackenzie River basin and into the Arctic Ocean. Now, Condron and research partner Peter Winsor, from the University of Alaska Fairbanks, have developed a new high-resolution ocean-ice circulation computer model that shows strong support for the latter idea. This computer model, the most powerful so far created, runs on a supercomputer at the National Energy Research Science Computing Center in Berkeley, California. "With this higher resolution modeling, our ability to capture narrow ocean currents dramatically improves our understanding of where the fresh water may be going." said Condron and Winsor. "The results we obtain are only possible by using a much higher computational power available with faster computers. Older models weren't powerful enough to model the different pathways because they contained too few data points to capture smaller-scale, faster-moving coastal currents." added Condron. Condron and Winsor's simulations showed that were the waters from the Laurentide ice sheet to have flowed down the St. Lawrence, they would have entered the Atlantic Ocean waters around 3000 kilometres too far south to have disrupted the ocean circulation enough to cause the Big Freeze. However, simulations showing the water flowing into the Arctic Ocean via the Mackenzie river basin showed that the currents in the Arctic Ocean would have transported the cold, fresh waters to exactly where they were needed to cause the event — the sub-polar Atlantic Ocean, off the coast of Greenland. "Dumping water in the Arctic is a very efficient way to … cool the Northern Hemisphere," says W. Richard Peltier, according to Science News. Peltier is a professor of physics at the University of Toronto, and director of UofT's Centre for Global Change Science. Although Broecker's hypothesis had wide support, there was a lack of physical evidence along the St. Lawrence River, however there is evidence in boulders and gravel along the Mackenzie River basin that supports the idea of a massive flood around the time of the Big Freeze. "This whole thing now hangs together beautifully," said Peltier. "Our results are particularly relevant for how we model the melting of the Greenland and Antarctic Ice sheets now and in the future," said Condron. "It is apparent from our results that climate scientists are artificially introducing fresh water into their models over large parts of the ocean that freshwater would never have reached. In addition, our work points to the Arctic as a primary trigger for climate change. This is especially relevant considering the rapid changes that have been occurring in this region in the last 10 years."

https://ca.news.yahoo.com/blogs/geekquinox/blame-canada-ancient-massive-1-300-big-freeze-191925701.html

Earth’s average temperature has remained more or less steady since 2001, despite rising levels of atmospheric carbon dioxide and other greenhouse gases—a trend that has perplexed most climate scientists. A new study suggests that the missing heat has been temporarily stirred into the relatively shallow waters in the western Pacific by stronger-than-normal trade winds. Over the past 20 years or so, trade winds near the equator—which generally blow from east to west—have driven warm waters of the Pacific ahead of them, causing larger-than-normal volumes of cool, deep waters to rise to the surface along the western coasts of Central America and South America. (Cooler-than-average surface waters are depicted in shades of blue, image from late July and early August 2007.) Climate simulations suggest that that upwelling has generally cooled Earth’s climate, stifling about 0.1°C to 0.2°C in warming that would have occurred by 2012 if winds hadn’t been inordinately strong, the researchers reported online yesterday in Nature Climate Change. Both real-world observations and the team’s simulations reveal that the abnormally strong winds—driven by natural variation in a long-term climate cycle called the Interdecadal Pacific Oscillation—have, for the time being, carried the “missing” heat to intermediate depths of the western Pacific Ocean. Eventually, possibly by the end of this decade, the inevitable slackening of the trade winds will bring the energy back to the ocean’s surface to be released to the atmosphere, fueling rapid warming, the scientists contend.

http://www.sciencemag.org/news/2014/02/scienceshot-pacific-ocean-keeping-earth-cool-now?mobile_switch=mobile

A classic rhyme, Simple Simon and the Pie-Man, introduces students to the concepts of consumer and producer. Students learn that consumers are the people who buy and use goods and services. Producers make the goods and provide the services. When producers are working, they often use goods and services provided by other producers. These goods and services are called resources. An interactive activity helps students distinguish between consumers and producers. In a second activity, the students match producers with the resources needed to provide goods and services. A suggested follow-up lesson is We are Consumers and Producers which examines how students and their families function as consumers and producers in their homes and communities. Students will be able to distinguish between people who produce goods and people who provide services to a community. This lesson will help students become good consumers and producers by taking turns buying and selling things in a classroom-created market. Students will establish prices for items and observe what happens during the sale of those items. The following lessons come from the Council for Economic Education's library of publications. Clicking the publication title or image will take you to the Council for Economic Education Store for more detailed information. This publication contains 16 stories that complement the K-2 Student Storybook. Specific to grades K-2 are a variety of activities, including making coins out of salt dough or cookie dough; a song that teaches students about opportunity cost and decisions; and a game in which students learn the importance of savings. 9 out of 18 lessons from this publication relate to this EconEdLink lesson. Designed primarily for elementary and middle school students, each of the 15 lessons in this guide introduces an economics concept through activities with modeling clay. 1 out of 17 lessons from this publication relate to this EconEdLink lesson. This interdisciplinary curriculum guide helps teachers introduce their students to economics using popular children's stories. 1 out of 13 lessons from this publication relate to this EconEdLink lesson.

http://www.econedlink.org/economic-standards/EconEdLink-related-publications.php?lid=457

A double standard is the application of different sets of principles for similar situations. A double standard may take the form of an instance in which certain concepts (often, for example, a word, phrase, social norm, or rule) are perceived as acceptable to be applied by one group of people, but are considered unacceptable—taboo—when applied by another group. The concept of a double standard has long been applied (as early as 1872) to the fact that different moral structures are often applied to men and women in society. An example being that a man going out to bars and picking up a different woman to have sex with every night for two weeks, will probably be considered "macho", a "stud", or a "ladies' man". All positive terms, but a woman who went home with 14 different men in a two-week period in the same way as the male example, typically would be labeled a "slut" or a "whore", both being pejorative. Conversely, if a man cries, he is commonly seen as "weak" or "pathetic", denoting a negative connotation; but if a woman cries, she is commonly seen as "innocent" and "sensitive", denoting a compassionate connotation. A double standard can therefore be described as a biased or morally unfair application of the principle that all are equal in their freedoms. Such double standards are seen as unjustified because they violate a basic maxim of modern legal jurisprudence: that all parties should stand equal before the law. Double standards also violate the principle of justice known as impartiality, which is based on the assumption that the same standards should be applied to all people, without regard to subjective bias or favoritism based on social class, rank, ethnicity, gender, religion, sexual orientation, age, or other distinctions. A double standard violates this principle by holding different people accountable according to different standards. History of the concept |This section is empty. You can help by adding to it. (December 2015)| Policy of double standards Policy of double standards is used to describe a situation when the assessment of the same phenomenon, process or event in the international relations, depends on character of the relations of the estimating parties with assessment objects. At identical intrinsic filling of action of one country get support and a justification, and other – is condemned and punished. The phrase became a classical example of policy of double standards: "One man's terrorist is another man's freedom fighter", entered into use by the British writer Gerald Seymour in his work "Harry's Game" in 1975. - "Double standard" Dictionary.com - "Unjust Judgments on Subjects of Morality". The Ecclesiastical Observer (London: Arthur Hall and Co.) XXV: 167–170. April 1, 1872. - Josephine E. Butler (Nov 27, 1886). "The Double Standard of Morality". Friends' Intelligencer and Journal (Philadelphia: Friends' Intelligencer Association). XLIII (48): 757–758. - Satish Chandra Pandey. International Terrorism and the Contemporary World. Sarup & Sons, 2006. С. 17. - Who said one man’s terrorist is another man’s revolutionary?

https://en.wikipedia.org/wiki/Double_standard

Where did plants come from? Plants' Adaptations for Life on Land The first photosynthetic organisms were bacteria that lived in the water. So, where did plants come from? Evidence shows that plants evolved from freshwater green algae, a protist (Figure below). The similarities between green algae and plants is one piece of evidence. They both have cellulose in their cell walls, and they share many of the same chemicals that give them color. So what separates green algae from green plants? The ancestor of plants is green algae. This picture shows a close up of algae on the beach. There are four main ways that plants adapted to life on land and, as a result, became different from algae: - In plants, the embryo develops inside of the female plant after fertilization. Algae do not keep the embryo inside of themselves but release it into water. This was the first feature to evolve that separated plants from green algae. This is also the only adaptation shared by all plants. - Over time, plants had to evolve from living in water to living on land. In early plants, a waxy layer called a cuticle evolved to help seal water in the plant and prevent water loss. However, the cuticle also prevents gases from entering and leaving the plant easily. Recall that the exchange of gasses—taking in carbon dioxide and releasing oxygen—occurs during photosynthesis. - To allow the plant to retain water and exchange gases, small pores (holes) in the leaves called stomata also evolved (Figure below). The stomata can open and close depending on weather conditions. When it's hot and dry, the stomata close to keep water inside of the plant. When the weather cools down, the stomata can open again to let carbon dioxide in and oxygen out. - A later adaption for life on land was the evolution of vascular tissue. Vascular tissue is specialized tissue that transports water, nutrients, and food in plants. In algae, vascular tissue is not necessary since the entire body is in contact with the water, and the water simply enters the algae. But on land, water may only be found deep in the ground. Vascular tissues take water and nutrients from the ground up into the plant, while also taking food down from the leaves into the rest of the plant. The two vascular tissues are xylem and phloem. Xylem is responsible for the transport of water and nutrients from the roots to the rest of the plant. Phloem carries the sugars made in the leaves to the parts of the plant where they are needed. Stomata are pores in leaves that allow gasses to pass through, but they can be closed to conserve water. - Plants evolved from freshwater green algae. - Plants have evolved several adaptations to life on land, including embryo retention, a cuticle, stomata, and vascular tissue. Use the resources below to answer the questions that follow. - The Role of Xylem Tissue and Stomata at http://www.youtube.com/watch?v=QBMkiLIyETc (3:34) - The Phloem at http://www.youtube.com/watch?v=M4onP3_4ERU (3:03) - In what groups of plants do you find xylem and phloem? Hint: refer to previous lesson if necessary. - What are the main components of sap? - Compare and contrast xylem and phloem. - What does each transport? - How are their structures similar? - What is "transpirational pull"? How is it key to the functioning of xylem? - How are plants different from green algae? How are they the same? - What is the purpose of vascular tissue? - How do plants prevent excess water loss? - Compare xylem to phloem. - What is the role of stomata?

http://www.ck12.org/life-science/plants-adaptations-for-life-on-land-in-life-science/lesson/Plants-Adaptations-for-Life-on-Land/

The deep-water canyons, seamounts, and underwater mountain ranges in the coastal waters of New England are gaining recognition for their importance to the health of fish populations like the struggling Atlantic cod. But these unique geological formations are also critical for the marine mammals that call the North Atlantic home. Hail the Whales The Atlantic coast is a veritable highway for migrating whales, which travel from breeding grounds in the south to feeding grounds in the north each year. But with many species facing reduced habitat, diminished populations, and increased boat traffic, this annual journey has become more and more difficult. These growing threats make areas of food abundance and shelter, such as Cashes Ledge and the New England Canyons and Seamounts, ever more critical to the success of migrating whales’ journeys. Cashes Ledge and the canyons and seamounts are unique in the Atlantic because their topography creates ideal conditions for plankton, zooplankton, and copepods – the main food for migrating minke, right, and humpback whales – to thrive. They also serve as spawning ground for larger food sources – including many squish, fish, and crustaceans. Altogether, this rich abundance of species adds up to a bountiful buffet for whales and other marine mammals. Sperm whales have often been spotted in the waters of seamounts, taking advantage of the reliable food, and Cashes Ledge serves as an oasis for hungry whales on their journey north. The healthy kelp domino effect These areas are not only crucial to whales; other marine mammals depend on them as well. Cashes Ledge boasts the largest coldwater kelp forest on the Atlantic seaboard, a habitat that creates ideal spawning grounds for cod, herring, and hake. The abundance of fish in turn feeds seals and porpoises, as well as whales. Scientists have noted a positive correlation between the size of an undersea kelp forest and populations of marine mammals, suggesting that more, healthy kelp means more marine mammals. That makes protecting areas with large kelp forests such as Cashes Ledge even more important. Even marine mammals that don’t visit Cashes Ledge itself still benefit from the protection of the area’s kelp forest, thanks to the “spillover effect:” Fish spawned in the shelter of the rocky crevasses and havens of the kelp forests disperse beyond Cashes Ledge and feed sea animals throughout the Gulf of Maine. Across the globe, underwater mountain and canyon habitats have proved to be important areas where marine mammals congregate to feed – and the canyons, seamounts, and ledges off the coast of New England are no different. Unfortunately, these important ecosystems are delicate and facing threats from harmful fishing gear and climate change. With so much at stake, it is vital to protect these places – not only for their inherent ecological value, but also so that they may sustain the mammals that depend on them. Imagine it’s 20 years from now, and your grandchild is about to head to bed – but first, she wants to hear a favorite bedtime story, “the one about the fish.” You pull it off the shelf – Mark Kurlansky’s The Cod’s Tale – and begin reading. Unbidden, her eyes widen at the vivid illustrations of the fish with a single chin whisker, at how it has millions of babies, and at how it gave birth to this country. Every time you read her the story, she asks the same question: “Can we go catch a cod tomorrow?” Every time, you have to tell her there aren’t any more cod in New England. And, every time she asks: “Why?” But you never really have a good answer for her. No Happy Ending in Sight for Cod The crisis in New England’s cod fishery was once again on the agenda at the New England Fishery Management Council’s December meeting in Portland, Maine. And once again, managers failed to take the basic actions needed for a concerted effort to restore this iconic fish. In addition to the collapse of the cod stock in the Gulf of Maine, New England is facing even greater declines of cod on Georges Bank, the historically important fishing area east of Cape Cod. The outlook for cod keeps getting worse, and the “actions” taken by the Council are so unlikely to make a difference that we must continue our call to save cod. The Worst of the Worst Some recent analyses have concluded that the cod population on Georges Bank is the lowest ever recorded – roughly 1 percent of what scientists would consider a healthy population. Other estimates put the population at only about 3 to 5 percent of the healthy target. The cod stock in the Gulf of Maine is hovering for the second year in a row at roughly 3 percent of the targeted healthy population. At its meeting last week, the Council did set new, lower catch limits for the severely depleted Georges Bank cod, but, true to form, those limits don’t go far enough. The Council is clearly in denial about the state of this fishery. If there is even a chance the number is 1 percent, this should be cause for major distress among Council members and fishermen alike. The Council’s actions (or, really, lack of action) leave me wondering, again, whether anyscience would ever be “enough” to compel them to halt the fishing of cod entirely. Habitat Loss Adds Fuel to the Fire Astoundingly, the Council also decided earlier this year to strip protection for important cod habitat on Georges Bank – amounting to a loss of some 81 percent of the formerly protected cod habitat. To recover, depleted fish populations need large areas protected from fishing and fishing gears; they need protected habitat where they can find food and shelter and reproduce; and they need large areas where female cod can grow old and reproduce prolifically. However, our fisheries managers – who are entrusted with safeguarding these precious resources for future generations as well as for current fishermen – ignore this science and continue to stubbornly deny the potential scope of this problem. This is an especially irresponsible stance in light of climate change. Not only are New England’s cod struggling to recover from decades of overfishing and habitat degradation, now the rapid rise in the region’s sea temperatures is further stressing their productivity. Protected habitats help marine species survive ecological stresses like warming waters. If a Cod Fish Dies But No One Records It, Did It Ever Really Exist? As if matters couldn’t get worse, the Council also voted to cut back significantly on the numbers of observers that groundfishing boats would have to have on-board to record what fish are actually coming up in their nets. This is little more than the Council’s blessing of unreported discards of cod and flounder and other depleted fish. We should be protecting more of these areas, not fewer; we should be doing more for these iconic fish, not less. So why is the Council making it so much harder for cod to recover? Perhaps it is simply contrary to human nature to expect the Council’s fishermen members to impose harsh measures on themselves when the benefits may only be seen by future generations. Perhaps federal fishery councils comprising active fishermen only work well with healthy fisheries. Federal officials at NOAA Fisheries will have the final say on these Council decisions to strip habitat protections, cutback on monitoring, and continue fishing on cod. We can only hope those officials will start taking the tough but necessary actions, giving New Englanders at least a semblance of hope that our grandchildren will be able to catch a codfish, not just read about one in a book. The Gulf of Maine is warming fast — faster than almost any other ocean area in the world. To say this is alarming is an understatement, and action is needed today to permanently protect large areas of the ocean, which scientists say is one of the best buffers against the disastrous effects of climate change. To that end, a diverse group of marine-oriented businesses, hundreds of marine scientists, aquaria, conservation organizations and members of the public are calling on the Obama administration to designate the Cashes Ledge Closed Area and the New England Coral Canyons and Seamounts as the first Marine National Monuments in the Atlantic. Conservation Law Foundation has worked for years to permanently protect the remarkable Cashes Ledge area. This biodiversity hotspot provides refuge for a stunning array of ocean wildlife — from cod to endangered right whales, bluefin tuna to Atlantic wolffish — and a rare lush kelp forest. The New England canyons and seamounts similarly shelter an incredible breadth of sea life, including spectacular ancient coral formations. Public support is widespread and growing. In September, more than 600 people attended a sold-out event hosted by the New England Aquarium and National Geographic Society where scientists discussed why these places are unique natural treasures. More than 160,000 people have electronically petitioned the president for monument protection. America has a long tradition of protecting our remarkable natural heritage and biological bounty. In contrast to our public lands and the Pacific Ocean, there are no areas in the Atlantic that are fully protected as national monuments. But why monument protection? Unlike fishery management closed areas or national marine sanctuaries, national monument designation protects against all types of commercial extraction that are harmful and can damage critical habitat: fishing, oil and gas exploration, sand and gravel mining, and more. Scientists say large-scale marine habitat protection is necessary to increase ocean resiliency in the face of climate change. Undisturbed underwater “laboratories” in places with relatively pristine habitats, like the Cashes Ledge area and the canyons and seamounts area, will be key in studying how — and how well — we are managing these already changing ocean ecosystems. These irreplaceable habitats can only play that role when protected in their entirety. Current protections by the New England Fishery Management Council are critical but not sufficient, as they are temporary, only limited to commercial fish species, and any coral protections are only discretionary. A monument designation protects all sea life and makes that protection permanent. It would be managed by scientists and others with ecological expertise (including but not limited to fisheries expertise). Fishery management councils were not designed and are not in the business of protecting scientifically unique and ecologically critical areas in the ocean. Permanent closure will also benefit collapsed fish populations like Atlantic cod, which would be able to rebuild and sustain themselves at healthy levels. Research is beginning to show that refuges could help struggling species like cod produce larger, older and significantly more productive females that could help recovery when their offspring eventually spill out to restock fishing in surrounding waters. The fishing industry is poised to benefit in the long term when commercially important fish are able to rebound. Protecting the few unique marine places we have left is good for the fishermen and communities that rely on a healthy and abundant ocean for their livelihoods and is our obligation to future generations. In recent weeks, we learned more sobering news for New England’s cod population. A paper published in Science detailed how rapidly increasing ocean temperatures are reducing cod’s productivity and impacting – negatively – the long-term rebuilding potential of New England’s iconic groundfish. The paper confirmed both the theoretical predictions associated with climate change and the recent scientific federal, state, and Canadian trawl surveys that reported a record-low number of cod caught in recent months. To be clear, the Science authors do not conclude that ocean temperature changes associated with climate change have caused the collapse of cod. We have management-approved overfishing of cod to thank for that. What rising ocean temperatures do seem to be doing, according to the Science paper, is dramatically changing the productivity of the remaining cod stocks. This makes it more difficult for cod to recover from overfishing today than at any other time in history, and perhaps reduces the ultimate recovery potential even if all fishing were halted. Stock assessments conducted without taking these productivity reductions into full account will dramatically overestimate cod populations and, in turn, fishing quotas. The Science paper is potentially very important, with major implications for fishing limits on cod for decades to come, But stock assessment scientists have warned for years that their recent models were likely overestimating the amount of cod actually in the water – and the corresponding fishing pressure the stock could withstand. Unfortunately, those warnings have fallen on deaf ears at the New England Fishery Management Council. In fact, the managers at the Council, dominated by fishermen and state fisheries directors with short-term economic agendas, could hardly have done more than they already have to jeopardize Atlantic cod’s future—climate change or not. Overfishing, a Weakened Gene Pool, and the Loss of Productive Female Fish As a result of chronic overfishing, New England’s cod population is likely facing what geneticists call a “population bottleneck,” meaning that the diversity of the remaining cod gene pool is now so greatly reduced that the fish that are left are less resilient to environmental stresses like increasing sea temperatures. Overfishing has also caused the collapse of the age structure of the cod populations by removing almost all of the larger, more reproductive females (also known as the Big, Old, Fat, Fecund Females, or BOFFFS). Scientists have previously warned that losing these old spawners is a problem for cod productivity, but this new research suggests that the potential damage from their elimination may be significantly greater than imagined as a result of poor, climate change–related ecological conditions. The Science paper hypothesizes that an underlying factor in the productivity decline of cod this past decade was the correlation between extremely warm spikes in ocean temperatures and the drop in zooplankton species that are critical to the survival of larval cod. With fewer zooplankton, fewer cod larvae make it to their first birthday. The impacts of this zooplankton decline on cod productivity, however, could be exacerbated by the loss of the BOFFFs. Here’s why: Cod start to spawn at three to four years old, but young females produce significantly fewer and weaker eggs and cod larvae than their older counterparts. Those elder female fish, on the other hand, produce larger, more viable eggs – sometimes exponentially more healthy eggs – over longer periods of time. If the older female cod population had still been plentiful, they might have produced larvae more capable of surviving variations in zooplankton abundance. Perhaps the continued presence of larger, older, spawning females to the south of New England (where there is no commercial cod fishery) is one of the reasons that the cod fishery in the nearby warm waters off New Jersey is healthier now than it has been in recent history. The Cod Aren’t Completely Cooked Yet: Four Potential Solutions Cod have been in trouble since the 1990s, and now climate change is magnifying these troubles. This new reality, however, is not cause for us to throw in the towel. There are actions that our fishery managers can take now that will make a difference. First, large cod habitat areas have to be closed to fishing – permanently. This is the only way to protect the large females and increase their number. Designating cod refuges such as the Cashes Ledge Closed Area as a marine national monument will remove the temptation for fishery councils – always under pressure to provide access to fish – to reopen them in the future. Such monuments would also sustain a critical marine laboratory where more of these complex interactions between cod and our changing ocean environment can be studied and understood. Second, managers need to gain a better understanding of the cod populations south of Cape Cod. While it is well and good to land “monster” female cod on recreational boat trips, those fish may be the key to re-populating Georges Bank. Caution, rather than a free-for-all, is the best course of action until the patterns of movement of those cod populations, as related to ocean temperature increases, are better understood. Third, as observed in the Science paper, stock assessment models as well as guidance from the Council’s Science and Statistical Committee must start incorporating more ecosystem variables and reflecting a more appropriate level of scientific precaution in the face of the reality of climate change shifts. Enough talk about scientific uncertainty and ecosystem-based fisheries management; action is needed, and science should have the lead in guiding that action. Finally, the importance of funding data collection and fishery science is evident from this important Science paper, which was supported by private, philanthropic dollars. NOAA should be undertaking this sort of work – but it is not in a position to even provide adequate and timely stock assessments, because limited funding forces the agency to use the existing outdated models. NOAA’s funding limitations are constraining both collection of the essential field data needed to understand our changing world as well as the analysis of that data into meaningful and appropriate management advice. If Congress can find $33 million to give fishermen for the most recent “groundfish disaster,” it ought to be able to find money to prevent such avoidable disasters in the future. Ultimately, the Science paper shines some much needed light on our climate change–related fishery issues in New England, but we can’t let it overshadow decades of mismanagement or justify a fatalistic attitude toward cod rebuilding. Steps can and must be taken, and fishery managers are still on the hook for the success or failure of our current and future cod stocks. One of the North Atlantic’s smallest ocean critters is making big waves in New England. Over the last decade, we’ve seen the collapse of our iconic Atlantic cod fishery due to extreme overfishing. Now, a new study is showing a potentially disastrous link between the effects of climate change and the ailing species’ chance of recovery. Warming waters are bound to be bad news for a cold water fish, but the problem goes much deeper than that, affecting the entire life cycle of the species. Some of this is due to tiny, microscopic creatures called zooplankton. So what are these little guys, and why are they so important? Zooplankton is a categorization of a type of ocean organism that includes various species, including Pseudocalanus spp, and Centropages typicus. These two species happen to be the major food source of larval cod in the Gulf of Maine. Zooplankton, which are usually smaller than 1/10 of an inch, play a major role in the Atlantic’s food web. When there are lots of them, things are pretty good. Young fish prey on them and grow to be healthy, adult fish. But when there aren’t enough plankton to go around, species like Atlantic cod can suffer. When cod larvae aren’t easily able to find the food they need to grow, fewer of them make it to their first birthday. And without lots of cod that survive to be at least 4 years old (the age at which females begin spawning), the recovery of the entire stock can stall. The stock needs larger, older, more productive females to thrive in order to have any hope for recovery. Warming and shifting But why would the plankton be in such short supply? This is where climate change comes in. According to NOAA, temperature changes can cause the redistribution of plankton communities. In the Gulf of Maine, scientists have found fewer plankton in the same areas where cod populations have been found to be struggling. The shifts in temperature lead to the displacement of a critical food source, making it difficult for young cod to survive. With the Gulf of Maine is warming faster than 99% of ocean areas, this is an enormously alarming problem. More temperature changes and the shifting of plankton populations could make it even harder for New England cod populations to return to healthy, sustainable levels. While the cod crisis is the result of many factors – but the loss of tiny zooplankton is a big problem. When considering how to best help cod stocks recover, fishery managers must take into account the effects of climate change, or else risk the total collapse of the species. In honor of Halloween, we’ve decided to highlight one of the more creepy looking fish that can be found in the waters off of New England. The monkfish (Lophius americanus), also known as goose-fish, anglerfish, and sea-devil, is considered a delicacy abroad, but until recently has been overlooked in America, perhaps due to its obtrusive appearance. The monkfish is highly recognizable, with its brown, tadpole-shaped body, and its gaping, fang-filled mouth. These eerie-looking fish can be found from Newfoundland to Georges Bank, and all the way down to North Carolina. They prefer to dwell on the sandy or muddy ocean-floor, where they feed on a variety of small lobsters, fish, and eels. Monkfish are typically found at depths of 230-330 feet, but have been caught in waters as deep as 2,700 feet; they have also been known to occasionally rise to the surface and consume small, unsuspecting birds. Females can grow up to forty inches and males up to thirty-five inches, and both can weigh up to seventy pounds. The average market size fish is around seventeen to twenty inches long. Before the 1960s, monkfish were considered to be undesirable bycatch. However, in the wake of the collapse of the New England Atlantic Cod fishery, the monkfish has slowly started to become a more common alternative, in part due to awareness campaigns about “underutilized species” in New England. Now, monkfish is caught to supply both international and domestic demand – the tail is prized for its firm texture and sweet taste, perfect for baking and poaching, and the liver is used in Japanese sushi. In fact, in the last two decades, fishing has increased so dramatically that monkfish stocks started to decline. Landings peaked in 1997 at sixty million pounds. However, thanks to the quick action of both the United States and Canada, a management plan was put in place and the stock population started to increase and stabilize. Landings now average around thirty-five million pounds annually. Monkfish are caught using trawls, gillnets, and dredges. The fishery is managed by the National Oceanic and Atmospheric Administration (NOAA), the New England Regional Fishery Management Council, and the Mid-Atlantic Fishery Management Council. These organizations do not impose annual catch limits, but do limit daily catches as well as limit access to the fishery. Nevertheless, the catch is still exceeding target catch levels in certain locations. Current threats to monkfish are common among New England marine species: warming temperatures, ocean acidification, and habitat loss. NOAA Fishwatch considers monkfish to be well managed and a “smart seafood choice” – however, it is still vulnerable, and the fishery should continue to be closely monitored, or it could suffer the same fate as other groundfish fisheries. So, if you are looking for a spooky-themed seafood dish for this weekend’s festivities, it might be time to give monkfish a try… It would also make one unique Halloween costume! October is National Seafood Month! To celebrate, I spoke with Andrea Tomlinson, General Manager of New Hampshire Community Seafood, an organization committed to supporting the state’s fishing industry and ensuring community access to fresh, locally caught seafood. We hear a lot about sustainable seafood in New England, but what does it really mean, and how can we, as consumers (and seafood lovers), impact the future of the fishing industry – all the while eating more healthy fish? AY: What is “sustainable seafood”? AT: I think few people understand what it really means – as more people use the term, it seems to have lost meaning. For me, sustainable seafood simply means that our fishermen are only taking an amount of a particular population that does not prevent parent fish from reproducing at the same level the following year. If fishermen leave the pregnant and older fish alone, and take just the younger fish, it’s more likely to be sustainable. The fish population must be able to sustain itself while also being fished for commercial purposes. AY: Do you think most of the industry fishes this way? AT: No. In the past, it was a free for all. Fishermen took whatever they wanted — cod was our fish, there was lots of it, so we took lots of it. Today, our small New England fishermen are still fishing the same amount (and taking the parent fish), but there are other, bigger players in the game. Once cod was shown to be a successful industry, the number of fishermen increased – and now the populations are suffering because of it. Our local fishermen never had to be conscious about [the amount they could catch] before. In order to stay in business, you want to take the biggest and most fish you can. When you take this traditional way of fishing and compound it with new catch regulations (and a perceived lack of communication from those enforcing the regulations), and more and bigger players fishing in the area, that’s how we ended up where we are today, with the fishing industry in crisis. AY: What are “underutilized fish” (formerly called “trash fish”) and how could they help the industry and/or economy? AT: In New England, there are certain types of fish that we have a lot of, but that just aren’t as popular as cod or haddock. There’s the dogfish shark, which is a shark but they are small – about three-and-a-half to four feet in length. In Europe, they are commonly used in fish and chips. Here in New England, we have lots of it. So much so, that they are almost considered overpopulated, making it a great alternative for consumers, especially since whatever you can do with cod, you can do with dogfish. AY: But it doesn’t have quite the same ring to it. AT: Right. When people hear “shark” and “dogfish,” they don’t like that. But as soon as you tell them how to prepare it, and that it holds up well in the freezer, and it seasons well, and is cheap – that makes a difference. AY: Are there other underutilized fish in New England? AT: There’s the King Whiting, a type of Silver Hake. It’s a delectable, thick, firm white fish that’s high in protein and omega-3s. It’s good for grilling or sautéing, and the fillet is just as large as one from a cod or haddock. And there’s also the Monkfish, which is an incredibly scary-looking fish on the outside – and delicious on the inside. We hear it called the “poor man’s lobster.” It tastes just like lobster, but for a fraction of the price. AY: How does a Community Supported Fishery work? Is this model feasible in other places? AT: The way fishing in New England works now, most fishermen sell everything they catch all at once at an auction, instead of buying directly “off the boat.” So, as a Community Supported Fishery, or CSF, New Hampshire Community Seafood gives the fishermen an incentive – we’d give them, say, an extra $0.25 per pound of a certain fish that’s higher than what they would receive at an auction. For dogfish, it’s actually a $1.10 per pound incentive! A CSF is really the only way to buy off the boat now. We buy a small portion of what the local fishermen catch, but it’s something. There are about 50 CSFs in the United States. On land, we’ve seen a growing popularity in supporting the local farmer, and this fits in well with that model. You pay up front, and get what’s ripe each week – it works the same way with fish. Community members can support local fishermen and the local economy in this way. So, the challenge is to get people to realize that underutilized fish are just as delicious as cod and haddock. In New England in particular, when people hear that the fishing industry is in crisis, that affects them. Many who grew up here are enamored by our iconic fishing traditions – maybe they have good memories of fishing, or they feel that it’s a big part of the culture. When you add in the “locavore” mentality, as well as those who are trying to eat healthier, we see a real opportunity to appeal to a lot of people. AY: So consumers can have a real impact here. AT: Yes. The fish are there – all we need is more consumers and more buyers, and it can make a greater impact. We are also working with restaurants and chefs; they will buy underutilized fish and put it on the menu, creating more exposure and making it easier for consumers to try something new. Right now we are in 10 restaurants and a hospital cafeteria, and are continuing to expand. AY: How can people get involved? AT: We are mostly based in Portsmouth, NH, but our CSF has 17 pickup locations in New Hampshire, one in Northern Massachusetts, and we’re partnering with Monadnock food cooperative in Keene, NH. (All of these are listed on the New Hampshire Community Seafood website). We also have a newsletter that informs locals about what’s new, how to cook underutilized fish, recipes, and more. AY: Anything else you would like to add? AT: Three years ago, there were 26 local fishermen in New Hampshire, and now there are only 9 left. We buy fish from all of them. The industry is in desperate need of support, both from communities and from the NMFS [regulators]. In addition to community-supported fishing organizations like NHCS, the Gulf of Maine Research Institute’s Out of the Blue series aims to educate the public about abundant fish that are well-managed and are not harvested primarily due to low market demand. And NOAA recently announced the public availability of fishwatch.gov, a resource that provides up-to-date information about fish, including the ability to look up a certain fish to see where it’s available, whether it’s a smart and sustainable option, nutrition information, and more. Would you (or have you) tried dogfish, whiting, or monkfish? Leave a comment below! If there’s one thing we can be sure of, it’s that New Englanders love lobster. It’s weaved into our culture and history, and it’s unimaginable to think of New England without this famed summer seafood. Few know that lobsters were once so plentiful in New England that Native Americans used them as fertilizer for their fields, and as bait for fishing. And before trapping was common, “catching” a lobster meant picking one up along the shoreline! During World War II, lobster was viewed as a delicacy, so it wasn’t rationed like other food sources. Lobster meat filled a demand for protein-rich sources, and continued to increase in popularity in post-war years, which encouraged more people to join the industry. Popular ever since, now when most people are asked what comes to mind when they think of New England, seafood – especially lobster – is typically at the top of the list. An industry under threat We love our New England lobster, but there’s evidence suggesting they’re in danger of moving away from their longtime home. That’s because lobster is under threat from climate change, the effects of which can already be seen on this particular species. The Gulf of Maine is warming faster than 99% of ocean areas. Until last winter’s uncharacteristically cold temperatures, the prior few years saw an increase in catchable lobster – as the warmer temperatures cause them to molt early, and they move toward inshore waters after molting. However, continued warming will ultimately encourage the lobsters to move north to find colder waters, where they spend the majority of their time. This is already happening in southern New England, where the industry is already suffering, seeing lobsters migrating northward. And we’re still learning about the potential for damage caused by ocean acidification, as well as how lobsters may be affected by an increase in colder than usual New England winters. As we celebrate one of New England’s iconic species on National Lobster Day, let’s remember that slowing down climate change is an important priority for ensuring that future generations can enjoy not Canadian or Icelandic lobster, but New England lobster. Click here to support Conservation Law Foundation’s efforts on fighting climate change.

http://newenglandoceanodyssey.org/category/talking-fish/

Constitution of Denmark |This article needs additional citations for verification. (June 2014)| |This article is part of a series on the politics and government of The Constitutional Act of Denmark (Danish: Danmarks Riges Grundlov) is the main part of the constitution of Kingdom of Denmark. First written in 1849, it establishes a sovereign state in the form of a constitutional monarchy, with a representative parliamentary system. The later sections of the Constitution guarantee fundamental human rights and lay out the duties of citizens. The current Constitution was signed on 5 June 1953 as "the existing law, for all to unswerving comply with, the Constitutional Act of Denmark". Idea and structure The main principle of the Constitution was to limit the monarch's power (section 2). The Constitution of 1849 established a bicameral parliament, the Rigsdag, consisting of the Landsting and the Folketing. It also secured civil rights, which remain in the current constitution, such as habeas corpus (section 71), private property rights (section 72) and freedom of speech (section 77). The Constitution is based on the separation of powers into the three branches of government, the legislative, the executive and the judicial branches. The Constitution is heavily influenced by the French philosopher Montesquieu, whose separation of powers was aimed at achieving mutual monitoring of each of the branches of government. This is achieved through the Constitution's section 3, although the division between legislative and executive power is not as sharp as in the United States. The original constitution of Denmark was signed on 5 June 1849 by King Frederick VII. The event marked the country's transition to constitutional monarchy, putting an end to the absolute monarchy that had been introduced in Denmark in 1660. The Constitution has been rewritten 4 times since 1849. Before the first constitutions, the power of the king was tempered by a håndfæstning, a charter each king had to sign before being accepted as king by the land things. This tradition was abandoned in 1665 when Denmark got its first constitution Lex Regia (The Law of The King, Danish: Kongeloven) establishing absolute power for King Frederick III of Denmark, and replacing the old feudal system. This is Europe's only formal absolutist constitution. Absolute power was passed along with a succession of Danish monarchs until Frederick VII, who agreed to sign the new constitution into law on 5 June 1849, which has since been a Danish national holiday. Frederick VII's immediate predecessor, his father Christian VIII, ruled Denmark from 1839 to 1848, and had been king of Norway until the political turmoil of 1814 forced him to abdicate after a constitutional convention. Those who supported similar constitutional reforms in Denmark were disappointed by his refusal to acknowledge any limitations to his inherited absolute power, and had to wait for his successor to put through the reforms. Ditlev Gothard Monrad, who became Secretary in 1848, drafted the first copy of the Constitution, based on a collection of the constitutions of the time, sketching out 80 paragraphs, whose basic principles and structure resembles the current constitution. The language of the draft was later revised by Secretary Orla Lehmann among others, and treated in the Constitutional Assembly of 1848 (Danish: Grundlovsudvalget af 1848). Sources of inspiration for the Constitution include the Constitution of Norway of 1814 and the Constitution of Belgium. The constitution's civil rights are based on the Constitution of the United States of 1787, especially the Bill of Rights. The government's draft was laid before the Constitutional Assembly of the Realm (Danish: Den Grundlovgivende Rigsforsamling), part of which had been elected on 5 October 1848, the remainder having been appointed by the King. The 152 members were mostly interested in the political aspects, the laws governing elections and the composition of the two chambers of Parliament. The Constitution was adopted during a period of strong national unity, namely the First Schleswig War, which lasted from 1848–1851. The Danish constitution has been written five times, in 1849, 1866, 1915, 1920 and 1953. No Danish constitution has ever been amended; each time, a new constitution replaced the existing constitution. According to section 88 of the 1953 Constitution, changes require a majority in two consecutive Parliaments: before and after a general election. In addition, the Constitution must pass a popular vote, with the additional demand that at least 40% of voting age population must vote in favour. The Constitution sets out only the basic principles, with more detailed regulation left over to the legislative branch of government, currently the Danish parliament Folketinget. The four changes can be summed up as follows: - In 1866, the defeat in the Second Schleswig War, and the loss of Schleswig-Holstein led to tightened election rules for the Upper Chamber, which paralyzed legislative work, leading to provisional laws. The conservative Højre had pressed for a new constitution, giving the upper chamber of parliament more power, making it more exclusive and switching power to the conservatives from the original long standing dominance of the National liberals, who lost influence and was later disbanded. This long period of dominance of the Højre party under the leadership of Jacob Brønnum Scavenius Estrup with the backing of the king Christian IX of Denmark was named the provisorietid (provisional period) because the government was based on provisional laws instead of parliamentary decisions. This also gave rise to a conflict with the Liberals (farm owners) at that time and now known as Venstre (Left). This constitutional battle concluded in 1901 with the so-called systemskifte (change of system) with the liberals as victors. At this point the king and Højre finally accepted parliamentarism as the ruling principle of Danish political life. This principle was not codified until the 1953 constitution. - In 1915, the tightening from 1866 was reversed, and women were given the right to vote. Also, a new requirement for changing the constitution was introduced. Not only must the new constitution be passed by two consecutive parliaments, it must also pass a referendum, where 45% of the electorate must vote yes. This meant that Prime Minister Thorvald Stauning's attempt to change the Constitution in 1939 failed. - In 1920, a new referendum was held to change the Constitution again, allowing for the reunification of Denmark following the defeat of Germany in World War I. This followed a referendum held in the former Danish territories of Schleswig-Holstein regarding how the new border should be placed. This resulted in upper Schleswig becoming Danish, today known as Southern Jutland, and the rest remained German. - In 1953, the fourth constitution abolished the Upper Chamber (the Landsting), giving Denmark a unicameral parliament. It also enabled females to inherit the throne (see Succession), but the change still favored boys over girls (this was changed by a referendum in 2009 so the first-born inherits the throne regardless of sex). Finally, the required number of votes in favor of a change of the Constitution was decreased to the current value of 40% of the electorate. The Constitution of Denmark outlines certain human rights in sections 71–80. Several of these are of only limited scope and thus serve as a sort of lower bar. The European Convention on Human Rights was introduced in Denmark by law on 29 April 1992 and supplements the mentioned paragraphs. Symbolic status of the king When reading the Danish Constitution, it is important to bear in mind that the King is meant to be read as the government because of the monarch's symbolic status. This is a consequence of sections 12 and 13, by which the King executes his power through his ministers, who are responsible for governing. An implication of these sections is that the monarch cannot act alone in disregard of the ministers, so the Danish monarch does not interfere in politics. Section 4 establishes that the Evangelical Lutheran Church is "the people's church" (folkekirken), and as such is supported by the state. Freedom of religion is granted in section 67, and official discrimination based on faith is forbidden in section 70. Christianity is a major religion Section 20 of the current constitution establishes that specified parts of national sovereignty can be delegated to international authorities if the Parliament or the electorate votes for it. This section has been debated heavily in connection with Denmark's membership of the European Union, as critics hold that changing governments have violated the Constitution by surrendering too much power. In 1996, Prime Minister Poul Nyrup Rasmussen was sued by 12 euroskeptics for violating this section. The Danish Supreme Court (Danish: Højesteret) acquitted Rasmussen (and thereby earlier governments dating back to 1972) but reaffirmed that there are limits to how much sovereignty can be surrendered. Other constitutional laws of Denmark The Danish constitution contains these additional parts: - The parts of Kongeloven, the former absolute monarchist constitution from 1665, that were not superseded. - The Act of Succession to the Danish Throne of 27 March 1953 also has status as a constitutional law, as it is directly referred to in Article 2 of the Constitutional Act. Therefore, amendments to the Act of Succession require adherence to the constitutional amendment procedure as provided for in Article 88 of the Danish Constitution Act. An amendment to abolish male preference to the throne (bill no. 1, Folketing session of 2005–06) was passed by a referendum in 2009. - To an extent the laws granting self government to the Faroe Islands and Greenland can be considered constitutional. - Certain particular customs, not explicitly referred to in the Constitutional Act itself, have been recognised as carrying constitutional legal weight (such as the right of the Finance Committee to authorise public expenditure outside of the national budget), also form part of Danish Constitutional law. - Constitutional law - Constitutional economics - Index of Denmark-related articles - Outline of Denmark - Danish Realm - "CIA World Factbook: Denmark: Government". Retrieved 8 July 2009. - Constitutional Act of Denmark, 5 June 1953 (WikiSource) - Folketinget Archived 7 May 2009 at the Wayback Machine - Grundloven 1849 by Erik Strange Petersen Aarhus University in danish - "Chronology". Constitute. Retrieved 29 April 2015. - CIA – The World Factbook - Søren Mørch: 24 statsministre. ISBN 87-02-00361-9. - Grundloven, Mikael Witte 1997 ISBN 87-7724-672-1 - The Constitutional Act of Denmark | Folketinget (Danish Parliament) |Wikisource has original text related to this article:|

https://en.wikipedia.org/wiki/Constitution_of_Denmark

You Are Here Activity 1: Forgiveness in History Activity time: 15 minutes Materials for Activity - Leader Resource 1, Truth and Reconciliation Match Ups - Leader Resource 2, Histories - Basket or box Preparation for Activity - Cut apart the names in Leader Resource 1, Truth and Reconciliation Match Ups. Put the slips of paper into a basket or box. - Cut apart the historical data in Leader Resource 2, Histories. Be prepared to distribute the individual case histories to volunteers for reading. Description of Activity Youth look at forgiveness on a large scale: nations or organizations seeking forgiveness for oppression of a group of people. Ask for a volunteer to look up the word "forgive" in the dictionary and read the definition to the group. Say, in your own words: We generally offer an apology to someone when we are seeking forgiveness or a pardon. This can be done by individuals. However, sometimes groups, even nations, issue apologies for wrongs committed against an entire group of people. Sometimes the apology is a long time coming. Sometimes, it includes reparations, which are payments for an injury or a wrong. Show the group the basket with the names from Leader Resource 1, Truth and Reconciliation Match Ups. Tell them that they are going to play a matching game. Everyone should take a slip of paper that has the name of one party of an apology. They need to find their counterpart. They will do this by asking other youth "yes and no questions" until they believe they have found their match. Assist youth as needed. After everyone has correctly found a match, ask for volunteers to read the case histories from Leader Resource 2, Histories. Lead a group discussion with questions such as: - Do you think every member of the oppressed group accepts the apologies? Why or why not? - How would you feel if a nation or organization issued an apology, but no reparations or other efforts to try to repair the damage? - How would you feel if the nation or organization offered reparation, but did not accept wrongdoing or offer an apology? - Why do you think the responsible parties are hesitant to accept responsibility or offer reparations? - Can you think of other cases where a government has addressed its previous wrongdoing? Including All Participants Be aware that youth who identify as a member of an oppressed group covered in the histories might be in the room. If you think this youth might find the activity difficult, delete that case history. However, do not assume that will be the case. Use your judgment, based on the experiences you have shared with the youth so far. You might also ask the youth beforehand.

http://www.uua.org/re/tapestry/youth/call/workshop11/173101.shtml

The Treaty of Paris (see Paris, Treaty of) formally recognized the new nation in 1783, although many questions were left unsettled. The United States was floundering through a postwar depression and seeking not too successfully to meet its administrative problems under the Articles of Confederation (see Confederation, Articles of). The leaders in the new country were those prominent either in the council halls or on the fields of the Revolution, and the first three Presidents after the Constitution of the United States was adopted were Washington, Adams, and Jefferson. Some of the more radical Revolutionary leaders were disappointed in the turn toward conservatism when the Revolution was over, but liberty and democracy had been fixed as the highest ideals of the United States. The American Revolution had a great influence on liberal thought throughout Europe. The struggles and successes of the youthful democracy were much in the minds of those who brought about the French Revolution, and most assuredly later helped to inspire revolutionists in Spain's American colonies. Sections in this article: The Columbia Electronic Encyclopedia, 6th ed. Copyright © 2012, Columbia University Press. All rights reserved. See more Encyclopedia articles on: U.S. History

http://www.infoplease.com/encyclopedia/history/american-revolution-aftermath.html

Forests play an important an important role in climate change. The destruction and degradation of forests contributes to the problem through the release of CO2. But the planting of new forests can help mitigate against climate change by removing CO2 from the atmosphere. Combined with the sun's energy, the captured carbon is converted into trunks, branches, roots and leaves via the process of photosynthesis. It is stored in this "biomass" until being returned back into the atmosphere, whether through natural processes or human interference, thus completing the carbon cycle. Tree planting and plantation forestry are well established both in the private and public sectors. The most recent data released by the UN's Food and Agriculture Organisation suggest that plantation forests comprised an estimated 7% of global forest area in 2010. Most of these forests were established in areas that were previously not under forest cover, at least in recent years. Trees are also planted as part of efforts to restore natural forests as well as in agroforestry, which involves increasing tree cover on agricultural land and pastures. Under certain conditions plantations can grow relatively fast, thus absorbing CO2 at higher rates than natural forests. In the absence of major disturbances, newly planted or regenerating forests can continue to absorb carbon for 20–50 years or more. In comparison to preventing the loss of natural forests, however, tree planting has the potential to make only a limited contribution to reducing CO2 levels in the atmosphere. In 2000, the IPCC gathered the available evidence for a special report which concluded that tree-planting could sequester (remove from the atmosphere) around 1.1–1.6 GT of CO2 per year. That compares to total global greenhouse gas emissions equivalent to 50 GT of CO2 in 2004. Unlike measures to reduce deforestation, tree planting and reforestation were included as activities eligible for finance under the Kyoto protocol. Kyoto's rules and procedures, however, restricted the scale and scope of these activities. As a result, projects have struggled to get off the ground and the carbon sequestered has been almost negligible. Outside of Kyoto, some tree-planting projects established to absorb CO2 have turned out to be nonviable due to the cost of acquiring inputs or protecting young trees from fire, drought, pests or diseases. The cost of land is another barrier to widespread tree-planting, especially where there is competition with other land uses such as food or biofuel production. As negotiations over the future of Kyoto continue, the extent of the possible role of tree planting in a future climate change framework remains unclear. Tree planting is, however, unlikely to be implemented on a scale to reach even the relatively modest potential contribution outlined by the IPPC – especially in the absence of a high carbon price. • This article was written by Dr Charles Palmer of the Grantham Research Institute on Climate Change and the Environment at LSE in collaboration with the Guardian The ultimate climate change FAQ This editorial is free to reproduce under Creative Commons This post by The Guardian is licensed under a Creative Commons Attribution-No Derivative Works 2.0 UK: England & Wales License. Based on a work at theguardian.com

http://www.theguardian.com/environment/2012/nov/29/planting-trees-climate-change?view=mobile

Swedish is descended from Old Norse. Compared to its progenitor, Swedish grammar is much less characterized by inflection. Modern Swedish has two genders and no longer conjugates verbs based on person or number. Its nouns have lost the morphological distinction between nominative and accusative cases that denoted grammatical subject and object in Old Norse in favor of marking by word order. Swedish uses some inflection with nouns, adjectives, and verbs. It is generally a subject–verb–object (SVO) language with V2 word order. - 1 Nouns - 2 Pronouns - 3 Adjectives - 4 Comparison - 5 Numerals - 6 Verbs - 7 Adverbs - 8 Prepositions - 9 Syntax - 10 Notes - 11 References - 12 External links Nouns have two grammatical genders: common (utrum) and neuter (neutrum), which determine their definite forms as well as the form of any adjectives used to describe them. Noun gender is largely arbitrary and must be memorized; however, around three quarters of all Swedish nouns are common gender. Living beings are often common nouns, like in en katt, en häst, en fluga, etc. Swedish once had three genders—masculine, feminine and neuter. Though traces of the three-gender system still exist in archaic expressions and certain dialects, masculine and feminine nouns have today merged into the common gender. A remnant of the masculine gender can still be expressed in the singular definite form of adjectives according to natural gender (male humans), in the same way as personal pronouns, han/hon, are chosen for representing nouns in Contemporary Swedish (male/female humans and optionally animals). There are traces of the former four-case system for nouns evidenced in that pronouns still have a subject, object (based on the old accusative and dative form) and genitive forms. Nouns make no distinction between subject and object forms, and the genitive is formed by adding -s to the end of a word. This -s genitive functions more like a clitic than a proper case and is nearly identical to the possessive suffix used in English. Note, however, that in Swedish this genitive s is appended directly to the word and is not preceded by an apostrophe. Swedish nouns are inflected for number and definiteness and can take a genitive suffix. They exhibit the following morpheme order: |Noun stem||(Plural)||(Definite article)||(Genitive -s)| Nouns form the plural in a variety of ways. It is customary to classify Swedish nouns into five declensions based on their plural indefinite endings: -or, -ar, -er, -n, and unchanging nouns. - Nouns of the first declension are all of the common gender. The majority of these nouns end in -a in the singular and replace it with -or in the plural. For example: en flicka (a girl), flickor (girls). A few nouns of the first declension end in a consonant, such as: en våg (a wave), vågor (waves); en ros (a rose), rosor (roses). - Nouns of the second declension are also of the common gender, with the exception of finger (finger). They all have the plural ending -ar. Examples include: en arm (an arm), armar (arms); en hund (a dog), hundar (dogs); en sjö (a lake), sjöar (lakes); en pojke (a boy), pojkar (boys); en sjukdom (an illness), sjukdomar (illnesses); en främling (a stranger), främlingar (strangers). A few second declension nouns have irregular plural forms, for instance: en afton (an evening), aftnar (evenings); en sommar (a summer), somrar (summers). - The third declension includes both common and neuter nouns. The plural ending for nouns of this declension is -er or, for some nouns ending in a vowel, -r. For example: en park (a park), parker (parks); ett museum (a museum), museer (museums); en sko (a shoe), skor (shoes); en fiende (an enemy), fiender (enemies). Some third declension nouns modify their stem vowels in the plural: en hand (a hand), händer (hands); ett land (a country), länder (countries); en bok (a book), böcker (books). - All nouns in the fourth declension are of the neuter gender and end in a vowel in the singular. Their plural ending is -n. For example: ett bi (a bee), bin (bees); ett äpple (an apple), äpplen (apples). Two nouns in this declension have irregular plural forms: ett öga (an eye), ögon (eyes); ett öra (an ear), öron (ears). - Fifth declension nouns have no plural ending and they can be of common or neuter gender. Examples of these include: ett barn (a child), barn (children); ett djur (an animal), djur (animals); en lärare (a teacher), lärare (teachers). Some fifth declension nouns show a vowel change in the plural: en mus (a mouse), möss (mice); en gås (a goose), gäss (geese); en man (a man), män (men). Articles and definite forms The definite article in Swedish is mostly expressed by a suffix on the head noun, while the indefinite article is a separate word preceding the noun. This structure of the articles is shared by the Scandinavian languages. Articles differ in form depending on the gender and number of the noun. The indefinite article, which is only used in the singular, is "en" for common nouns, and "ett" for neuter nouns, e.g. en flaska (a bottle), ett brev (a letter). The definite article in the singular is generally the suffixes "-en" or "-n" for common nouns (e.g. flaskan "the bottle"), and "-et" or "-t" for neuter nouns (e.g. brevet "the letter"). The definite article in the plural is "-na", "-a" or "-en", depending on declension group, for example flaskorna (the bottles), breven (the letters). When an adjective or numeral is used in front of a noun with the definite article, an additional definite article is placed before the adjective(s). This additional definite article is det for neuter nouns, den for common nouns, and de for plural nouns, e.g. den nya flaskan (the new bottle), det nya brevet (the new letter), de fem flaskorna (the five bottles). A similar structure involving the same kind of circumfixing of the definite article with the words där (there) or här (here) is used to mean "this" and "that", e.g. den här flaskan (this bottle), det där brevet (that letter) as a demonstrative article. The five declension classes may be named -or, -ar, -er, -n, and null after their respective plural indefinite endings. Each noun has eight forms: singular/plural, definite/indefinite and caseless/genitive. The caseless form is sometimes referred to as nominative, even though it is used for grammatical objects as well as subjects. The genitive is always formed by appending -s to the caseless form. In the second, third and fifth declensions words may end with an -s already in the caseless form. These words take no extra -s in genitive use: the genitive (indefinite) of hus ("house") is hus. Morpheme boundaries in some forms may be analyzed differently by some scholars. The Swedish genitive is not considered a case by all scholars today,[who?] as the -s is usually put on the last word of the noun phrase even when that word is not the head noun, mirroring English usage (e.g. Mannen som står där bortas hatt. "The man standing over there's hat."). This use of -s as a clitic rather than a suffix has traditionally been regarded as ungrammatical, but are today dominant to the point where putting an -s on the head noun is considered old fashioned. The Swedish Language Council recommends putting the ending after the phrase, except when making temporary constructions, where one should instead try to reformulate. These examples cover all regular Swedish caseless noun forms. First declension: -or (common gender) Second declension: -ar (common gender) Third declension: -er, -r (mostly common gender nouns, some neuter nouns) Words taking only -r as a marker for plural is regarded as a declension of its own by some scholars. However, traditionally these have been regarded as a special version of the third declension. Fourth declension: -n (neuter) This is when a neuter noun ends in a vowel. Fifth declension: unmarked plural (mostly neuter nouns that don't end in vowels and common gender nouns ending in certain derivation suffixes) The Swedish personal pronoun system is almost identical to that of English. Pronouns inflect for person, number, and, in the third person singular, gender. Swedish is different, inter alia, as it has a separate third-person reflexive pronoun sig (himself, herself, itself, themselves) analogous to French se, and distinct 2nd person singular forms du ("thou") and ni ("you", formal/respectful), and their objective forms, which have all merged to "you" in English, while the third person plurals are becoming merged in Swedish instead[clarification needed]. Some aspects of personal pronouns are simpler in Swedish: reflexive forms are not used for the first and second person, although själv ("self") and egen/eget/egna ("own") may be used for emphasis, and there are no absolute forms for the possessive. The Swedish personal pronouns are: |Person||Nominative||Objective||Possessive: com./neut./pl.||Person||Nominative||Objective||Possessive: com./neut./pl.| |2 (familiar)||du||dig||din/ditt/dina1||2 (formal: sg. or pl.)||ni3||er||er/ert/era1 |3 Gen-Neu. (neologism)||hen4||hen/henom4||hens4| |3 Indef.||man ("one", Fr. "on")||en||ens| |(3 Refl.)||—||sig||sin/sitt/sina1||(3 Refl.)||—||sig||sin/sitt/sina1| 1These possessive pronouns are inflected similarly to adjectives, agreeing in gender and number with the item possessed. The other possessive pronouns (i.e. those listed without slashes) are genitive forms that are unaffected by the item possessed. 2de (they) and dem (them) are both usually pronounced "dom" (/dɔm/) in colloquial speech, while in formal speech, "dom" may optionally replace just "dem". In some dialects (especially Finnish ones) there is still a separation between the two; de is then commonly pronounced /di/. Also, mig, dig, sig are pronounced as if written "mej", "dej", "sej", and are also sometimes spelled that way in less formal writing or to signal spoken language, but this is not appreciated by everyone. 3ni is derived from an older pronoun I, "ye", for which verbs were always conjugated with the ending -en. I became ni when this conjugation was dropped; thus the n was moved from the end of the verb to the beginning of the pronoun. 4hen and its inflections are neologisms: they are gender-neutral pronouns used by some to avoid a preference for female or male, when a person's gender is not known, or to refer to people whose gender is not defined as female or male. They are relatively new in widespread use, but since 2010 have appeared frequently in traditional and online media, legal documents, and literature. The use of these words has prompted a political and linguistic debate in Sweden, and their use is not universally accepted by Swedish speakers. Demonstrative, interrogative, and relative pronouns - including related words not strictly considered pronouns - den här, det här, de här: this, these (may qualify a noun in the definite form.) - den där, det där, de där: that, those (may qualify a noun in the definite form.) - denne/denna/detta/dessa: this/these (may qualify a noun in the indefinite form.) - som: as, that, which, who (strictly speaking, a subordinating conjunction rather than a pronoun, som is used as an all-purpose relative pronoun whenever possible in Swedish.) - vem: who, whom (interrogative) - vilken/vilket/vilka: which, what, who, whom, that - vad: what - vems: whose (interrogative) - vars: whose (relative) - när: when - då1: then, when (relative) - här, där, var1: here, there, where (also form numerous combinations such as varifrån, "where from", and därav, "thereof".) - hit, dit, vart1: hither, thither, whither (not archaic as in English) - vem som helst, vilket som helst, vad som helst, när som helst, var som helst: whoever, whichever, whatever, whenever, wherever, etc. - hädan, dädan, vadan, sedan1: hence, thence, whence, since (The contractions hän and sen are common. These are all somewhat archaic and formal-sounding except for sedan.) - någon/något/några, often contracted to and nearly always said as nån/nåt/nåra2: some/any, a few; someone/anyone, somebody/anybody, something/anything (The distinction between "some" in an affirmative statement and "any" in a negative or interrogative context is actually a slight difficulty for Swedes learning English.) - ingen/inget/inga2: no, none; no one, nobody, nothing - annan/annat/andra: other, else - någonstans, ingenstans, annanstans, överallt: somewhere/anywhere, nowhere, elsewhere, everywhere; (more formally någonstädes, ingenstädes, annorstädes, allestädes) - någorlunda, ingalunda, annorlunda: somehow/anyhow, in no wise, otherwise - någonting, ingenting, allting: something/anything, nothing, everything 1 då, där, dit, and dädan, (then, there, thither, and thence,) and any compounds derived from them are used not only in a demonstrative sense, but also in a relative sense, where English would require the "wh-" forms when, where, whither and whence. 2 Animacy is implied by gender in these pronouns: non-neuter implies a person (-one or -body) and neuter implies a thing. Swedish adjectives are declined according to gender, number, and definiteness of the noun. In singular indefinite, the form used with nouns of the common gender is the undeclined form, but with nouns of the neuter gender a suffix -t is added. In plural indefinite an -a suffix is added irrespective of gender. This constitutes the strong adjective inflection, characteristic of Germanic languages: |Common||en stor björn, a large bear||stora björnar, large bears| |Neuter||ett stort lodjur, a large lynx||stora lodjur, large lynxes| In standard Swedish, adjectives are inflected according to the strong pattern, by gender and number of the noun, in complement function with är, is, such as - lodjuret är skyggt, the lynx is shy, and - björnarna är bruna, the bears are brown. In some dialects of Swedish, the adjective is uninflected in complement function with är, so becoming: - lodjuret är skygg, the lynx is shy, and - björnarna är brun, the bears are brown. In definite form we instead have a weak adjective inflection, originating from a Proto-Germanic nominal derivation of the adjectives. The adjectives now invariably take on an -a suffix irrespective of case and number, which was not always the case, cf. Proto-Germanic adjectives: |Common||den stora björnen, the large bear||de stora björnarna, the large bears| |Neuter||det stora lodjuret, the large lynx||de stora lodjuren, the large lynxes| As the sole exception to this -a suffix is that naturally masculine nouns (replaceable with han/honom) take the -e ending in singular. Colloquially, however, the usual -a-ending is possible in these cases in some Swedish dialects: |den store mannen, the large man||de stora männen, the large men| |den stora mannen, the large man| Adjectives with comparative and superlative forms ending in -are and -ast, which is a majority, also, and so by rule, use the -e suffix for all persons on definite superlatives: den billigaste bilen ("the cheapest car"). Another instance of -e for all persons is the plural forms and definite forms of adjectival verb participles ending in -ad: en målad bil ("a painted car") vs. målade bilar ("painted cars") and den målade bilen ("the painted car"). The cardinal numbers from zero to twelve in Swedish are: The number 1 is the same as the indefinite article, and its form (en/ett) depends on the gender of the noun that it modifies. The Swedish numbers from 13 to 19 are: The form aderton is archaic, and is nowadays only used in poetry and some official documents. It is still common in Finland Swedish. The numbers for multiples of ten from 20 to 1000 are: |tjugo||trettio||fyrtio||femtio||sextio||sjuttio||åttio||nittio||(ett) hundra||(ett) tusen| When trettio (30), fyrtio (40), femtio (50), sextio (60), sjuttio (70), åttio (80), nittio (90) are compounded with another digit, they form a compound number. In some dialects, numbers are not always pronounced the way they are spelled. With the numbers nio (9), tio (10) and tjugo (20), the -o is often pronounced as an -e, e.g. /tjuge/. In some northern dialects the -o is pronounced as a /-u/, /tjugu/, and in some middle dialects the -o is pronounced as an /-i/, /tjugi/. In spoken language, tjugo usually drops the final syllable when compounded with another digit and is pronounced as /tju/ + the digit, e.g. tjugosju (27) may be pronounced /tjusju/. Words ending in -io (trettio, fyrtio, etc.) are most often pronounced without the final -o. The "y" in fyrtio (40) is always pronounced as an /ö/. The ett preceding hundra (100) and tusen (1000) is optional, but in compounds it is usually required. Higher numbers include: |1 000 000||en miljon| |10 000 000||tio miljoner| |100 000 000||(ett) hundra miljoner| |1 000 000 000||en miljard ¹| ¹ Swedish uses the long scale for large numbers. The cardinal numbers from miljon and larger are true nouns and take the -er suffix in the plural. They are separated in written Swedish from the preceding number. Any number can be compounded by simply joining the relevant simple cardinal number in the same order as the digits are written. Written with digits, a number is separated with a space between each third digit from the right. The same principle is used when a number is written with letters, although using letters becomes less common the longer the number is. However, round numbers, like tusen, miljon and miljard are often written with letters as are small numbers (below 20). Numbers between 21-99 are written in the following format: (big number)(small number) example: 63 - "sextiotre" 48 - "fyrtioåtta" (note that the a is taken out between numbers 40-49) although 30-39 are slightly special, an extra t is added to these numders: 31 - "trettioett" 33 - "trettiotre" |Written form||In components (do not use in written Swedish)| |21||tjugoett / tjugoen||(tjugo-ett) / (tjugo-en)| |1 975||ettusen niohundrasjuttifem |10 874||tiotusen åttahundrasjuttifyra |100 557||etthundratusen femhundrafemtisju |1 378 971||en miljon trehundrasjuttiåtta tusen niohundrasjuttiett en miljon trehundrasjuttioåtta tusen niohundrasjuttioett |(en miljon tre-hundra-sjuttio-åtta tusen nio-hundra-sjuttio-ett)| The decimal point is written as "," (comma) and written and pronounced komma. The digits following the decimal point may be read individually or as a pair if there are only two. When dealing with monetary amounts (usually with two decimals), the decimal point is read as och, i.e. "and": 3,50 (tre och femtio), 7,88 (sju och åttioåtta). Rational numbers are read as the cardinal number of the numerator followed by the ordinal number of the denominator compounded with -del or -delar (part(s)). If the numerator is more than one, logically, the plural form of del is used. For those ordinal numbers that are three syllables or longer and end in -de, that suffix is usually dropped in favour of the de in -del. There are a few exceptions. |1⁄2||en halv, one half| |1⁄8||en åttondel or en åttondedel| |8⁄9||åtta niondelar or åtta niondedelar| |1⁄10||en tiondel or en tiondedel| |1⁄13||en trettondel or en trettondedel| |1⁄14||en fjortondel or en fjortondedel| |1⁄15||en femtondel or en femtondedel| |1⁄16||en sextondel or en sextondedel| |1⁄17||en sjuttondel or en sjuttondedel| |1⁄18||en artondel or en artondedel| |1⁄19||en nittondel or en nittondedel| |1⁄20||en tjugondel or en tjugondedel| First to twelfth: Thirteen to nineteen: - As cardinal numerals, but with the suffix -de, e.g., trettonde (13:e), fjortonde (14:e). Even multiples of ten (20th to 90th): - As cardinal numerals, but with the suffix -nde, e.g., tjugonde (20:e), trettionde (30:e) - As cardinal numerals, but with the suffix -de, e.g., hundrade (100:e, hundredth), tusende (1000:e, thousandth) - As cardinal numerals, but with the suffix -te, e.g., miljonte (millionth). There is no ordinal for "miljard" (billion). Verbs do not inflect for person or number in modern standard Swedish. They inflect for the present and past tense and imperative, subjunctive, and indicative mood. Other tenses are formed by combinations of auxiliary verbs with infinitives or a special form of the participle called the "supine". In total there are 6 spoken active-voice forms for each verb: infinitive, imperative, present, preterite/past, supine, and past participle. The only subjunctive form used in everyday speech is vore, the past subjunctive of vara ("to be"). It is used as one way of expressing the conditional ("would be", "were"), but is optional. Except for this form, subjunctive forms are considered archaic. Verbs may also take the passive voice. The passive voice for any verb tense is formed by appending -s to the tense. For verbs ending in -r, the -r is first removed before the -s is added. Verbs ending in -er often lose the -e- as well, other than in very formal style: stärker ("strengthens") becomes stärks or stärkes ("is strengthened") (exceptions are monosyllabic verbs and verbs where the root ends in -s). Swedish uses the passive voice more frequently than English. Swedish verbs are divided into four groups: |1||regular -ar verbs| |2||regular -er verbs| |3||short verbs, end in -r| |4||strong and irregular verbs, end in -er or -r| About 80% of all verbs in Swedish are group 1 verbs, which is the only productive verb group. Swenglish variants of English verbs can be made by adding -a to the end of an English verb, sometimes with minor spelling changes. The verb is then treated as a group 1 verb. Examples of modern loan words within the IT field are chatta and surfa. Swenglish variants from the IT field that may be used but are not considered Swedish include maila, mejla ([ˈmejˌla], to email or mail) and savea, sejva ([ˈsejˌva] to save). The stem of a verb is based on the present tense of the verb. If the present tense ends in -ar, the -r is removed to form the stem, e.g., kallar → kalla-. If the present tense ends in -er, the -er is removed, e.g., stänger → stäng-. For short verbs, the -r is removed from the present tense of the verb, e.g., syr → sy-. The imperative is the same as the stem. For group 1 verbs, the stem ends in -a, the infinitive is the same as the stem, the present tense ends in -r, the past tense in -de, the supine in -t, and the past participle in -d, -t, and de. For group 2 verbs, the stem ends in a consonant, the infinitive ends in -a, and the present tense in -er. Group 2 verbs are further subdivided into group 2a and 2b. For group 2a verbs, the past tense ends in -de and the past participle in -d, -t, and -da. For group 2b verbs, the past tense ends in -te and the past participle in -t, -t, and -ta. This is in turn decided by whether the stem ends in a voiced or a voiceless consonant. E.g. The stem of Heta (to be called) is het, and as t is a voiceless consonant the past tense ends in -te, making hette the past tense. If the stem ends in a voiced consonant however, as in Stör-a (to disturb), the past tense ends in -de making störde the past tense. For group 3 verbs, the stem ends in a vowel that is not -a, the infinitive is the same as the stem, the present tense ends in -r, the past tense in -dde, the supine in -tt, and the past participle in -dd, -tt, and -dda. Group 4 verbs are strong and irregular verbs. Many commonly used verbs belong to this group. For strong verbs, the vowel changes for the past and often the supine, following a definite pattern, e.g., stryka is a strong verb that follows the u/y, ö, u pattern (see table below for conjugations). Irregular verbs, such as vara (to be), are completely irregular and follow no pattern. As of lately, an increasing number of verbs formerly conjugated with a strong inflection has been subject to be conjugated with its weak equivalent form in colloquial speech. |to strike out |4 (irregular)||var-||var!||vara||är||var||varit||-||to be| *often new vowel Examples of tenses with English translations |Infinitive||To work||(Att) arbeta| |Present Tense||I work||Jag arbetar| |Past Tense, Imperfect Aspect||I worked||Jag arbetade| |Past Tense, Perfect Aspect||I have worked||Jag har arbetat| |Future Tense, Futurum Simplex||I will work||Jag ska arbeta| The irregular verb gå |Infinitive||To walk||(Att) gå| |Present Tense||I walk||Jag går| |Past Tense, Imperfect Aspect||I walked||Jag gick| |Past Tense, Perfect Aspect||I have walked||Jag har gått| |Future Tense, Futurum Simplex||I will walk||Jag ska gå| As in all Germanic languages, strong verbs change their vowel sounds in the various tenses. For most Swedish strong verbs that have a verb cognate in English or German, that cognate is also strong. For example, "to bite" is a strong verb in all three languages as well as Dutch: |Swedish||bita||jag biter||jag bet||jag har bitit||biten, bitet, bitna| |Dutch||bijten||ik bijt||ik beet||ik heb gebeten||gebeten| |German||beißen||ich beiße||ich biss||ich habe gebissen||gebissen| |English||to bite||I bite||I bit||I have bitten||bitten| The supine (supinum) form is used in Swedish to form the composite past form of a verb. For verb groups 1-3 the supine is identical to the neuter form of the past participle. For verb group 4, the supine ends in -it while the past participle's neuter form ends in -et. Clear pan-Swedish rules for the distinction in use of the -et and -it verbal suffixes were missing though before the first official Swedish Bible translation, completed 1541. This is best shown by example: - Simple past: I ate (the) dinner - Jag åt maten (using preterite) - Composite past: I have eaten (the) dinner - Jag har ätit maten (using supine) - Past participle common: (The) dinner is eaten - Maten är äten (using past participle) - Past participle neuter: (The) apple is eaten - Äpplet är ätet - Past participle plural: (The) apples are eaten - Äpplena är ätna The supine form is used after ha (to have). In English this form is normally merged with the past participle, or the preterite, and this was formerly the case in Swedish, too (the choice of -it or -et being dialectal rather than grammatical); however, in modern Swedish, they are separate, since the distinction of -it being supine and -et being participial was standardised. The passive voice in Swedish is formed in one of four ways: - add an -s to the infinitive form of the verb - use a form of bli (become) + the perfect participle - use a form of vara (be) + the perfect participle - use a form of få (get) + the perfect participle Of the first three forms, the first (s-passive) tends to focus on the action itself rather than the result of it. The second (bli-passive) stresses the change caused by the action. The third (vara-passive) puts the result of the action in the centre of interest: - Dörren målas. (Someone paints the door right now.) - Dörren blir målad. (The door is being painted, in a new colour or wasn't painted before.) - Dörren är målad. (The door is painted.) The fourth form is different from the others, since it is analogous to the English "get-passive": Han fick dörren målad (He got the/his door painted). This form is used when you want to use a subject other than the "normal" one in a passive clause. In English you could say: "The door was painted for him", but if you want "he" to be the subject you need to say "He got the door painted." Swedish uses the same structure. The subjunctive mood is rarely used in modern Swedish and is mostly limited to fixed expressions like leve kungen, "long live the king". Present subjunctive is formed by adding the "-e" ending to the stem of a verb: |att tala, to speak||talar, speak(s)||tale, may speak| |att bliva, to become||bli(ve)r, become(s)||blive, may become| |att skriva, to write||skriver, write(s)||skrive, may write| |att springa, to run||springer, run(s)||sprunge, may run| Adjectival adverbs are formed by putting the adjective in neuter singular form. Adjectives ending in -lig may take either the neuter singular ending or the suffix -en, and occasionally -ligen is added to an adjective not already ending in -lig. |tjock, thick||tjockt, thick||tjockt, thickly| |snabb, fast||snabbt, fast||snabbt, fast| |avsiktlig, intentional||avsiktligt, intentional||avsiktligt, avsiktligen, intentionally| |stor, great, large||stort, great, large||storligen, greatly i stort sett, largely Adverbs of direction in Swedish show a distinction that is lacking in English: some have different forms exist depending on whether one is heading that way, or already there. For example: - Jag steg upp på taket. Jag arbetade där uppe på taket. - I climbed up on the roof. I was working up there on the roof. |Heading that way||Already there||English| Unlike in more conservative Germanic languages (e.g. German), putting a noun into a prepositional phrase doesn't alter its inflection, case, number or definiteness in any way, except for in a very small number of set phrases. Prepositions of location |på||on||Råttan dansar på bordet.||The rat dances on the table.| |under||under||Musen dansar under bordet.||The mouse dances under the table.| |i||in||Kålle arbetar i Göteborg.||Kålle works in Gothenburg.| |vid||by||Jag är vid sjön.||I am by the lake.| |till||to||Ada har åkt till Göteborg.||Ada has gone to Gothenburg.| Prepositions of time |på||at||Vi ses på rasten.||See you at the break.| |före||before||De var alltid trötta före rasten.||They were always tired before the break.| |om||in||Kan vi ha rast om en timme?||May we have a break in one hour?| |i||for||Kan vi ha rast i en timme?||May we have a break for one hour?| |på||for (in a negative statement)||Vi har inte haft rast på två timmar.||We have not had a break for two hours.| |under||during||Vi arbetade under helgdagarna.||We worked during the holidays.| The general rule is that prepositions are placed before the word they are referring to. However, there are a few ambipositions that may appear on either side of the head: |Adposition||Meaning||Succeeding adposition (postposition)||Preceding adposition (preposition)||Translation| |runt||around||riket runt||runt riket||around the Kingdom| |emellan||between||bröder emellan||emellan bröder||between brothers| |igenom||through||natten igenom||igenom natten||the night through / through the night| Being a Germanic language, Swedish syntax shows similarities to both English and German. Like English, Swedish has a subject–verb–object basic word order, but like German, utilizes verb-second word order in main clauses, for instance after adverbs, adverbial phrases and dependent clauses. Adjectives generally precede the noun they determine, though the reverse is not infrequent in poetry. Nouns qualifying other nouns are almost always compounded on the fly (as with German, but less so with English); the last noun is the head. A general word-order template may be drawn for a Swedish sentence, where each part, if it does appear, appears in this order. (Source—Swedish For Immigrants level 3). |Fundament||Finite verb||Subject (if not fundament)||Clausal Adverb/negation||Non-finite verb (in infinitive or supine)||Object(s)||Spatial Adverb||Temporal Adverb| |Conjunction||Subject||Clausal Adverb/Negation||Finite Verb||Non-finite verb (in infinitive or supine)||Object(s)||Spatial Adverb||Temporal Adverb| The "Fundament" can be whatever constituent that the speaker wishes to topicalize, emphasize as the topic of the sentence. In the unmarked case, with no special topic, the subject is placed in the fundament position. Common fundaments are an adverb or object, but it is also possible to topicalize basically any constituent, including constituents lifted from a subordinate clause into the fundament position of the main clause: Honom vill jag inte att du träffar. (Him I do not want you to meet.) or even the whole subordinate clause: Att du följer honom hem accepterar jag inte. (That you follow him home I do not accept.). An odd case is the topicalization of the finite verb, which requires the addition of a "dummy" finite verb in the V2 position, so that the same clause has two finite verbs: Arbetade gjorde jag inte igår. (Worked did I not yesterday.) - Källström, Roger. "Omarkerat neutrum?" (PDF). Göteborgs universitet. Retrieved 2008-03-26.[dead link] - Pettersson, 150-51. - Språkrådet. "Heter det Konungens av Danmark bröstkarameller eller Konungen av Danmarks bröstkarameller?" (in Swedish). Retrieved 21 July 2014. - Språktidningen, "Så snabbt ökar hen i svenska medier", 18 March 2013. Retrieved 27 July 2014. - "The Local", "Gender-neutral 'hen' makes its legal debut", 14 December 2012. Retrieved 27 July 2014. - Terese Allert, "Allt vanligare med hen i barnböcker", Aftonbladet, 15 March 2013. Retrieved 27 July 2014. - http://www.kristianstadsbladet.se/debatt/hall-hen-borta-fran-vara-barn/; in April 2015 it was added to Svenska Akademiens ordlista, the official glossary of the Swedish Academy - Holmes, Philip & Hinchliffe, Ian (2008) Swedish: An Essential Grammar Routledge: New York ISBN 0-415-45800-5 - Holmes, Philip & Hinchliffe, Ian (2003) Swedish: A Comprehensive Grammar Routledge: New York ISBN 0-415-27884-8 - Pettersson, Gertrud (1996) Svenska språket under sjuhundra år: en historia om svenskan och dess utforskande Lund: Studentlitteratur ISBN 91-44-48221-3

https://en.wikipedia.org/wiki/Swedish_grammar

As its name suggests, the elongate body of the small-scaled skink is covered in relatively small, glossy scales (2)(3). The background colour to the upperparts of the body is brownish grey, but a series of stripes extend lengthways from the snout towards the tail. Running down the middle of the back are consecutive segments of light and dark brown, adjoined on either side by a conspicuous pale stripe. A dark brown stripe, speckled above and below with pale markings, extends along the sides, while the belly is pale all over (2). Very little is known about the biology of the small-scaled skink other than it is an active diurnal forager (2)(4). In captivity, it will consume a wide variety of invertebrates (2), but most New Zealand skinks are omnivorous with fruit and insects known to form a large proportion of their diet (3). In captivity, the young are born from late January to early March with two to three offspring in each litter (2). As with other New Zealand skinks, habitat loss and introduced mammalian predators are thought to present the greatest threat to the small-scaled skink (3)(4). Owing to these impacts, the small-scaled skink population is believed to be undergoing a serious decline (4). With so many unknowns associated with the small-scaled skink, the immediate priority is to conduct further research into the species’ conservation status by obtaining data on its distribution, habitat use, relative abundance and threats, including the impact of mammalian predators. The collated information will then be used to determine the optimum means of ensuring the survival of this species (4). Embed this ARKive thumbnail link ("portlet") by copying and pasting the code below.

http://www.arkive.org/small-scaled-skink/oligosoma-microlepis/

You are hereHome › Now showing results 1-10 of 23 In this activity, children use common craft materials and ultraviolet (UV)-sensitive beads to construct a person (or dog or imaginary creature). They use sunscreen, foil, paper, and more to test materials that might protect UV Kid from being exposed... (View More) to too much UV radiation. Includes background for facilitators. This activity is part of the "Explore!" series of activities designed to engage children in space and planetary science in libraries and informal learning environments. (View Less) Each lesson or activity in this toolkit is related to NASA's Lunar Reconnaissance Orbiter (LRO). The toolkit is designed so that each lesson can be done independently, or combined and taught in a sequence. The Teacher Implementation Guide provides... (View More) recommendations for combining the lessons into three main strands: 1) Lunar Exploration. These lessons provide a basic introduction to Moon exploration. Note that this strand is also appropriate for use in social studies classes. 2) Mapping the Moon. These lessons provide a more in-depth understanding of Moon exploration through the use of scientific data and student inquiry. The lessons also include many connections to Earth science and geology. 3) Tools of Investigation. These higher-level lessons examine the role of technology, engineering and physics in collecting and analyzing data. (View Less) This project engages students in the science and engineering processes used by NASA Astrobiologists as they explore our Solar System and try to answer the compelling question, "Are we Alone?" Students will identify science mission goals and select... (View More) an astrobiologically significant target of interest: Mars, Europa, Enceladus or Titan. Students will then design their mission to this target in search of their chosen biosignature(s). Students will encounter the same considerations and challenges facing NASA scientists and engineers as they search for life in our Solar System. Students will need to balance the return of their science data with engineering limitations such as power, mass and budget. Risk factors play a role and will add to the excitement in this interactive science and engineering activity. Astrobiobound! will help students see how science and systems engineering are integrated to achieve a focused scientific goal. Includes an alignment document for NGSS and Common Core State Standards. (View Less) This activity focuses on the relationship between science of looking for life and the tools, on vehicles such as the Mars Rover, that make it possible. Learners will create their own models of a Mars rover. They determine what tools would be... (View More) necessary to help them better understand Mars (and something about life on Mars/its habitability). Then they work in teams to complete a design challenge where they incorporate these elements into their models, which must successfully complete a task. Teams may also work together to create a large-scale, lobby-sized version that may be put on display in the library to engage their community. The activity also includes specific tips for effectively engaging girls in STEM. This is activity 6 in Explore: Life on Mars? that was developed specifically for use in libraries. (View Less) This is an annotated, topical list of science fiction novels and stories based on more or less accurate astronomy and physics ideas. Learners can read fictional works that involve asteroids, astronomers, black holes, comets, space travel where... (View More) Einstein's ideas are used correctly, exploding stars, etc. (View Less) This is a set of four activities about spacecraft design. Learners will use the information learned in previous lessons, combined with their own creativity and problem-solving skills, to design and test a parachuting probe that will withstand a fall... (View More) from a high point, land intact, be able to descend slowly, float in liquid, and cost the least to launch into space. Includes a glossary, information for families, and guidance for deepening the science. This is lesson 7 of 8 in the Jewel of the Solar System: From Out-of-School to Outer Space an adaptation for afterschool programs of the Cassini-Huygens educational product Reading, Writing, and Rings. (View Less) This is a series of three webpages about how humans and computers communicate. Learners will explore the binary and hexidecimal systems and how engineers use them to translate spacecraft data into images. This is a game about data compression. Learners will use virtual foam balls to explore the different compression methods (lossless, lossy, and superchannel) used by the Earth Observing 3 mission. This is a set of four activities about spacecraft design. Learners will think like engineers as they design, peer review, and then construct and present their spacecraft to travel to Saturn. Includes a glossary, information for families, and... (View More) guidance for deepening the science. This is lesson 5 of 8 in the Jewel of the Solar System: From Out-of-School to Outer Space an adaptation for afterschool programs of the Cassini-Huygens educational product Reading, Writing, and Rings. (View Less)

http://www.nasawavelength.org/resource-search?qq=&facetSort=1&educationalLevel=Informal+education&topicsSubjects=The+nature+of+technology&smdForumPrimary=Planetary+Science

NORD gratefully acknowledges Carole Samango-Sprouse, EdD, Executive Director and Chief Science Officer, The Focus Foundation, for assistance in the preparation of this report. Trisomy X is a disorder that affects females and is characterized by the presence of an additional X chromosome. Normally, females have two X chromosomes; however, females with trisomy X carry three X chromosomes in the nuclei of body cells. There are specific physical features (phenotype) associated with this chromosomal disorder. Common symptoms that can potentially occur include language-based learning disabilities, developmental dyspraxia, tall stature, low muscle tone (hypotonia), and abnormal bending or curving of the pinkies toward the ring fingers (clinodactyly). Trisomy X occurs randomly as a result from errors during the division of reproductive cells in one of the parents. This disorder occurs in one in 900 to 1,000 live births. The symptoms and physical features associated with trisomy X vary greatly from one person to another. Some females may have no symptoms (asymptomatic) or very mild symptoms and may go undiagnosed. Other women may have a wide variety of different abnormalities. It is important to note that affected individuals may not have all of the symptoms discussed below. Affected individuals should talk to their specialists and medical team about their specific case, associated symptoms and overall prognosis. Trisomy X is often associated with developmental differences and language-based learning disabilities. Intelligence is usually within the normal range. IQ may be 10-15 points below that of siblings or control groups if early intervention has not been successful or begun early enough. Infants and children with trisomy X experience delays in attaining developmental milestones, especially in the acquisition of motor and speech skills. For example, walking may be delayed and affected girls may exhibit poor coordination and clumsiness. Speech and language development is also commonly delayed and may become apparent by approximately one year to 18 months. Girls with trisomy X have an increased frequency of language-based learning disabilities including reading deficiencies such as dyslexia, reading comprehension deficits and/or reading fluency issues in conjunction with other language-based disabilities. They also have developmental dyspraxia which affects learning in every domain. Typically, motor planning skills are deficient, which affects gross and fine motor, speech and language as well as executive function. During early childhood or adolescence, girls with trisomy X usually exhibit increased height as compared with other girls their age (tall stature). Most girls are at or above the 75th percentile for height with an average height of 5 foot 7 inches. In some cases, infants with trisomy X may have mild facial abnormalities including vertical skin folds that may cover the eyes’ inner corners (epicanthal folds), widely spaced eyes (hypertelorism), and smaller than normal head circumference. Most infants also have decreased muscle tone (hypotonia) and the fifth finger may be abnormally bent or curved mildly, which is called clinodactyly. Individuals with trisomy X may have an increased incidence of anxiety and attention deficit hyperactivity disorder (ADHD). In some cases, such abnormalities improve with maturity and as the girls reach adulthood. Some individuals have minimal to no behavioral or emotional abnormalities; others have more issues that may necessitate intervention, which is typically only necessary short term. There are no controlled studies on behavioral or emotional abnormalities in trisomy X and the incidence of such conditions is unknown, although they are believed to occur with greater frequency than in the general population. Early detection and treatment are very beneficial for girls with trisomy X. In many cases, these girls have few issues later in life when identified early and treated appropriately. In most cases, sexual development and fertility are normal. However, reports indicate that some affected females may have abnormal development of the ovaries (ovarian dysgenesis) and/or the uterus; delayed puberty or early onset of puberty (precocious puberty), and/or fertility problems. There have been reports of women with trisomy X developing premature ovarian failure (POF). POF is the loss of function of the ovaries before the age where menopause is expected to begin. POF can cause a decrease in the production of certain hormones and eggs may no longer be released each month. Less often, additional abnormalities have been described in individuals with trisomy X including kidney abnormalities, such as absence of a kidney (unilateral renal agenesis) or malformation (dysplasia) of the kidneys; recurrent urinary tract infections; seizures; constipation; abdominal pain; flatfeet (pes planus); and pectus excavatum, a condition in which the breastbone is mildly depressed into the chest. Heart (cardiac) abnormalities have also been described in some isolated cases. Trisomy X is a chromosomal abnormality characterized by the presence of an extra X chromosome. Chromosomes are found in the nucleus of all body cells. They carry the genetic characteristics of each individual. Pairs of human chromosomes are numbered from 1 through 22, with an unequal 23rd pair that normally consists of an X and Y chromosome for males and two X chromosomes for females. Thus, females with a normal chromosomal make-up (karyotype) have 46 chromosomes, including two X chromosomes (46,XX karyotype); they receive one chromosome from the mother and one from the father in each of the 23 pairs. However, females with trisomy X have 47 chromosomes, three of which are X chromosomes (47,XXX karyotype). Trisomy X is a genetic disorder, but it is not inherited. The presence of the extra X chromosome results from errors during the normal division of reproductive cells in one of the parents (nondisjunction during meiosis). These errors occur randomly for no apparent reason (sporadically). Studies have shown that the risk of such errors increases with advanced paternal age. In most cases, the additional X chromosome comes from the mother. In approximately 20 percent of cases, nondisjunction events occur after conception in the developing fetus (postzygotic nondisjunction). In some cases, only a certain percentage of an individual’s cells may have three X chromosomes, while others have a normal chromosomal make-up (46,XX/47,XXX mosaicism). Evidence suggests that such cases may be associated with milder symptoms and fewer developmental and learning problems, but further research is needed. Variants have also been described in which cells contain four or five X chromosomes (tetra X syndrome and penta X syndrome). Such variants are typically associated with more severe symptoms and findings. (For further information, please see the “Related Disorders” section of this report below.) Researchers believe that the symptoms and physical features associated with trisomy X develop because of overexpression of the genes that escape normal X-inactivation. Although females have two X chromosomes, one of the X chromosomes is “turned off” and all of the genes on that chromosome are inactivated (X-inactivation). Researchers suspect that the presence of a third X chromosome allows genes normally “turned off” to be expressed. However, the exact manner in which the extra X chromosome ultimately causes the symptoms and physical features of trisomy X is not fully understood. Trisomy X is a chromosomal disorder that affects only females. Reported estimates concerning the disorder’s frequency have varied with the most common estimate being one in 1,000 female births. Because many females with the disorder may have few or no symptoms, they may never be diagnosed. Researchers believe that the disorder is underdiagnosed and that the reported number of cases as reflected in the medical literature is inappropriately low. Researchers believe that only approximately 10 percent of cases are diagnosed. With increased detection, more in depth studies may be conducted and more girls with triple X can be appropriately treated. Trisomy X may be suspected based upon the identification of characteristic developmental, behavioral or learning disabilities. A diagnosis may be confirmed by a thorough clinical evaluation, a detailed family history, and certain specialized tests such as chromosomal analysis performed on blood samples that can reveal the presence of an extra X chromosome in body cells. In addition, trisomy X is increasingly being diagnosed before birth (prenatally) based on chromosomal analysis performed subsequent to amniocentesis or chorionic villus sampling (CVS). During amniocentesis, a sample of fluid that surrounds the developing fetus is removed and analyzed, while CVS involves the removal of tissue samples from a portion of the placenta. Approximately 5-15 percent of women with Turner syndrome also have a 47,XXX karyotype found in certain white blood cells (blood lymphocytes), but the characteristic Turner syndrome karyotype (45,X) in other cells. Specific therapeutic strategies depend upon several factors including the age of an affected individual upon diagnosis, the specific symptoms that are present and the overall severity of the disorder in each case. Early intervention services are recommended for infants and children diagnosed with trisomy X. Experts advise developmental assessment by age four months to evaluate muscle tone and strength; language and speech assessment by 12 months of age to evaluate expressive and receptive language development; and pre-reading assessment during preschool years prior to first grade to look for early signs of reading dysfunction. An evaluation is recommended to help assess additional learning disabilities and social and emotional problems. Evidence suggests that affected children are greatly responsive to early intervention services and treatment. Such services can include speech therapy, occupational therapy, physical therapy, and developmental therapy and counseling. Infants and children with trisomy X should also receive kidney (renal) and heart (cardiac) evaluations to detect abnormalities of those organs potentially associated with the disorder. Adolescent and adult women who exhibit late periods (menarche), menstrual abnormalities, or fertility issues should be evaluated for primary ovarian failure. Genetic counseling will be of benefit for affected individuals and their families. Additional treatment for this disorder should be targeted at infancy for physical therapy, between 12 and 15 months for speech delay, prior to first grade for early signs of reading dysfunction and by third grade for anxiety and ADHD. Adolescence is challenging for children, and girls with triple X often struggle as they enter middle school years so counseling short term may be necessary to help them during these turbulent years. Information on current clinical trials is posted on the Internet at www.clinicaltrials.gov. All studies receiving U.S. Government funding, and some supported by private industry, are posted on this government web site. For information about clinical trials being conducted at the NIH Clinical Center in Bethesda, MD, contact the NIH Patient Recruitment Office: Tollfree: (800) 411-1222 TTY: (866) 411-1010 For information about clinical trials sponsored by private sources, contact: For information about clinical trials conducted in Europe, contact: (Please note that some of these organizations may provide information concerning certain conditions potentially associated with this disorder [e.g., learning disabilities].) Speicher MR, Antonarakis SE, Motulsky AG. Eds. Vogel and Motulsky’s Human Genetics: Problems and Approaches. 4th ed. Springer. New York, NY; 2009:124. Samango-Sprouse CA. XXX Syndrome (Triple X Syndrome). NORD Guide to Rare Disorders. Lippincott Williams & Wilkins. Philadelphia, PA. 2003:89. Samango-Sprouse CA Frontal Lobe Development in Childhood. The Human Frontal Lobe: Functions and Disorders, 2nd Edition, Eds. BL Miller, and JL Cummings, Guilford Press, New York, 2007. Rimoin D, Connor JM, Pyeritz RP, Korf BR. Eds. Emory and Rimoin’s Principles and Practice of Medical Genetics. 4th ed. Churchill Livingstone. New York, NY; 2002:1195-1196. Otter M, Schrander-Stumpel CT, Curfs LM. Triple X syndrome: a review of the literature. Eur J Hum Genet. 2010;18:265-271. Krusinskie V, Alvesalo L, Sidlauskas A. The craniofacial complex in 47,XXX females. Eur J Orthod. 2005;27:396-401. Liebezeit BU, Rohrer TR, Singer H, Doerr HG. Tall stature as presenting symptom in a girl with triple X syndrome. J Pediatr Endocrinol Metab. 2003;16:233-235. Rovet J, Netley C, Bailey J, Keenan M, Stewart D. Intelligence and achievement in children extra X aneuploidy: a longitudinal perspective. Am J Med Genet. 1995;60:356-363. Raticliffe SG, Pan H, McKie M. The growth of XXX females: population-based studies. Ann Hum Biol. 1994;21:57-66. Samango-Sprouse CA, Rogol A. XXY: The Hidden Disability and Prototype for Infantile Presentation of Developmental Dyspraxia (IDD). Infants and Young Children. 2002;15:11-18. Tartaglia NR, Howell S, Sutherland A, Wilson R, Wilson L. A review of trisomy X (47,XXX). Orphanet encyclopedia, 2010. Available at: http://www.ojrd.com/content/5/1/8 Accessed March 25, 2014. Mayo Clinic for Medical Education and Research. Triple X Syndrome. Nov. 08, 2012. Available at: http://www.mayoclinic.com/health/triple-x-syndrome/DS01090 Accessed March 25, 2014. The information in NORD’s Rare Disease Database is for educational purposes only and is not intended to replace the advice of a physician or other qualified medical professional. The content of the website and databases of the National Organization for Rare Disorders (NORD) is copyrighted and may not be reproduced, copied, downloaded or disseminated, in any way, for any commercial or public purpose, without prior written authorization and approval from NORD. Individuals may print one hard copy of an individual disease for personal use, provided that content is unmodified and includes NORD’s copyright. National Organization for Rare Disorders (NORD) 55 Kenosia Ave., Danbury CT 06810 • (203)744-0100

http://rarediseases.org/rare-diseases/trisomy-x/

Introduction to Named Pipes Bash uses named pipes in a really neat way. Recall that when you enclose a command in parenthesis, the command is actually run in a “subshell”; that is, the shell clones itself and the clone interprets the command(s) within the parenthesis. Since the outer shell is running only a single “command”, the output of a complete set of commands can be redirected as a unit. For example, the command: (ls -l; ls -l) >ls.out writes two copies of the current directory listing to the file ls.out. Command substitution occurs when you put a < or > in front of the left parenthesis. For instance, typing the command: cat <(ls -l) results in the command ls -l executing in a subshell as usual, but redirects the output to a temporary named pipe, which bash creates, names and later deletes. Therefore, cat has a valid file name to read from, and we see the output of ls -l, taking one more step than usual to do so. Similarly, giving >(commands) results in Bash naming a temporary pipe, which the commands inside the parenthesis read for input. If you want to see whether two directories contain the same file names, run the single command: cmp <(ls /dir1) <(ls /dir2) The compare program cmp will see the names of two files which it will read and compare. Command substitution also makes the tee command (used to view and save the output of a command) much more useful in that you can cause a single stream of input to be read by multiple readers without resorting to temporary files—bash does all the work for you. The command: ls | tee >(grep foo | wc >foo.count) \ >(grep bar | wc >bar.count) \ | grep baz | wc >baz.count counts the number of occurrences of foo, bar and baz in the output of ls and writes this information to three separate files. Command substitutions can even be nested: cat <(cat <(cat <(ls -l))))works as a very roundabout way to list the current directory. As you can see, while the unnamed pipes allow simple commands to be strung together, named pipes, with a little help from bash, allow whole trees of pipes to be created. The possibilities are limited only by your imagination. Practical books for the most technical people on the planet. Newly available books include: - Agile Product Development by Ted Schmidt - Improve Business Processes with an Enterprise Job Scheduler by Mike Diehl - Finding Your Way: Mapping Your Network to Improve Manageability by Bill Childers - DIY Commerce Site by Reven Lerner Plus many more.

http://www.linuxjournal.com/article/2156?page=0,1

|Classification and external resources| Intestinal parasites are parasites that can infect the gastro-intestinal tract of humans and other animals. They can live throughout the body, but most prefer the intestinal wall. Means of exposure include: ingestion of undercooked meat, drinking infected water, and skin absorption. The two main types of intestinal parasites are those helminths and protozoa that reside in the intestines (not all helminths and protozoa are intestinal parasites). An intestinal parasite can damage or sicken its host via an infection which is called helminthiasis in the case of helminths. Signs and symptoms These depend on the type of infection. The major groups of parasites include protozoans (organisms having only one cell) and parasitic worms (helminths). Of these, protozoans, including cryptosporidium, microsporidia, and isospora, are most common in HIV-infected persons. Each of these parasites can infect the digestive tract, and sometimes two or more can cause infection at the same time. Parasites can get into the intestine by going through the mouth from uncooked or unwashed food, contaminated water or hands, or by skin contact with larva infected soil; they can also be transferred by the sexual act of anilingus in some cases. When the organisms are swallowed, they move into the intestine, where they can reproduce and cause symptoms. Children are particularly susceptible if they are not thoroughly cleaned after coming into contact with infected soil that is present in environments that they may frequently visit such as sandboxes and school playgrounds. People in developing countries are also at particular risk due to drinking water from sources that may be contaminated with parasites that colonize the gastrointestinal tract. Due to the wide variety of intestinal parasites, a description of the symptoms rarely is sufficient for diagnosis. Instead, two common tests are used: stool samples may be collected to search for the parasites, and an adhesive may be applied to the anus in order to search for eggs. Good hygiene is necessary to avoid reinfection. The Rockefeller Foundation's hookworm campaign in Mexico in the 1920s was extremely effective at eliminating hookworm from humans with the use of antihelminthics. However, preventative measures were not adequately introduced to the people that were treated. Therefore, the rate of reinfection was extremely high and the project evaluated through any sort of scientific method was a marked failure. More education was needed to inform the people of the importance of wearing shoes, using latrines (better access to sanitation), and good hygiene. Drugs are frequently used to kill parasites in the host. In earlier times, turpentine was often used for this, but modern drugs do not poison intestinal worms directly. Rather, antihelmintic drugs now inhibit an enzyme that is necessary for the worm to make the substance that prevents the worm from being digested. For example, tapeworms are usually treated with a medicine taken by mouth. The most commonly used medicine for tapeworms is Praziquantel. Praziquantel is also used to treat infections of certain parasites (e.g., Schistosoma and liver flukes). - Loukopoulos P, Komnenou A, Papadopoulos E , Psychas V. Lethal Ozolaimus megatyphlon infection in a green iguana (Iguana iguana rhinolopa). Journal of Zoo and Wildlife Medicine 2007; 38:131-134 - Birn, Anne-Emanuelle, and Armando Solórzano. 1999. Public health policy paradoxes: science and politics in the Rockefeller Foundation's hookworm campaign in Mexico in the 1920s. Social Science & Medicine 49 (9):1197-1213

https://en.wikipedia.org/wiki/Intestinal_worms

Print this page. Home / Browse / Mississippi Alluvial Plain The Mississippi Alluvial Plain (a.k.a. Delta) is a distinctive natural region, in part because of its flat surface configuration and the dominance of physical features created by the flow of large streams. This unique physiography occupies much of eastern Arkansas including all or parts of twenty-seven counties. The Alluvial Plain, flatter than any other region in the state, has elevations ranging from 100 to 300 feet above sea level. In Arkansas, the Alluvial Plain extends some 250 miles in length from north to south and varies in width from east to west from only twelve miles in Desha County to as much as ninety-one miles measured from Little Rock (Pulaski County) to the Mississippi River. The work of large rivers (including the Mississippi, Arkansas, White, and St. Francis rivers) and other smaller rivers and streams has played an important role in forming the character of the landscape. These rivers eroded older deposits and built up deep layers of soil, gravel, and clay transported from slopes as far away as the Rocky Mountains to the west and the Appalachians to the east. The result of these alluvial processes is a terrain and soil suitable for large-scale farming. In fact, the Mississippi Alluvial Plain is one of the most agriculturally productive regions in the world. Alluvial (stream-deposited) material covers almost the entire region. Interestingly, terraces are found throughout the Alluvial Plain, frequently paralleling streams but at a slightly higher elevation than the adjacent stream banks. These terraces are older than present bottomlands and represent former levels of bottomland through which streams have now eroded. The so-called recent alluvium has been deposited over the last 12,000 years and contains fertile “water-washed” material, especially silt. The deep, fertile soils of the Mississippi Alluvial Plain are sometimes extremely dense and poorly drained. The combination of flat terrain and poor drainage creates conditions suitable for wetlands. Wetlands, areas where the periodic or permanent presence of water controls the characteristics of the environment and associated plants and animals, now cover approximately eight percent of Arkansas’s land surface. While some wetland areas remain intact, many have been drained and converted to agricultural land uses. Protecting the remaining wetlands and encouraging the restoration of some former wetland areas are significant natural resource conservation issues. At one time, wetlands were very abundant across the Mississippi Alluvial Plain. The decline in wetlands began years ago when the first ditches were dug to drain extensive areas of the Alluvial Plain. Clearing bottomland hardwoods for agriculture and other activities has resulted in the loss of more than seventy percent of the original wetlands. The majority of Arkansas’s wetlands, occupying a diverse physiographic setting, are often riverine and depressional wetlands associated with the floodplains of the Mississippi River and its major tributaries. Some of the most significant wetlands are referred to as “bottoms” or “bottomland hardwood forests.” Of particular importance is the Cache River and lower White River area, where impressive stands of bottomland hardwoods are found. It represents the largest continuous expanse of bottomland hardwoods in the Lower Mississippi Valley. Nearly one-third of the remaining bottomland hardwoods in the Arkansas Delta are found within the ten-year floodplain of the Cache and lower White rivers. The wetlands of the Delta offer an internationally important winter habitat for migratory water fowl. The White River National Wildlife Refuge alone is a temporary home to between 3,000 and 10,000 Canada geese and up to 300,000 ducks per year. These large numbers account for one-third of the total found in Arkansas and ten percent of the Mississippi Flyway total. Wetlands of Arkansas serve many important functions in addition to being a vital wildlife habitat, including flood storage and flood prevention, natural water quality improvement (sediment traps, for example), shoreline erosion protection, groundwater recharge, recreational opportunities, and aesthetic beauty. The original natural vegetation of the region was significantly different from the other natural regions in Arkansas in part because of the region’s wetland characteristics. It was largely southern floodplain forest suited to the wet, poorly drained soils. Cypress-tupelo-gum types occupied the wettest sites. The willow oak and overcup oak were found on flat and poorly drained locations, and oak-hickory on higher and better drained terrace sites of the floodplain. Currently, the Mississippi Alluvial Plain has been widely cleared and drained for cultivation. The widespread loss or degradation of forest and wetland habitat has impacted wildlife and reduced bird populations. Relatively small plots of natural vegetation remain along streams, in areas unsuitable for agriculture, or within areas protected from clearing and development. The most significant of these protected areas are the Big Lake National Wildlife Refuge, the Sunken Lands Wildlife Management Areas, the Wapanocca National Wildlife Refuge, the St. Francis National Forest, and the White River National Wildlife Refuge. A rather unique feature of this region is the Grand Prairie, an area of prairie soils and grasses that are found primarily in Arkansas and Prairie counties in eastern Arkansas. These prairie soils, with their very compact clay subsoil, are more suitable for grass than trees as the natural vegetation cover. The Grand Prairie is an extremely productive agricultural region and is noted for its high yields of rice. Stuttgart (Arkansas County) is known as the rice capital and duck capital of the world. Another important characteristic of the Alluvial Plain is related to a significant natural hazard, earthquakes. These may occur along the New Madrid Seismic Zone. This seismic zone is a prolific source of intra-plate earthquakes (earthquakes within a tectonic plate) in the southern and mid-western United States. This seismic zone was responsible for the 1811–1812 New Madrid Earthquakes and has the potential to produce large earthquakes in the future. Several relatively small earthquakes have been recorded in the region since 1812, but an important question remains concerning the next “big” earthquake in terms of when it will occur and at what magnitude. As of 2011, according to some experts, there is a ten percent chance of a magnitude 7.0 quake within the next fifty years along the fault that extends from New Madrid, Missouri, to Marked Tree (Poinsett County) and beyond. In addition to a unique physical landscape, the Mississippi Alluvial Plain has a number of distinctive cultural/demographic characteristics. In an article titled “Delta Population Trends: 1990–2000,” Jason Combs discusses the significant population decline that has occurred within counties bordering the Mississippi River. Population decline, economic depression, and other negative socioeconomic factors characterize many of these Delta counties. Data released from the 2010 census shows that the population decline is continuing within several Arkansas counties that are adjacent to or near the Mississippi River. The most significant population declines (between -10.1 and -20.5 percent) from 2000 to 2010 were in Mississippi, Lee, Phillips, Desha, Chicot, Monroe, and Woodruff counties. These counties have relatively high rates of unemployment and few or no positive features to reverse the trend of population decline, according to Combs. These and other Delta counties have experienced population decline for a variety of reasons, in addition to high unemployment. According to Combs, part of the problem is the image of the Delta. Strained race relations, poverty, and resistance to social change have “tarnished” the Delta’s image and contributed to the absence of substantial economic development. Moreover, most Americans perceive the Delta as “flat and uninteresting, not a place to go for recreation, retirement, or a glamorous job,” according to an article by Richard Lonsdale and J. Clark Archer. Another contributing factor to the population loss in the Delta is agricultural mechanization. Improvement in mechanization and modern science allowed fewer farmers to produce as much or more agricultural output on the same amount of land with far less labor. The need for fewer farm workers coupled with the absence of other job opportunities has been a significant contributing factor to the population decline and the economic depression that many Delta counties are experiencing. In summary, the Mississippi Alluvial Plain is a natural region with several distinguishing characteristics, including an extremely flat surface topography; deep alluvial soils; poor drainage; wetland areas; widely scattered bottomland and hardwood forests; large and highly productive farms; counties plagued by economic depression and population loss; and the Mississippi Flyway, with ideal locations for hunting, fishing, and other water-related sports activities. The result is a region marked by sharp social contrast: pockets of prosperity and wealth exist aside poverty and economic despair. For additional information: Arkansas Department of Planning. Arkansas Natural Area Planning. Little Rock: State of Arkansas, 1974. Collins, Janelle, ed. Defining the Delta: Multidisciplinary Perspectives on the Lower Mississippi River Delta. Fayetteville: University of Arkansas Press, 2015. Combs, Jason. “Delta Population Trends: 1990–2000.” Arkansas Review: A Journal of Delta Studies 34 (April 2004): 26–35. “Delta Geography.” Delta Cultural Center. http://www.deltaculturalcenter.com/geography/ (accessed July 18, 2011). “Ecoregions of the Mississippi Alluvial Plain.” The Encyclopedia of Earth. http://www.eoearth.org/article/Ecoregions_of_the_Mississippi_Alluvial_Plain_%28EPA%29 (accessed July 18, 2011). Hagge, Patrick David. “The Decline and Fall of a Cotton Empire: Economic and Land-Use Change in the Lower Mississippi River ‘Delta’ South, 1930–1970.” PhD diss., Pennsylvania State University, 2013. Lonsdale, Richard, and J. Clark Archer. “Emptying Areas of the United States, 1990–1995.” Journal of Geography 97 (1998): 108–122. Stroud, Hubert B., and Gerald T. Hanson. Arkansas Geography: The Physical Landscape and the Historical-Cultural Setting. Little Rock: Rose Publishing Company, 1981. “Wetlands in Arkansas.” Arkansas Multi-Agency Wetland Planning Team. http://www.mawpt.org/wetlands/ (accessed July 18, 2011). Whayne, Jeannie, and Willard B. Gatewood, eds. The Arkansas Delta: Land of Paradox. Fayetteville: University of Arkansas Press, 1993. Hubert B. Stroud Arkansas State University Last Updated 12/11/2015 About this Entry: Contact the Encyclopedia / Submit a Comment / Submit a Narrative

http://www.encyclopediaofarkansas.net/encyclopedia/entry-detail.aspx?entryID=444

palette(1) In computer graphics, a palette is the set of available colors. For a given application, the palette may be only a subset of all the colors that can be physically displayed. For example, a SVGA system can display 16 million unique colors, but a given program would use only 256 of them at a time if the display is in 256-color mode. The computer system's palette, therefore, would consist of the 16 million colors, but the program's palette would contain only the 256-color subset. A palette is also called a CLUT (color look-up table). (2) In paint and illustration programs, a palette is a collection of symbols that represent drawing tools. For example, a simple palette might contain a paintbrush, a pencil, and an eraser. Stay up to date on the latest developments in Internet terminology with a free weekly newsletter from Webopedia. Join to subscribe now. Webopedia's student apps roundup will help you to better organize your class schedule and stay on top of assignments and homework. Read More »20 Ways to Shorten a URL If you need to shorten a long URL try this list of 20 free online redirection services. Read More »Top 10 Tech Terms of 2015 The most popular Webopedia definitions of 2015. Read More » This Webopedia study guide describes the different parts of a computer system and their relations. Read More »Network Fundamentals Study Guide Networking fundamentals teaches the building blocks of modern network design. Learn different types of networks, concepts, architecture and... Read More »The Five Generations of Computers Learn about each of the five generations of computers and major technology developments that have led to the current devices that we use today. Read More »

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So You Want to Be President? Lesson Plan Activities to do before and after reading the book by Judith St. George - Grades: 3–5 About this book Humorous and just slightly off-kilter, this book is sure to entertain while it enlightens. - Learn interesting and often little-known facts about the political leaders who have governed our nation from its beginning. - Understand how democratic values came to be, and how they have been exemplified by people, events, and symbols. Before Reading the Book The Name Game Got any history buffs in your class? Test everyone's knowledge with a quick and easy name game. - On your blackboard, write the numbers 1–42. - Ask students to pick their brains and see how many presidents they can name. - Using the list at the back of So You Want to be President? (if necessary), write each president your class can name in his proper spot. - Try to spot any trends or patterns in the list — lots of men named James, some relatives, etc. - You may want to fill in the list, as a class, after you've read the book. Teaching Plan: Activities Even United States Presidents started out as regular kids! You can help students understand presidential legacies by imagining their own. - Ask students to think about any interesting or important facts about their own lives. - On a piece of paper, ask each to create a time line of his or her life to this point. (For example, born on this day, little sister arrived on this day, started piano lessons on this day, etc.). - Next, have your class think about what they'd like their future to hold. Ask them to record these anticipated future events in a different color pencil or ink. - When each has finished, have students share their "legacies," both current and anticipated, with the class. - Post timelines on a classroom bulletin board. My Favorite President So You Want to Be President? is chock-full of interesting tidbits about our Presidents. Use them to captivate your students and encourage them to learn more! - Using the information gleaned from the book, ask each student to choose a favorite president. The choice should leave aside political associations; it should be based solely on the trivia presented. - Have students create a list of facts about their chosen president. - Ask each student to give an oral presentation about his or her favorite. The very brief report could begin with a list of facts about the president and conclude with reasons why the student found these qualities interesting or appealing. - As a class, discuss each student's choice; were some presidents chosen more often than others? Analyze the outcome. Other Books About the Electoral Process Presidential Elections and Other Cool Facts by Syl Sobel This handy resource book outlines the legal requirements for electing a president, discusses the Electoral College, campaign rules and practices, and much more! Landslide! A Kid's Guide to the U.S. Elections by Dan Gutman How does a voting machine work? Who belongs to the Electoral College? What happens if there's a tie? Find answers to these questions, and many more. Other Books by Judith St. George Teaching Plan written by Rebecca Gómez.

http://www.scholastic.com/teachers/lesson-plan/so-you-want-be-president-lesson-plan

|This article needs additional citations for verification. (September 2009)| - details about the world itself and the experiences of its characters are revealed explicitly through narrative - the story is told or recounted, as opposed to shown or enacted. In diegesis the narrator tells the story. The narrator presents the actions (and sometimes thoughts) of the characters to the readers or audience. In contrast to mimesis Diegesis (Greek διήγησις "narration") and mimesis (Greek μίμησις "imitation") have been contrasted since Plato's and Aristotle's times. Mimesis shows rather than tells, by means of action that is enacted. Diegesis is the telling of the story by a narrator. The narrator may speak as a particular character or may be the invisible narrator or even the all-knowing narrator who speaks from "outside" in the form of commenting on the action or the characters. In Book III of his Republic (c. 373 BC), Plato examines the "style" of "poetry" (the term includes comedy, tragedy, epic and lyric poetry): All types narrate events, he argues, but by differing means. He distinguishes between narration or report (diegesis) and imitation or representation (mimesis). Tragedy and comedy, he goes on to explain, are wholly imitative types; the dithyramb is wholly narrative; and their combination is found in epic poetry. When reporting or narrating, "the poet is speaking in his own person; he never leads us to suppose that he is any one else"; when imitating, the poet produces an "assimilation of himself to another, either by the use of voice or gesture". In dramatic texts, the poet never speaks directly; in narrative texts, the poet speaks as him or herself. In his Poetics, the ancient Greek philosopher Aristotle argues that kinds of "poetry" (the term includes drama, flute music, and lyre music for Aristotle) may be differentiated in three ways: according to their medium, according to their objects, and according to their mode or "manner" (section I); "For the medium being the same, and the objects the same, the poet may imitate by narration — in which case he can either take another personality as Homer does, or speak in his own person, unchanged — or he may present all his characters as living and moving before us" (section III). In filmmaking, the term is used to name the story depicted on screen—as opposed to the story in real life time that the screen narrative is about. Diegesis may concern elements, such as characters, events and things within the main or primary narrative. However, the author may include elements which are not intended for the primary narrative, such as stories within stories; characters and events that may be referred to elsewhere or in historical contexts and that are therefore outside the main story and are thus presented in an extradiegetic situation. For narratologists, all parts of narratives — characters, narrators, existents, actors — are characterized in terms of diegesis. For definitions of diegesis, one should consult Aristotle's Poetics; Gerard Genette's Narrative Discourse: An Essay in Method (Cornell University Press, 1980); or (for a readable introduction) H. Porter Abbott's The Cambridge Introduction to Narrative (Cambridge University Press 2002). In literature, discussions of diegesis tend to concern discourse/sjužet (in Russian Formalism) (vs. story/fabula). Diegesis is multi-levelled in narrative fiction. Genette distinguishes between three "diegetic levels". The extradiegetic level (the level of the narrative's telling) is, according to Prince, "external to (not part of) any diegesis." One might think of this as what we commonly understand to be the narrator's level, the level at which exists a narrator who is not part of the story being told. The diegetic level or intradiegetic level is understood as the level of the characters, their thoughts and actions. The metadiegetic level or hypodiegetic level is that part of a diegesis that is embedded in another one and is often understood as a story within a story, as when a diegetic narrator himself/herself tells a story. The classical distinction between the diegetic mode and the mimetic mode relate to the difference between the epos (or epic poetry) and drama. The "epos" relates stories by telling them through narration, while drama enacts stories through direct embodiment (showing). In terms of classical poetics, the cinema is an epic form that utilizes dramatic elements; this is determined by the technologies of the camera and editing. Even in a spatially and temporally continuous scene (mimicking the theatrical situation, as it were), the camera chooses where to look for us. In a similar way, editing causes us to jump from one place (and/or time) to another, whether it be somewhere else in the room, or across town. This jump is a form of narration; it is as if a narrator whispers to us: "meanwhile, on the other side of the forest". It is for this reason that the "story-world" in cinema is referred to as "diegetic"; elements that belong to the film's narrative world are diegetic elements. This is why, in the cinema, we may refer to the film's diegetic world. "Diegetic", in the cinema, typically refers to the internal world created by the story that the characters themselves experience and encounter: the narrative "space" that includes all the parts of the story, both those that are and those that are not actually shown on the screen (such as events that have led up to the present action; people who are being talked about; or events that are presumed to have happened elsewhere or at a different time). Thus, elements of a film can be "diegetic" or "non-diegetic". These terms are most commonly used in reference to sound in a film, but can apply to other elements. For example, an insert shot that depicts something that is neither taking place in the world of the film, nor is seen, imagined, or thought by a character, is a non-diegetic insert. Titles, subtitles, and voice-over narration (with some exceptions) are also non-diegetic. Film sound and music Sound in films is termed diegetic (termed source music by professionals in the radio, film and television industry) if it is part of the narrative sphere of the film. For instance, if a character in the film is playing a piano, or turns on a CD player, the resulting sound is diegetic. The cantina band sequence in the original Star Wars is an example of diegetic music in film, with the band playing instruments and swaying to the beat, as patrons are heard reacting to the second piece the band plays. If, on the other hand, music plays in the background but cannot be heard by the film's characters, it is termed non-diegetic or extradiegetic. Songs are commonly used in various film sequences to serve different purposes. They can be used to link scenes in the story where a character progresses through various stages toward a final goal. An example of this is in Rocky: Bill Conti's "Gonna Fly Now" plays non-diegetically as Rocky makes his way through his training regimen finishing on the top steps of the Philadelphia Museum of Art with his hands raised in the air. Mickey Mousing is an example of extradiegetic music. This distinction may be toyed with in order to break the fourth wall. In the Archer episode "Sea Tunt Part 1", Cheryl begins to hear music that would otherwise seem to be non-diegetic. She even comments: "Just ignore it; it's non-diegetic." In musical theater In musical theater, as in film, the term "diegesis" refers to the context of a musical number in a work's theatrical narrative. In typical operas or operettas, musical numbers are non-diegetic; characters are not aware that they are singing. In contrast, when a song occurs literally in the plot, the number is considered diegetic. Diegetic numbers are often present in backstage musicals. For example, in The Sound of Music, the song "Edelweiss" is diegetic, since the characters are aware they are singing. The character Maria is using the song to teach the children how to sing. In contrast, the song "How Do You Solve A Problem Like Maria?" is non-diegetic, since the musical material is external to the narrative. In both the 1936 and the 1951 film versions of Show Boat, as well as in the original stage version, the song "Bill" is diegetic. The character Julie LaVerne sings it during a rehearsal in a nightclub. A solo piano (played onscreen) accompanies her, and the film's offscreen orchestra (presumably not heard by the characters) sneaks in for the second verse of the song. Julie's other song in the film, "Can't Help Lovin' Dat Man" is also diegetic. In the 1936 film, it is supposed to be an old folk song known only to blacks; in the 1951 film it is merely a song which Julie knows; however, she and the captain's daughter Magnolia are fully aware that Julie is singing. When Julie, Queenie, and the black chorus sing the second chorus of the song in the 1936 version, they are presumably unaware of any orchestral accompaniment, but in the 1951 film, when Magnolia sings and dances this same chorus, she does so to the accompaniment of two deckhands on the boat playing a banjo and a harmonica. Two other songs in the 1936 Show Boat are also diegetic: "Goodbye My Lady Love" (sung by the comic dancers Ellie and Frank), and "After the Ball", sung by Magnolia. Both are interpolated into the film, and both are performed in the same nightclub in which Julie sings Bill. In the television series Buffy the Vampire Slayer, the episode entitled "Once More, with Feeling" toys with the distinction between diegetic and non-diegetic musical numbers. In this episode, the Buffy characters find themselves compelled to burst into song in the style of a musical. The audience is led to assume that this is a "musical episode", in which the characters are unaware that they are singing. It becomes clear that the characters are all too aware of their musical interludes, and that determining the supernatural causes of the singing is the focus of the episode's story. In video games In video games, "diegesis" comprises the narrative game world, its characters, objects and actions which can be classified as "intra-diegetic". Status icons, menu bars and other UI which are not part of the game world itself can be considered as "extra-diegetic"; a game character does not know about them even though for the player they may present crucial information. In this respect, these elements can be considered part of the narration provided by the game itself, although this will usually be a separate and distinct voice from that of the Story narrator, if there is one. A noted example of a diegetic interface in video games is that of the Dead Space series, in which the player-character is equipped with an advanced survival suit that projects holographic images to the character within the game's rendering engine that also serve as the game's user-interface to the player to show weapon selection, inventory management, and special actions that can be taken. - Gerald Prince, A Dictionary of Narratology, 2003, University of Nebraska Press, ISBN 0-8032-8776-3 - An etext of Plato's Republic is available from Project Gutenberg. The most relevant section is the following: "You are aware, I suppose, that all mythology and poetry is a narration of events, either past, present, or to come? / Certainly, he replied. / And narration may be either simple narration, or imitation, or a union of the two? / [...] / And this assimilation of himself to another, either by the use of voice or gesture, is the imitation of the person whose character he assumes? / Of course. / Then in this case the narrative of the poet may be said to proceed by way of imitation? / Very true. / Or, if the poet everywhere appears and never conceals himself, then again the imitation is dropped, and his poetry becomes simple narration."(Plato, Republic, Book III.) - Plato, Republic, Book III. - See also Pfister (1977, 2-3) and Elam: "classical narrative is always oriented towards an explicit there and then, towards an imaginary "elsewhere" set in the past and which has to be evoked for the reader through predication and description. Dramatic worlds, on the other hand, are presented to the spectator as "hypothetically actual" constructs, since they are "seen" in progress "here and now" without narratorial mediation. [...] This is not merely a technical distinction but constitutes, rather, one of the cardinal principles of a poetics of the drama as opposed to one of narrative fiction. The distinction is, indeed, implicit in Aristotle's differentiation of representational modes, namely diegesis (narrative description) versus mimesis (direct imitation)" (1980, 110-111). - Elam (1980, 110-111). - https://archive.org/details/GregoryKurczynskiOnOutsightRadioHours, Interview with a filmmaker on the diegetic role of music in film - Tach, Dave (13 March 2013). "Deliberately diegetic: Dead Space's lead interface designer chronicles the UI's evolution at GDC". Polygon. Retrieved 15 April 2015. - Aristotle. 1974. "Poetics". Trans. S.H. Butcher. In Dramatic Theory and Criticism: Greeks to Grotowski. Ed. Bernard F. Dukore. Florence, KY: Heinle & Heinle. ISBN 0-03-091152-4. p. 31-55. - Bunia, Remigius. 2010. "Diegesis and Representation: Beyond the Fictional World, on the Margins of Story and Narrative," Poetics Today 31.4, 679–720. doi:10.1215/03335372-2010-010. - Elam, Keir. 1980. The Semiotics of Theatre and Drama. New Accents Ser. London and New York: Methuen. ISBN 0-416-72060-9. - Pfister, Manfred. 1977. The Theory and Analysis of Drama. Trans. John Halliday. European Studies in English Literature Ser. Cambridige: Cambridge University Press, 1988. ISBN 0-521-42383-X. - Plato. c. 373 BC. Republic. Retrieved from Project Gutenberg on 2 September 2007. - Coyle, R. (2004). Pop goes the music track. Metro Magazine, 140, 94-95. - An Introduction to Film Analysis: Technique and Meaning in Narrative Film: Michael Ryan, Melissa Lenos: 9780826430021: Amazon.com: Books. The Continuum International Publishing Group, n.d. Web. 3 May 2013 - The dictionary definition of diegesis at Wiktionary

https://en.wikipedia.org/wiki/Diegesis

New information about what is inside Mars shows the Red Planet has a molten liquid-iron core, confirming the interior of the planet has some similarity to Earth and Venus. Researchers at NASA's Jet Propulsion Laboratory (JPL), Pasadena, Calif., analyzing three years of radio tracking data from the Mars Global Surveyor spacecraft, concluded Mars has not cooled to a completely solid iron core; rather its interior is made up of either a completely liquid iron core or a liquid outer core with a solid inner core. Their results are published in the March 7, 2003, online issue of the journal Science. "Earth has an outer liquid-iron core and solid inner core. This may be the case for Mars as well," said Dr. Charles Yoder, a planetary scientist at JPL and lead author on the paper. "Mars is influenced by the gravitational pull of the sun. This causes a solid body tide with a bulge toward and away from the sun (similar in concept to the tides on Earth). However, for Mars this bulge is much smaller, less than one centimeter. By measuring this bulge in the Mars gravity field we can determine how flexible Mars is. The size of the measured tide is large enough to indicate the core of Mars can not be solid iron but must be at least partially liquid," he explained. The team used Doppler tracking of a radio signal emitted by the Global Surveyor spacecraft to determine the precise orbit of the spacecraft around Mars. "The tidal bulge is a very small but detectable force on the spacecraft. It causes a drift in the tilt of the spacecraft's orbit around Mars of one-thousandth of a degree over a month," said Dr. Alex Konopliv, a planetary scientist at JPL and co-author on the paper. The researchers combined information from Mars Pathfinder on the Mars precession with the Global Surveyor tidal detection to draw conclusions about the Mars core, according to Dr. Bill Folkner, another co-author on the paper at JPL. The precession is the slow motion of the spin-pole of Mars as it moves along a cone in space (similar to a spinning top). For Mars it takes 170,000 years to complete one revolution. The precession rate indicates how much the mass of Mars is concentrated toward the center. A faster precession rate indicates a larger dense core compared to a slower precession rate. In addition to detection of a liquid core for Mars, the results indicate the size of the core is about one-half the size of the planet, as is the case for Earth and Venus, and the core has a significant fraction of a lighter element such as sulfur. In addition to measuring the Mars tide, Global Surveyor has been able to estimate the amount of ice sublimated, changed directly into a gaseous state, from one pole into the atmosphere and then accreted onto the opposite pole. "Our results indicate the mass change for the southern carbon- dioxide ice cap is 30 to 40 percent larger than the northern ice cap, which agrees well with the predictions of the global atmosphere models of Mars," said Yoder. The amount of total mass change depends on assumptions about the shape of the sublimated portion of the cap. The largest mass exchange occurs if one assumes the cap change is uniform or flat over the entire cap, while the lowest mass exchange corresponds to a conically shaped cap change. Mars at JPL Subscribe To SpaceDaily Express Mars May Be Much Older Or Younger Than Thought Buffalo - Jan 24, 2003 Research by a University at Buffalo planetary geologist suggests that generally accepted estimates about the geologic age of surfaces on Mars -- which influence theories about its history and whether or not it once sustained life -- could be way off. |The content herein, unless otherwise known to be public domain, are Copyright 1995-2016 - Space Media Network. All websites are published in Australia and are solely subject to Australian law and governed by Fair Use principals for news reporting and research purposes. AFP, UPI and IANS news wire stories are copyright Agence France-Presse, United Press International and Indo-Asia News Service. ESA news reports are copyright European Space Agency. All NASA sourced material is public domain. Additional copyrights may apply in whole or part to other bona fide parties. Advertising does not imply endorsement, agreement or approval of any opinions, statements or information provided by Space Media Network on any Web page published or hosted by Space Media Network. Privacy Statement All images and articles appearing on Space Media Network have been edited or digitally altered in some way. Any requests to remove copyright material will be acted upon in a timely and appropriate manner. Any attempt to extort money from Space Media Network will be ignored and reported to Australian Law Enforcement Agencies as a potential case of financial fraud involving the use of a telephonic carriage device or postal service.|

http://www.spacedaily.com/news/mars-general-03b.html

[/caption]Pierre Janssen was a French astronomer who discovered helium in 1868. He was observing a solar eclipse in India when he noticed the yellow spectral emission lines of the element. An English astronomer by the name of Norman Lockyer observed the same spectra and proposed the name helium after the Greek name for the sun, Helios. Helium can be observed at 587.49 nanometres in the spectrum of the chromosphere of the Sun. It was first thought that helium could only exist in/on the Sun because the spectral results could not be produced in the lab. That did not stop researchers form looking for it. In 1895, Sir William Ramsay discovered helium after treating cleveite, a uranium mineral, with mineral acids. Ramsay was looking for argon but, after separating nitrogen and oxygen from the gas liberated by sulfuric acid, noticed a bright-yellow line that matched the spectral line observed in the the Sun. Ramsey sent samples of the gas to Sir William Crookes and Sir Norman Lockyer who verified that it was helium. It was independently isolated from cleveite the same year by chemists Per Teodor Cleve and Abraham Langlet in Uppsala, Sweden, who were able to accurately determine its atomic weight. In a bit of irony or opportunity lost, American geochemist William Francis Hillebrand found the element prior to Ramsay’s discovery while testing a sample of the mineral uraninite. He attributed the lines to nitrogen and lost the claim to the discovery in the process. Several interesting properties of helium have been discovered in the ensuing years. In 1907, Ernest Rutherford and Thomas Royds demonstrated that an alpha particle is actually a helium nucleus. In 1908, helium was first liquefied by Dutch physicist Heike Kamerlingh Onnes by cooling the gas to less than one kelvin. He tried to solidify it by reducing the temperature more but failed because helium does not have a triple point temperature where the solid, liquid, and gas phases are at equilibrium. The element was eventually solidified in 1926 by his student Willem Hendrik Keesom. He managed to do so by subjecting helium to 25 atmospheres of pressure. Helium was one of the first elements to be found to have superfluidity. In 1938, Russian physicist Pyotr Leonidovich Kapitsa discovered that helium-4 has almost no viscosity at temperatures near absolute zero(superfluidity). In 1972, the same phenomenon was observed in helium-3 by American physicists Douglas D. Osheroff, David M. Lee, and Robert C. Richardson. Here on Universe Today we have a couple of great articles related to helium. One is about the possibility that white dwarfs can merge and form helium stars and the other is about liquid metal helium. Astronomy Cast offers a good episode about the energy spectra that we have been talking about in this article.

http://www.universetoday.com/53563/who-discovered-helium/

Introduced SpeciesThe Everglades National Park was established to protect the diverse natural habitats of the region which include freshwater marshes, hardwood hammocks, pinelands, cypress swamps, mangrove swamps, and estuaries. However, despite its status as a national park, the Everglades is threatened by introduced plants and animals. Introduced species are those organisms that are native to somewhere else that have been introduced to new areas through human activities. Many introduced species have detrimental effects on native flora and fauna due to lack of population controls such as predators and disease. As population numbers grow out of control, these introduced species are often referred to as invasive species. The introduction of species began in the late 1800s and has escalated since that time. These species continue to spread due to a lack of predators and disease, outcompeting native species for food and space. There are over 200 introduced species of plants that have been documented in the Everglades. These plants, including melaleuca (Melaleuca quinquenervia), Brazilian pepper (Schinus terebinthifolius), Australian pine (Casuarina equisetifolia), and Old World Climbing Fern (Lygodium microphyllum) displace native species and alter the natural habitat. Also detrimental to the habitats and communities of the Everglades are introduced species of wildlife. People have released unwanted pets into the Everglades including aquarium fishes, pythons, boa constrictors, parakeets, and parrots. Feral hogs also pose a major disturbance within the Everglades by digging native vegetation and disturbing archeological sites.Many species of fish originating from tropical and subtropical regions have been introduced into the freshwaters of the Everglades. Most can tolerate low to moderate salinities, allowing them to become established in brackish water estuaries. These fish have been introduced primarily through aquarium and aquaculture facilities, while some species have been released on purpose in hopes of establishing breeding populations. These fish include the Mayan cichlid (Cichlasoma urophthalmus), walking catfish (Clarias batrachus), Asian swamp eel (Monopterus albus), black acara (Cichlasoma bimaculatum), pike killifish (Belonesox belizanus), blue tilapia (Oreochromis aureus), spotted tilapia (Tilapia mariae), and oscar (Astronotus ocellatus). For more information, visit:

http://www.flmnh.ufl.edu/southflorida/regions/everglades/introduced-species

Madrid Teacher Resources Find Madrid educational ideas and activities Showing 1 - 20 of 333 resources Commonly Confused Words Exercise Accept or except? Advice or advise? Eminent or imminent? Which is which witch? In order to select the correct word to complete 20 sentences, learners get out their dictionaries and check the meaning and usage of the commonly confused pairs. 4th - 6th English Language Arts Comparing the New Madrid and San Andreas Fault Zones For this faults worksheet, students use an earthquake reference sheet to find the numbers for a modified Mercalli and Richter scale. They compare the San Andreas Fault zone and the New Madrid fault zone on the United States map. They... 7th - 10th Science Modal verbs of probability express what could or may happen. The class will look at 15 sentences and then choose which verb of probability fits best in each phrase. Then they write four phrases using accurate verbs in the present tense.... 4th - 8th English Language Arts Big Grammar Book With this comprehensive language arts resource in your arsenal, you'll never have to look for another grammar worksheet! Whether you're teaching kindergartners how to write the upper- and lower-case letters of the alphabet, or helping... K - 8th English Language Arts CCSS: Adaptable Do Journalists Shape or Report the News? Analyze the presence of negative stereotypes and biased reporting in news media, and how this affects one's understanding of other cultures. Learners read newspaper excerpts and quotes from famous personalities to discuss the power of... 9th - 12th Social Studies & History CCSS: Adaptable The Great Age of Exploration (1400-1550) Delve into the Age of Exploration with this activity-packed resource! Complete with a pre-test, discussion questions and quiz for a 30-minute video on the period, map activities, timeline of discoveries, vocabulary, etc. this is a... 7th - 12th Social Studies & History CCSS: Adaptable One Man’s Terrorist…Another Man’s Freedom Fighter Why is there no universal definition for terrorism? What tactics and objectives do terrorist groups share? Through an engaging and collaborative activity, as well as using rich informational texts and guided notes, lead your class... 10th - 12th Social Studies & History CCSS: Adaptable

http://www.lessonplanet.com/lesson-plans/madrid

One hundred years ago, two teams of explorers raced to be the first to reach the South Pole. Roald Engelbregt Gravning Amundsen reached the South Pole on December 14, 1911. Thirty-three days later on 17 January 1912 the Terra Nova Expedition led by Robert Falcon Scott arrived at the Pole in second place. At the same time in East Antarctica, the Australasian Antarctic Expedition led by Douglas Mawson was searching for the South Magnetic Pole. On their expeditions for King and country, Scott and Mawson carried out some of the first scientific studies in Antarctica. Scott's ill-fated expedition found fossils of Gondwanaland trees showing that Antarctica was once covered in lush forests. Even today, we tend to think of Antarctica as the last untouched wilderness preserved from human impact by International Treaty. However, despite its remoteness and vastness it is still affected by anthropogenic climate change. A paper to appear in the January issue of Global Change Biology shows how the dominant plants in Antarctica have been affected by modern climate change. In a handful of coastal Antarctic 'oases' void of permanent ice cover, lush moss beds grow during the short summer season from December to February using melt water from streams and lakes. Up until now, measuring the seasonal growth rate of these plants has been extremely difficult and hence it was impossible to assess the impact of our changing climate. This research, conducted by a team of environmental scientists from the University of Wollongong (UOW) and nuclear physicists from the Australian Nuclear Science and Technology Organisation (ANSTO), shows how the increased concentration of radiocarbon in the atmosphere resulting from nuclear weapons testing (mostly in the late 1950s and early 1960s, called the 'the bomb spike') can be used to accurately date the age of the moss shoots along their stems in a similar way to tree-rings. Professor Sharon Robinson from UOW's Institute for Conservation Biology and Environmental Management (School of Biological Sciences) said the team found that that most of the plants were growing 50 years ago when nuclear testing was at its peak. In some species the peak of the radiocarbon bomb spike was found just 15 mm from the top of the 50 mm shoot suggesting that these plants may be more than 100 years old. 'Accurate dating along the moss stem allows us to determine the very slow growth rates of these mosses (ranging from 0.2 to 3.5 mm per year). Remarkably, these plants were already growing during the heroic age of Antarctic exploration. In terms of age these mosses are effectively the old growth forests of Antarctica -- in miniature," Professor Robinson said. Although increased temperature and precipitation in the polar regions due to climate change are predicted to increase growth rates, the scientists found that at some sites growth rates have declined since the 1980s. They suggest that this is likely due to moss beds drying out, which appears to be caused by increased wind speeds around Antarctica that are linked to the Antarctic ozone hole. In the 100 years since the start of scientific research in Antarctica, contamination of Earth's atmosphere with increased radioactivity due to nuclear weapons testing has led to radiocarbon labelling of Antarctic plants. "This has allowed scientists to show that climate change has made the driest continent on Earth an even harsher environment for plant life," Professor Robinson said. |Contact: Sharon Robinson|

http://www.bio-medicine.org/biology-news-1/Climate-change-stunting-growth-of-century-old-Antarctic-moss-shoots-22643-1/

A graben is a valley with a distinct escarpment on each side caused by the displacement of a block of land downward. Graben often occur side-by-side with horsts. Horst and graben structures indicate tensional forces and crustal stretching. Graben are produced from parallel normal faults, where the displacement of the hanging wall is downward, while that of the footwall is upward. The faults typically dip toward the center of the graben from both sides. Horsts are parallel blocks that remain between graben; the bounding faults of a horst typically dip away from the center line of the horst. Single or multiple graben can produce a rift valley. In many rifts the graben are asymmetric, with a major fault along only one of the boundaries, and these are known as half-graben. The polarity (throw direction) of the main bounding faults typically alternate along the length of the rift. The asymmetry of a half-graben strongly affects syntectonic deposition. Comparatively little sediment enters the half-graben across the main bounding fault, due to the effects of footwall uplift on the drainage systems. The exception is at any major offset in the bounding fault, where a relay ramp may provide an important sediment input point. Most of the sediment will enter the half-graben down the unfaulted hanging wall side (e.g. Lake Baikal). - The Basin and Range Province of southwestern North America is an example of multiple horst/graben structures, including Death Valley, with Salt Lake Valley being the easternmost and Owens Valley being the westernmost. - The Rio Grande Rift Valley in Colorado/New Mexico/Texas of the United States - The Rhine valley to the north of Basel, Switzerland - The Oslo graben around Oslo, Norway - The East African Rift Valley - The Saguenay Graben, Quebec, Canada - The Narmada River valley in central India - The lower Godavari River valley in southern India - The Ottawa-Bonnechere Graben in Ontario and Quebec, Canada - The Lambert Graben in Antarctica - Gulf St Vincent in South Australia, Australia - The Guanabara Bay in Rio de Janeiro, Brazil - The Central Lowlands (Midland Valley) of Scotland - Baikal Rift Zone, Siberia, Russia - Lake Tahoe, California and Nevada, US - Santa Clara Valley, California, US - Guatemala city valley, Guatemala, GT - Büyük Menderes Graben, Turkey - The Unzen Graben in Japan - The Republic Graben in Republic, Washington

https://en.wikipedia.org/wiki/Graben

One of the biggest knocks against cellphones is they require small amounts of rare earth elements: gallium, indium and arsenic, for example, that are both scarce and expensive. But what if you could make a phone out of a more common element, like carbon? Researchers are taking slow but sure steps toward building the innards of a cellphone out of carbon nanotubes, a structure that resembles a microscopic sheet of chicken wire rolled into a cylinder. These cylinders can be used to either conduct electricity or store energy. At the Technical University of Denmark, Jakob Wagner and colleagues have found a better way to build carbon nanotubes that could lead to their use as a semiconductor, a key component of all electronic circuit parts found in both cellphones and laptops. Carbon nanotubes have properties of both a metal and a semiconductor, depending on how they are rolled. “The breakthrough here is that we are able to control the production of nanotubes whether they are metallic or semiconducting,” Wagner said. “That’s important because if you want to use them in cellphones, we have to make sure they are either one or the other. The prospect is to use semiconducting carbon nanotubes as a substitute for gallium.” Warner published his work earlier this month in the Nature publication Scientific Reports. The next step is to be able to produce large amounts of semiconducting carbon nanotubes that could be made into an electronic device, Wagner said. “It will not be tomorrow, let’s say 10 years,” he said. But at IBM, researchers like James Hannon are working to speed up that lab-to-prototype timescale. Hannon says that Wagner’s finding is an important step, but it needs to be replicated on larger-diameter carbon nanotubes. "This is a nice scientific demonstration, but not in the range that would be used in a logic application," said Hannon, manager of IBM’s carbon electronics group in Yorktown Heights, N.Y. "I’d like to see if this technique could work for larger diameter tubes as well." Last year, Hannon and his IBM colleagues announced they had built memory and microprocessing chips using carbon nanotubes. He said the tough thing is getting them to lie down in straight lines, but they overcame this obstacle by creating special grooves etched into the silicon chip surface and a bonding agent. Hannon says the two challenges with carbon nanotubes is figuring out how to place them and how to separate the semiconducting ones from the metallic ones, which are thrown away. A separate team at North Carolina State University recently reported they were able to integrate carbon nanotubes into a flexible scaffold for a silicon-based battery that would last longer than existing lithium ion batteries. Hannon says he expects carbon nanotubes to play a big role in electronic devices in a few more years of testing. "Our mandate is that this stuff has to be ready pretty soon,” he said. Image: Flickr, Gonzalo Baeza H - Where Did Those Gravitational Waves Come From? There's a Map - New Heart-Shaped Hawaiian Fruits Discovered - DNews: Gravitational Waves Confirmed: A Historic Discovery - Fish Can Sense Touch with Their Fins This article originally published at Discovery News here

http://mashable.com/2013/08/30/carbon-based-cellphones/

Redistribution of income and wealth Redistribution of income and redistribution of wealth are respectively the transfer of income and of wealth (including physical property) from some individuals to others by means of a social mechanism such as taxation, charity, welfare, land reform, monetary policies, confiscation, divorce or tort law. The term typically refers to redistribution on an economy-wide basis rather than between selected individuals, and it always refers to redistributions from those who have more to those who have less. The desirability and effects of redistribution are actively debated on ethical and economic grounds. The subject includes analysis of its rationales, objectives, means, and policy effectiveness. A 2003 survey among American economists found that 71.2% of them support redistribution, while 20.4% oppose it, 7.2% had mixed feelings. The concept of wealth redistribution is old, and goes back as far as recorded human history. In ancient times, this was known as a Palace economy. These economies were centrally based around the administration, so the dictator or pharoah had both the ability and the right to say who did(and did not) get special treatment. Another early form of wealth redistribution occurred in the early American colonies under the leadership of William Bradford. Bradford records in his diary that this "common course" bred confusion, discontent, distrust, and the colonists looked upon it as a form of slavery. Role in economic systems Different types of economic systems feature vastly different levels of interventionism to redistribute income, depending on how unequal the initial distribution of income in their economies is. Free-market capitalist economies tend to feature high degrees of income redistribution, but Japan's government engages in much less redistribution because its initial wage distribution is much more equal. Likewise, the socialist planned economies of the former Soviet Union and Eastern bloc had very little income redistribution because private capital and land income, the major drivers of income inequality in capitalist systems, did not exist in these economies; and the government set wages in these economies. Modern forms of redistribution Today, income redistribution occurs in some form in most democratic countries. In a progressive income tax system, a high income earner will pay a higher tax rate than a low income earner. Another taxation-based method of redistributing income is the negative income tax. Two other common types of governmental redistribution of income are subsidies and vouchers (such as food stamps). These transfer payment programs are funded through general taxation, but benefit the poor, who pay fewer or no taxes. While the persons receiving transfers from such programs may prefer to be directly given cash, these programs may be more palatable to society than cash assistance, as they give society some measure of control over how the funds are spent. Wealth redistribution can be implemented through land reform that transfers ownership of land from one category of people to another, or through inheritance taxes or direct wealth taxes. Before-and-after Gini coefficients for the distribution of wealth can be compared. The objectives of income redistribution are to increase economic stability and opportunity for the less wealthy members of society and thus usually include the funding of public services. One basis for redistribution is the concept of distributive justice, whose premise is that money and resources ought to be distributed in such a way as to lead to a socially just, and possibly more financially egalitarian, society. Another argument is that a larger middle class benefits an economy by enabling more people to be consumers, while providing equal opportunities for individuals to reach a better standard of living. Seen for example in the work of John Rawls, another argument is that a truly fair society would be organized in a manner benefiting the least advantaged, and any inequality would be permissible only to the extent that it benefits the least advantaged. Some[who?] argue that wealth and income inequality are a cause of economic crises, and that reducing these inequalities is one way to prevent or ameliorate economic crises, with redistribution thus benefiting the economy overall. This view was associated with the underconsumptionism school in the 19th century, now considered an aspect of some schools of Keynesian economics; it has also been advanced, for different reasons, by Marxian economics. It was particularly advanced in the US in the 1920s by Waddill Catchings and William Trufant Foster. There is currently a great debate concerning the extent to which the world's extremely rich have become richer over recent decades. Thomas Piketty's Capital in the Twenty-First Century is at the forefront, critiqued in certain publications such as The Economist. |This section requires expansion. (November 2015)| Economic effects of inequality Using statistics from 23 developed countries and the 50 states of the US, British researchers Richard G. Wilkinson and Kate Pickett show a correlation between income inequality and higher rates of health and social problems (obesity, mental illness, homicides, teenage births, incarceration, child conflict, drug use), and lower rates of social goods (life expectancy, educational performance, trust among strangers, women's status, social mobility, even numbers of patents issued per capita), on the other. The authors argue inequality leads to the social ills through the psychosocial stress, status anxiety it creates. A 2011 report by the International Monetary Fund by Andrew G. Berg and Jonathan D. Ostry found a strong association between lower levels of inequality and sustained periods of economic growth. Developing countries (such as Brazil, Cameroon, Jordan) with high inequality have "succeeded in initiating growth at high rates for a few years" but "longer growth spells are robustly associated with more equality in the income distribution." The socialist economists John Roemer and Pranab Bardhan criticize redistribution via taxation in the context of Nordic-style social democracy, highlighting its limited success at promoting relative egalitarianism and its lack of sustainability. They point out that social democracy requires a strong labor movement to sustain its heavy redistribution, and that it is unrealistic to expect such redistribution to be feasible in countries with weaker labor movements. They point out that, even in the Scandinavian countries, social democracy has been in decline since the labor movement weakened. Instead, Roemer and Bardham argue that changing the patterns of enterprise ownership and market socialism, obviating the need for redistribution, would be more sustainable and effective at promoting egalitarianism. Marxian economists argue that social democratic reforms - including policies to redistribute income - such as unemployment benefits and high taxes on profits and the wealthy create more contradictions in capitalism by further limiting the efficiency of the capitalist system via reducing incentives for capitalists to invest in further production. In the Marxist view, redistribution cannot resolve the fundamental issues of capitalism - only a transition to a socialist economy can. - Economic policy - Poverty reduction - Robin Hood - Robin Hood tax - Social inequality - Redistribution (cultural anthropology) - Wealth concentration - "Redistribution". Stanford Encyclopedia of Philosophy. Stanford University. 2 July 2004. Retrieved 13 August 2010. The social mechanism, such as a change in tax laws, monetary policies, or tort law, that engenders the redistribution of goods among these subjects - F.A. Cowell ( 2008). "redistribution of income and wealth,"The New Palgrave Dictionary of Economics, 2nd Edition, TOC. - Rugaber, Christopher S.; Boak, Josh (January 27, 2014). "Wealth gap: A guide to what it is, why it matters". AP News. Retrieved January 27, 2014. - Klein, Daniel B.; Stern, Charlotta (2006). "Economists’ policy views and voting" (PDF). Public Choice (Springer) 126 (3-4): 337. doi:10.1007/s11127-006-7509-6. - de Blois, Lukas; R.J. van der Spek; Susan Mellor (translator) (1997). An Introduction to the Ancient World. Routledge. pp. 56–60. ISBN 0-415-12773-4. Cite uses deprecated parameter - William Bradford - History of Plymouth Plantation, p. 135 - Rosser, Mariana V. and J Barkley Jr. (July 23, 2003). Comparative Economics in a Transforming World Economy. MIT Press. p. 11. ISBN 978-0262182348. Economies vary based on the extent to which and the methods by which governments intervene to redistribute income. This depends partly on how unequal income is to begin with before any redistributive policies are implemented. Thus the Japanese government does much less redistributing than the governments of many other capitalist countries because Japan has a more equal distribution of wages than most other capitalist countries. Command socialist economies also have had less income redistribution because governments initially control the distribution of income by setting wages and forbidding capital or land income. - Harvey S. Rosen & Ted Gayer, Public Finance pp. 271–72 (2010). - Marx, K. A Contribution to the Critique of Political Economy. Progress Publishers, Moscow, 1977 - (Dorfman 1959) - Allgoewer, Elisabeth (May 2002). "Underconsumption theories and Keynesian economics. Interpretations of the Great Depression" (PDF). Discussion paper no. 2002-14. University of St. Gallen. - Forget the 1%; Free Exchange, The Economist, 8 November 2014, p79. - Famine, Affluence, and Morality - Fighting Poverty - Statistics and graphs from Wilkinson and Pickett research. - The Spirit Level: how 'ideas wreckers' turned book into political punchbag| Robert Booth| The Guardian| 13 August 2010 - Inequality and Unsustainable Growth: Two Sides of the Same Coin? Andrew G. Berg and Jonathan D. Ostry| IMF STAFF DISCUSSION NOTE | April 8, 2011 - Berg, Andrew G.; Ostry, Jonathan D. (2011). "Equality and Efficiency". Finance and Development (International Monetary Fund) 48 (3). Retrieved September 10, 2012. - Plotnick, Robert (1986) "An Interest Group Model of Direct Income Redistribution", The Review of Economics and Statistics, vol. 68, #4, pp. 594–602. - Market socialism, a case for rejuvenation, by Pranab Bardhan and Johen E. Roemer. 1992. Journal of Economic Perspectives, Vol. 6, No. 3, pp. 104: "Since it (social democracy) permits a powerful capitalist class to exist (90 percent of productive assets are privately owned in Sweden), only a strong and unified labor movement can win the redistribution through taxes that is characteristic of social democracy. It is idealistic to believe that tax concessions of this magnitude can be effected simply through electoral democracy without an organized labor movement, when capitalists organize and finance influential political parties. Even in the Scandinavian countries, strong apex labor organizations have been difficult to sustain and social democracy is somewhat on the decline now." - Market Socialism: The Debate Among Socialists, by Schweickart, David; Lawler, James; Ticktin, Hillel; Ollman, Bertell. 1998. (P.60-61): "The Marxist answers that...it involves limiting the incentive system of the market through providing minimum wages, high levels of unemployment insurance, reducing the size of the reserve army of labour, taxing profits, and taxing the wealthy. As a result, capitalists will have little incentive to invest and the workers will have little incentive to work. Capitalism works because, as Marx remarked, it is a system of economic force (coercion)." - Levy, Frank (2008). "Distribution of Income". In David R. Henderson (ed.). Concise Encyclopedia of Economics (2nd ed.). Indianapolis: Library of Economics and Liberty. ISBN 978-0865976658. OCLC 237794267. - Small calculus of inequality measures

https://en.wikipedia.org/wiki/Property_redistribution

Discover what myths reveal about ancient and contemporary cultures. - Grades: PreK–K, 1–2, 3–5, 6–8, 9–12 Describes a lesson in identifying and charting the characteristics of a myth through reading and making inferences. Allow your class to take a journey into the world of Greek myths. Students learn vocabulary from Ancient Greece that will help them to understand roots of modern English words. Mythology is not my strong suit, so when I stumbled across a Greek mythology readers theater book, I was ecstatic. Read on to find out how to incorporate this activity into your classroom. Presents a lesson using a mythological hero chart. Students chart character traits based on readings of Ancient Greek myths. Proposes a lesson in which students write journalist pieces about events from Greek myths. In this lesson unit on ancient Greece, students compare three myths and create their own original myth. Students combine their journalistic skills with their knowledge of Greek myth to write a fictional article about mythical characters taking over modern L.A. Online Learning Activities This four-step workshop hosted by an award-winning myth writer offers writing strategies and exercises to help students craft successful myths. In this blend of Greek mythology and modern pop culture, things become clearer when Percy discovers he is the son of Poseidon. But trouble starts all over again when Percy is sent on a quest to prevent war on Mount Olympus. Three exciting learning activities to complement author Rick Riordan's modern myth Use these 15 questions to help students get more out of the experience of reading Rick Riordan's book, Sea of Monsters. Two drawing assignments and a memory activity to follow reading the book by Rick Riordan Students choose a myth from Mary Pope Osbourne's book to dramatize. Booktalk for author Ross Collins' Medusa Jones, the story of a young outcast named Medusa, complete with hilarious illustrations! Find folk tales and legends from all over the world in this book list for grades PreK-5. These books for middle- and high-school students put a new spin on fairy tales and legends.

http://www.scholastic.com/teachers/collection/myths

Our galaxy's dark matter is clumpier than once thought, according to a new computer simulation. The model, created by one of the most powerful supercomputers in the world, shows that the spherical halo of dark matter that envelopes the Milky Way contains dense clumps and streams of the mysterious stuff, even in the neighborhood of our solar system. "In previous simulations, this region came out smooth, but now we have enough detail to see clumps of dark matter," said researcher Piero Madau, an astrophysicist at the University of California, Santa Cruz. Dark matter, which scientists can only detect by noting its gravitational effect, is thought to make up about 85 percent of the matter in the universe. Its composition remains a mystery, though some scientists think it's made up of hypothetical particles called WIMPs (weakly interacting massive particles), which could annihilate each other and emit gamma rays when they collide. The new simulation, described in the Aug. 7 issue of the journal Nature, implies that dark matter could be detected by the recently launched Gamma-ray Large Area Space Telescope (GLAST). "That's what makes this exciting," Madau said. "Some of those clumps are so dense they will emit a lot of gamma rays if there is dark matter annihilation, and it might easily be detected by GLAST." So far, though many teams have been looking for WIMP particles, no one has conclusively detected them. "There are several candidate particles for cold dark matter, and our predictions for GLAST depend on the assumed particle type and its properties," said Juerg Diemand, a postdoctoral fellow at UCSC who led the new research. "For typical WIMPs, anywhere from a handful to a few dozen clear signals should stand out from the gamma-ray background after two years of observations. That would be a big discovery for GLAST." The model took about one month to run on the Jaguar supercomputer at Oak Ridge National Laboratory in Tennessee. By following the gravitational interactions of more than a billion parcels of dark matter over 13.7 billion years, the computer could predict how the dark matter in the universe developed over time based on leading theories of how dark matter interacts. "It simulates the dark matter distribution from near the time of the Big Bang until the present epoch, so practically the entire age of the universe, and focuses on resolving the halo around a galaxy like the Milky Way," Diemand said. The research was funded by the U.S. Department of Energy, NASA and the Swiss National Science Foundation. - Video: Dark Matter in 3-D - Vote: The Strangest Things in Space Mysteries: Where is the Rest of the Universe?

http://www.space.com/5705-milky-dark-matter-clumpier-thought.html

Definition of Saxon in English: 1A member of a people that inhabited parts of central and northern Germany from Roman times, many of whom conquered and settled in much of southern England in the 5th-6th centuries. - There was relative peace with British rule over the western half of the country and Germanic rule in the east for the next fifty years, and it seems likely that the Britons may even have regained some areas of central England from the Saxons. - Faced with invasion by a coalition of Picts and Saxons, the Roman citizens of Britain appeal to the Emperor for help; but Honorius is in no position to aid them. - When Charlemagne conquered the Saxons, he extended his empire to the borders of Viking realms: specifically, to Friesland in southern Denmark. 1.1A native of modern Saxony in Germany. adjectiveBack to top 1Relating to the Anglo-Saxons, their language (Old English), or their period of dominance in England (5th-11th centuries). - Wales is contiguous to England and had been the subject of Saxon raids for centuries. - For much of the Saxon period it was probably fairly wide and marshy, perhaps acting as a separator between Westwyk and Conesford. - Across much of midland England wide-ranging changes took place in the countryside in the late Saxon period. 1.1Relating to or denoting the style of early Romanesque architecture preceding the Norman in England. - The site develops with the construction of an aisled Late Saxon timber hall, which was one of King Cnut's royal manors. - Within the church, parts of the Saxon north wall can be seen above the Norman arcade. - On the outside of the north wall, (about a third of the way down the Nave), the remains of a Saxon doorway can be seen, complete with round headed arch and jambs of flint. Pronunciation: /ˈsaks(ə)nʌɪz/(also Saxonise) verb Words that rhyme with Saxonflaxen, Jackson, klaxon, Sachsen, waxen Definition of Saxon in: - US English dictionary What do you find interesting about this word or phrase? Comments that don't adhere to our Community Guidelines may be moderated or removed.

http://www.oxforddictionaries.com/definition/english/saxon

This concept introduces students to inductive reasoning and provides many examples of inductive reasoning. Inductive Reasoning from Patterns Interactive This video gives more detail about the mathematical principles presented in Inductive Reasoning. This video shows how to work step-by-step through one or more of the examples in Inductive Reasoning. A list of student-submitted discussion questions for Inductive Reasoning from Patterns. To activate prior knowledge, make personal connections, reflect on key concepts, encourage critical thinking, and assess student knowledge on the topic prior to reading using a Quickwrite. To stress understanding of a concept by summarizing the main idea and applying that understanding to create visual aids and generate questions and comments using a Concept Matrix. To activate prior knowledge, to generate questions about a given topic, and to organize knowledge using a KWL Chart. Learn how inductive reasoning is used throughout the sciences, from medicine to zoology. Symbolic notation used in logic, inductive reasoning (patterns), and deductive reasoning are the focus of this study guide.

http://www.ck12.org/geometry/Inductive-Reasoning-from-Patterns/

Brood parasites are organisms that rely on others to raise their young. The strategy appears among insects, fishes, and birds. The brood parasite manipulates a host, either of the same or of another species, to raise its young as if it were its own. Brood parasitism relieves the parasitic parents from the investment of rearing young or building nests for the young, enabling them to spend more time on other activities such as foraging and producing further offspring. Bird parasite species mitigate the risk of egg loss by distributing eggs amongst a number of different hosts. As this behaviour damages the host, it often results in an evolutionary arms race between parasite and host. - 1 Birds - 2 Parental-care parasitism - 3 Fish - 4 Insects - 5 See also - 6 References - 7 External links In many monogamous bird species, there are extra-pair matings resulting in males outside the pair bond siring offspring and used by males to escape from the parental investment in raising their offspring. This form of cuckoldry is taken a step further when females of the goldeneye, Bucephala clangula often lay their eggs in the nests of other individuals. Intraspecific brood parasitism is seen in a number of duck species, where females often lay their eggs in the nests of others. Interspecific brood-parasites include the Old World cuckoo, cowbirds, black-headed ducks, and some New World cuckoos in the Americas, and indigobirds, whydahs, and honeyguides in Africa. Seven independent origins of obligate interspecific brood parasitism in birds have been proposed. While there is still some controversy over when and how many origins of interspecific brood parasitism have occurred, recent phylogenetic analyses suggest two origins in Passeriformes (once in New World cowbirds: Icteridae, and once in African Finches: Viduidae); three origins in Old World and New World cuckoos (once in Cuculinae, Phaenicophaeinae, and in Neomorphinae-Crotophaginae); a single origin in Old World honeyguides (Indicatoridae); and in a single species of waterfowl, the black-headed duck (Heteronetta atricapilla). Most avian brood parasites are specialists which will only parasitize a single host species or a small group of closely related host species, but four out of the five parasitic cowbirds are generalists, which parasitize a wide variety of hosts; the brown-headed cowbird has 221 known hosts. They usually only lay one egg per nest, although in some cases, particularly the cowbirds, several females may use the same host nest. The common cuckoo presents an interesting case in which the species as a whole parasitizes a wide variety of hosts, including the reed warbler and dunnock, but individual females specialize in a single species. Genes regulating egg coloration appear to be passed down exclusively along the maternal line, allowing females to lay mimetic eggs in the nest of the species they specialize in. Females generally parasitize nests of the species which raised them. Male common cuckoos will fertilize females of all lines, maintaining sufficient gene flow among the different maternal lines. The mechanisms of host selection by female cuckoos are somewhat unclear, though several hypotheses have been suggested in attempt to explain the choice. These include genetic inheritance of host preference, host imprinting on young birds, returning to place of birth and subsequently choosing a host randomly ("natal philopatry"), choice based on preferred nest site (nest-site hypothesis), and choice based on preferred habitat (habitat-selection hypothesis). Of these hypotheses the nest-site selection and habitat selection have been most supported by experimental analysis. Adaptations for parasitism Among specialist avian brood parasites, mimetic eggs are a nearly universal adaptation. There is even some evidence that the generalist brown-headed cowbird may have evolved an egg coloration mimicking a number of their hosts. Most avian brood parasites will remove a host egg when they lay one of their own in a nest. Depending upon the species, this can happen either in the same visit to the host nest or in a separate visit before or after the parasitism. This both prevents the host species from realizing their nest has been parasitized and reduces competition for the parasitic nestling once it hatches. Most avian brood parasites have very short egg incubation periods and rapid nestling growth. This gives the parasitic nestling a head start on growth over its nestmates, allowing it to outcompete them. In many brood parasites, such as cuckoos and honeyguides, this short egg incubation period is due to internal incubation periods up to 24 hours longer in cuckoos than hosts. Some non-parasitic cuckoos also have longer internal incubation periods, suggesting that this longer internal incubation period was not an adaptation following brood parasitism, but predisposed birds to become brood parasites. In cases where the host nestlings are significantly smaller than the parasite nestling, the host nestlings will often starve to death. Some brood parasites will eliminate all their nestmates shortly after hatching, either by ejecting them from the nest or killing them with sharp mandible hooks which fall off after a few days. It has often been a question as to why the majority of the hosts of brood parasites care for the nestlings of their parasites. Not only do these brood parasites usually differ significantly in size and appearance, but it is also highly probable that they reduce the reproductive success of their hosts. The "mafia hypothesis" evolved through studies in an attempt to answer this question. This hypothesis revolves around host manipulations induced by behaviors of the brood parasite. Upon the detection and rejection of a brood parasite's egg, the host's nest is depredated upon, its nest destroyed and nestlings injured or killed. This threatening response indirectly enhances selective pressures favoring aggressive parasite behavior that may result in positive feedback between mafia-like parasites and compliant host behaviors. There are two avian species that have been speculated to portray this mafia-like behavior: the brown-headed cowbird of North America, Molothrus ater, and the great spotted cuckoo of Europe, Clamator glandarius. The great spotted cuckoo lays the majority of its eggs in the nests of the European magpie, Pica pica. It has been observed that the great spotted cuckoo repeatedly visits the nests that it has parasitised, a precondition for the mafia hypothesis. An experiment was run by Soler et al. from April to July 1990 – 1992 in the high-altitude plateau Hoya de Guadix, Spain. They observed the effects of the removal of cuckoo eggs on the reproductive success of the magpie and measured the magpie's reaction; the egg was considered accepted if it remained in the nest, ejected if gone in between visits, or abandoned if the eggs were present but cold. If any nest contents were gone between consecutive visits, the nests were considered to have been depredated. The magpie's reproductive success was measured by number of nestlings that survived to their last visit, which was just before the nestling had been predicted to fledge from the nest. The results from these experiments show that after the removal of the parasitic eggs from the great spotted cuckoo, these nests are predated at much higher rates than those where the eggs were not removed. Through the use of plasticine eggs that model those of the magpie, it was confirmed that the nest destruction was caused by the great spotted cuckoo. This destruction benefits the cuckoo, for the possibility of re-nesting by the magpie allows another chance for the cuckoo egg to be accepted. Another similar experiment was done in 1996–2002 by Hoover et al. on the relationship between the parasitic brown-headed cowbird and a host, the prothonotary warbler, Protonotaria citrea. In their experiment, researchers manipulated the cowbird egg removal and cowbird access to the predator proof nests of the warbler. They found that 56% of egg ejected nests were depredated upon in comparison to 6% of non-ejected nests when cowbirds were not prevented from getting to the hosts' nest. Of the nests that were rebuilt by hosts that had previously been predated upon, 85% of those were destroyed. The number of young produced by the hosts that ejected eggs dropped 60% compared to those that accepted the cowbird eggs. In this hypothesis, female cuckoos select a group of host species with similar nest sites and egg characteristics to her own. This population of potential hosts is monitored and a nest is chosen from within this group. Research of nest collections has illustrated a significant level of similarity between cuckoo eggs and typical eggs of the host species. A low percentage of parasitized nests were shown to contain cuckoo eggs not corresponding to the specific host egg morph. In these mismatched nests a high percent of the cuckoo eggs were shown to correlate to the egg morph of another host species with similar nesting sites. This has been pointed to as evidence for nest- site selection. A criticism of the hypothesis is that it provides no mechanism by which nests are chosen, or which cues might be used to recognize such a site. Parental-care parasitism emphasizes the relationship between the host and the parasite in brood parasitism. Parental-care parasitism occurs when individuals raise offspring of other unrelated individuals. The host are the parents of offspring and the parasites are individuals who take advantage of either the nest or eggs within the family construct. Such dynamics occur when the parasites attempt to reduce their parental investment so they can invest the extra energy into other endeavors. Cost of the hosts Given the detrimental effects avian brood parasites can have on their hosts' reproductive success, host species have come up with various defenses against this unique threat. Given that the cost of egg removal concurrent with parasitism is unrecoverable, the best defense for hosts is avoiding parasitism in the first place. This can take several forms, including selecting nest sites which are difficult to parasitize, starting incubation early so they are sitting on the nests when parasites visit them early in the morning, and aggressive territorial defense. Birds nesting in aggregations can also benefit from group defense. The hosts reject offspring The host may be the one that ultimately ends up raising offspring after they return from foraging. Once parasitism has occurred, the next most optimal defense is to eject the parasitic egg. According to parental investment theory, the host can possibly adopt some defense to protect their own eggs if they distinguish which eggs are not theirs. Recognition of parasitic eggs is based on identifying pattern differences or changes in the number of eggs. This can be done by grasp ejection if the host has a large enough beak, or otherwise by puncture ejection. Ejection behavior has some costs however, especially when host species have to deal with mimetic eggs. In that case, hosts will inevitably mistake one of their own eggs for a parasite egg on occasion and eject it. In any case, hosts will sometimes damage their own eggs while trying to eject a parasite egg. Among hosts not exhibiting parasitic egg ejection, some will abandon parasitized nests and start over again. However, at high enough parasitism frequencies, this becomes maladaptive as the new nest will most likely become reparasitized. Other behavior can include modifying the nest to exclude the parasitic egg, either by weaving over the egg or in some cases rebuilding a new nest over the existing one. For instance, American coots might kick the parasites’ eggs out, or build a new nest beside the brood nests where the parasites’ babies starve to death due to lack of food. Cost of the parasites While parental-care parasitism significantly increased the breeding number of the parasite, only about half of the parasite eggs survived. Parasitism for the individual (the brood parasite) also has significant drawbacks. As an example, the parasitic offspring of the bearded tits, Panurus biarmicus, compared to the offspring in non-parasitic nests, tend to develop much more slowly and often don’t reach full maturity. Parasitic females however can adopt either floater traits or nesting traits. Floater females are entirely dependent on others to raise their eggs because they do not have their own nests. Hence, they reproduce significantly less because the hosts reject their ‘intruder’ eggs or they may just miss the egg-laying period of the bird they are trying to pass their eggs to. Nesting females who have their own nests may also be parasitic due to temporary situations like sudden loss of nests, or they lay surplus eggs, which overload their parental care ability. The hosts raise offspring Sometimes hosts are completely unaware that they are caring for a bird that is not their own. This most commonly occurs because the host cannot differentiate the parasitic eggs from their own. It may also occur when hosts temporarily leave the nest after laying the eggs. The parasites lay their own eggs into these nests so their nestlings share the food provided by the host. It may occur in other situations. For example, female eiders would prefer to lay eggs in the nests with one or two existing eggs of others because the first egg is the most vulnerable to predators. The presence of others’ eggs reduces the probability that a predator will attack her egg when a female eider leaves the nest after laying the first egg. Sometimes, the parasitic offspring kills the host nest-mates during competition for resources. As an example, the parasite offspring of the cowbird chick kill the host nest-mates if food intake for each of them is low, but do not do so if the food intake is adequate, as a result of their interactions with co-inhabitants of the nest. A mochokid catfish of Lake Tanganyika, Synodontis multipunctatus, is a brood parasite of several mouthbrooding cichlid fishes. The catfish eggs are incubated in the host's mouth, and in the manner of cuckoos hatch before the host's own eggs. The young catfish eat the host fry inside the host's mouth, effectively taking up virtually the whole of the host's parental investment. A cyprinid minnow, Pungtungia herzi is a brood parasite of the Serranid freshwater perch Siniperca kawamebari, which live in the south of the Japanese islands of Honshu, Kyushu and Shikoku, and in South Korea. Host males guard territories against intruders during the breeding season, creating a patch of reeds as a spawning site or "nest". Females (one or more per site) visit the site to lay eggs, which the male then defends. The parasite's eggs are smaller and stickier than the host's. 65.5% of host sites were parasitised in a study area. There are many different types of cuckoo bees, all of which lay their eggs in the nest cells of other bees, but they are normally referred to as kleptoparasites, rather than as brood parasites, because the immature stages are almost never fed directly by the adult hosts. Examples of cuckoo bees are Coelioxys rufitarsis, Melecta separata, Bombus bohemicus, Nomada and Epeoloides. Kleptoparasitism in insects is not restricted to bees; several lineages of wasp including most of the Chrysididae, the cuckoo wasps, are kleptoparasites. The cuckoo wasps lay their eggs in the nests of other wasps, such as those of the potters and mud daubers. Among the few exceptions, which are indeed fed by adult hosts, are cuckoo bumblebees in the subgenus Psithyrus. Their queens kill and replace the existing queen of a colony of the host species then use the host workers to feed their brood. An example of a true brood-parasitic wasp is Polistes sulcifer. This species of paper wasp has lost the ability to build their own nests, and relies on its host species, Polistes dominula, to raise its brood, with the adult hosts feeding the parasite larvae directly, unlike typical kleptoparasitic insects. In the bee species of Euglossa cordata, dominant reproductive females will display brood parasitism by replacing her daughter’s eggs with her own eggs, diverting her resources from producing grand-offspring to producing more of her own offspring. In addition, to increase her longevity and fecundity, a mother will also eat her daughter’s eggs to gain more nutrients. Host insects are sometimes tricked into bringing offspring of another species into their own nests, as is the case with the parasitic butterfly, Phengaris rebeli, and the host ant Myrmica schencki. The butterfly larvae release chemicals that confuse the host ant into believing that the P. rebeli larvae are actually ant larvae. Thus, the M. schencki ants bring back the P. rebeli larvae to their nests. - David Attenborough (1998) . The Life of Birds. New Jersey: Princeton University Press. p. 246. ISBN 0-691-01633-X. - Payne, R. B. 1997. Avian brood parasitism. In D. H. Clayton and J. Moore (eds.), Host-parasite evolution: General principles and avian models, 338–369. Oxford University Press, Oxford. - Rothstein, S.I (1990). "A model system for coevolution: avian brood parasitism". Annual Review of Ecology and Systematics 21: 481–508. doi:10.1146/annurev.ecolsys.21.1.481. - Stephen M. Yezerinac, Patrick J. Weatherhead 1997. Extra-Pair Mating, Male Plumage Coloration and Sexual Selection in yellow warblers (Dendroica petechia). Proc. R. Soc. London B. 264(1381):527–532 - Andersson, M.; Eriksson, M.O.G. (1982). "Nest parasitism in goldeneyes Bucephala clangula: some evolutionary aspects". American Naturalist 120: 1–16. doi:10.1086/283965. - Aragon, S.; Møller, A. P.; Soler, J. J.; Soler, M. (1999). "Molecular phylogeny of cuckoos supports a polyphyletic origin of brood parasitism". Journal of Evolutionary Biology 12: 495–506. doi:10.1046/j.1420-9101.1999.00052.x. - Sorenson, M.D; Payne, R.B. (2001). "A single ancient origin of brood parasitism in African finches: implications for host-parasite coevolution". Evolution 55: 2550–2567. doi:10.1554/0014-3820(2001)055[2550:asaoob]2.0.co;2. - Sorenson, M.D.; Payne, R.B. (2002). "Molecular genetic perspectives on avian brood parasitism". Integrative and Comparative Biology 42: 388–400. doi:10.1093/icb/42.2.388. PMID 21708732. - Vogl, W.; Taborsky, M.; Taborsky, B.; Teuschl, Y.; Honza, M. (2002). "Cuckoo females preferentially use specific habitats when searching for hot nests". Animal Behavior 64: 843–850. doi:10.1006/anbe.2003.1967. - Teuschl, Y; Taborsky, B; Taborsky, M (1998). "How do cuckoos find their hosts? The role of habitat imprinting". Animal Behavior 56: 1425–1433. doi:10.1006/anbe.1998.0931. - Brian Peer, Scott Robinson, and James Herkert in The Auk 117(4):892–901 - Birkhead, T. R.; Hemmings, N.; Spottiswoode, C. N.; Mikulica, O.; Moskát, C.; Ban, M.; Schulze-Hagen, K. (2011). "Internal incubation and early hatching in brood parasitic birds". Proceedings of the Royal Society Series B 278: 1019–1024. doi:10.1098/rspb.2010.1504. - Soler, M.; Soler, J. J.; Martinez, J. G.; Moller, A. P. (1995). "Magpie host manipulation by great spotted cuckoos: Evidence for an avian mafia?". Evolution 49: 770–775. doi:10.2307/2410329. - Hoover, J.P.; Robinson, S.K. (2007). "Retaliatory mafia behavior by a parasitic cowbird favors host acceptance of parasitic eggs". Proceedings of the National Academy of Sciences of the United States of America 104: 4479–4483. doi:10.1073/pnas.0609710104. - Moksnes, A; Roskaft, E (1995). "Egg-morphs and host preference in the common cuckoo (Cuculus canorus): an analysis of cuckoo and host eggs form European museums and collections". J. Zool 236: 625–648. doi:10.1111/j.1469-7998.1995.tb02736.x. - Vogl, W; Taborsky, M; Taborsky, B; Teuschl, Y; Honza, M (2002). "Cuckoo females preferentially use specific habitats when searching for hot nests". Animal Behavior 64: 843–850. doi:10.1006/anbe.2003.1967. - Roldán, M.; Soler, M. (2011). "Parental-care parasitism: How do unrelated offspring attain acceptance bv foster parents?". Behavioral Ecology 22 (4): 679–691. doi:10.1093/beheco/arr041. - Lyon, Bruce E (2003). "Egg recognition and counting reduce costs of avian conspecific brood parasitism". Nature 422: 495–499. doi:10.1038/nature01505. - Lyon, B. E. (1993). "Conspecific brood parasitism as a flexible female reproductive tactic in American coots". Animal Behaviour 46 (5): 911–928. doi:10.1006/anbe.1993.1273. - Hoi, H.; Darolová, A.; Krištofík, J. (2010). "Conspecific brood parasitism and anti-parasite strategies in relation to breeding density in female bearded tits". Behaviour 147 (12): 1533–1549. doi:10.1163/000579510X511060. - Robertson, G. J. (1998). "Egg adoption can explain joint egg-laying in common eiders". Behavioral Ecology and Sociobiology 43 (4-5): 289–296. doi:10.1007/s002650050493. - Gloag, R.; Tuero, D. T.; Fiorini, V. D.; Reboreda, J. C.; Kacelnik, A. (2012). "The economics of nestmate killing in avian brood parasites: A provisions trade-off". Behavioral Ecology 23 (1): 132–140. doi:10.1093/beheco/arr166. - Sato, Tetsu (4 September 1986). "A brood parasitic catfish of mouthbrooding cichlid fishes in Lake Tanganyika". Nature 323: 58–59. doi:10.1038/323058a0. PMID 3748180. - Baba, Reiko; Nagata, Yoshikazu; Yamagishi, Satoshi (October 1990). "Brood parasitism and egg robbing among three freshwater fish". Animal Behaviour 40: 776–778. doi:10.1016/s0003-3472(05)80707-9. - Pawelek, Jaime; Coville, Rollin. "Cuckoo Bees". UC Berkeley. Retrieved 24 February 2015. - "Cuckoo Wasps". Western Australian Museum. Retrieved 24 February 2015. - Kawakita, Atsushi; Sota, Teiji; Ito, Masao; Ascher, John S.; Tanaka, Hiroyuki; Kato, Makoto; Roubik, David W. (May 2004). "Phylogeny, historical biogeography, and character evolution in bumble bees (Bombus: Apidae) based on simultaneous analysis of three nuclear gene sequences". Molecular Phylogenetics and Evolution 31 (2): 799–804. doi:10.1016/j.ympev.2003.12.003. - Dapporto L, Cervo R, Sledge MF, Turillazzi S (2004). "Rank integration in dominance hierarchies of host colonies by the paper wasp social parasite Polistes sulcifer (Hymenoptera, Vespidae)". Journal of Insect Physiology 50 :217–223 - Ortolani, I.; Cervo, R. (2009). "Coevolution of daily activity timing in a host-parasite system". Biological Journal of the Linnean Society 96 (2): 399–405. doi:10.1111/j.1095-8312.2008.01139.x. - Akino, T; JJ Knapp; JA Thomas; GW Elmes (1999). "Chemical mimicry and host specificity in the butterfly Maculinea rebeli, a social parasite of Myrmica ant colonies". Proceedings of the Royal Society B 266 (1427): 1419–1426. doi:10.1098/rspb.1999.0796. Retrieved 28 September 2013. |Wikimedia Commons has media related to Brood parasitism.|

https://en.wikipedia.org/wiki/Brood_parasite

|This article does not cite any sources. (January 2009)| In physics, a charge carrier is a particle free to move, carrying an electric charge, especially the particles that carry electric charges in electrical conductors. Examples are electrons, ions and holes. In a conducting medium, an electric field can exert force on these free particles, causing a net motion of the particles through the medium; this is what constitutes an electric current. In different conducting media, different particles serve to carry charge: - In metals, the charge carriers are electrons. One or two of the valence electrons from each atom is able to move about freely within the crystal structure of the metal. The free electrons are referred to as conduction electrons, and the cloud of free electrons is called a Fermi gas. - In electrolytes, such as salt water, the charge carriers are ions, atoms or molecules that have gained or lost electrons so they are electrically charged. Atoms that have gained electrons so they are negatively charged are called anions, atoms that have lost electrons so they are positively charged are called cations. Cations and anions of the dissociated liquid also serve as charge carriers in melted ionic solids (see e.g. the Hall–Héroult process for an example of electrolysis of a melted ionic solid.) Proton conductors are electrolytic conductors employing positive hydrogen ions as carriers. - In a plasma, an electrically charged gas which is found in electric arcs through air, neon signs, and the sun and stars, the electrons and cations of ionized gas act as charge carriers. - In a vacuum, free electrons can act as charge carriers. These are sometimes called cathode rays. In a vacuum tube, the mobile electron cloud is generated by a heated metal cathode, by a process called thermionic emission. - In semiconductors (the material used to make electronic components like transistors and integrated circuits), in addition to electrons, the travelling vacancies in the valence-band electron population (called "holes"), act as mobile positive charges and are treated as charge carriers. Electrons and holes are the charge carriers in semiconductors. It can be seen that in some conductors, such as ionic solutions and plasmas, there are both positive and negative charge carriers, so an electric current in them consists of the two polarities of carrier moving in opposite directions. In other conductors, such as metals, there are only charge carriers of one polarity, so an electric current in them just consists of charge carriers moving in one direction. Charge carriers in semiconductors There are two recognized types of charge carriers in semiconductors. One is electrons, which carry a negative electric charge. In addition, it is convenient to treat the traveling vacancies in the valence band electron population (holes) as the second type of charge carrier, which carry a positive charge equal in magnitude to that of an electron. Carrier generation and recombination When an electron meets with a hole, they recombine and these free carriers effectively vanish. The energy released can be either thermal, heating up the semiconductor (thermal recombination, one of the sources of waste heat in semiconductors), or released as photons (optical recombination, used in LEDs and semiconductor lasers). The recombination means an electron which has been excited from the valence band to the conduction band falls back to the empty state in the valence band, known as the holes. The holes are the empty state created in the valence band when an electron gets excited after getting some energy to overpass the energy gap. Majority and minority carriers The more abundant charge carriers are called majority carriers, which are primarily responsible for current transport in a piece of semiconductor. In n-type semiconductors they are electrons, while in p-type semiconductors they are holes. The less abundant charge carriers are called minority carriers; in n-type semiconductors they are holes, while in p-type semiconductors they are electrons. In an intrinsic semiconductor, which does not contain any impurity, the concentrations of both types of carriers are ideally equal. If an intrinsic semiconductor is doped with a donor impurity then the majority carriers are electrons; if the semiconductor is doped with an acceptor impurity then the majority carriers are holes. Minority carriers play an important role in bipolar transistors and solar cells. Their role in field-effect transistors (FETs) is a bit more complex: for example, a MOSFET has both p-type and n-type regions. The transistor action involves the majority carriers of the source and drain regions, but these carriers traverse the body of the opposite type, where they are minority carriers. However, the traversing carriers hugely outnumber their opposite type in the transfer region (in fact, the opposite type carriers are removed by an applied electric field that creates an inversion layer), so conventionally the source and drain designation for the carriers is adopted, and FETs are called "majority carrier" devices. Free carrier concentration Free carrier concentration is the concentration of free carriers in a doped semiconductor. It is similar to the carrier concentration in a metal and for the purposes of calculating currents or drift velocities can be used in the same way. Free carriers are electrons (or holes) which have been introduced directly into the conduction band (or valence band) by doping and are not promoted thermally. For this reason electrons (holes) will not act as double carriers by leaving behind holes (electrons) in the other band.

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The Investiture Controversy or Investiture Contest was the most significant conflict between Church and state in medieval Europe. In the 11th and 12th centuries, a series of popes challenged the authority of European monarchies. The issue was whether the pope or the monarch would name (invest) powerful local church officials such as bishops of cities and abbots of monasteries. The conflict ended in 1122, when Emperor Henry V and Pope Calixtus II agreed on the Concordat of Worms. It differentiated between the royal and spiritual powers and gave the emperors a limited role in selecting bishops. The outcome seemed mostly a victory for the pope and his claim that he was God's chief representative in the world. However, the Emperor did retain considerable power over the Church. The investiture controversy began as a power struggle between Pope Gregory VII (1072–85) and Henry IV, Holy Roman Emperor (1056–1106). A brief but significant struggle over investiture also occurred between Henry I of England and Pope Paschal II in the years 1103 to 1107, and the issue played a minor role in the struggles between church and state in France, as well. By undercutting the Imperial power established by the Salian emperors, the controversy led to nearly 50 years of civil war in Germany, and the triumph of the great dukes and abbots. Imperial power was finally re-established under the Hohenstaufen dynasty. Historian Norman Cantor: The age of the investiture controversy may rightly be regarded as the turning-point in medieval civilization. It was the fulfillment of the early Middle Ages because in it the acceptance of the Christian religion by the Germanic peoples reached its final and decisive stage… The greater part of the religious and political system of the high Middle Ages emerged out of the events and ideas of the investiture controversy. After the decline of the Roman Empire, and prior to the Investiture Controversy, while theoretically a task of the church, investiture was in practice performed by members of the religious nobility. Many bishops and abbots were themselves usually part of the ruling nobility. Since the eldest son would inherit the title, younger siblings often found careers in the church. This was particularly true where the family may have established a proprietary church or abbey on their estate. Since Otto the Great (936-72) the bishops had been princes of the empire, had secured many privileges, and had become to a great extent feudal lords over great districts of the imperial territory. The control of these great units of economic and military power was for the king a question of primary importance, affecting as it did imperial authority. It was essential for a ruler or nobleman to appoint (or sell the office to) someone who would remain loyal. Since a substantial amount of wealth and land was usually associated with the office of a bishop or abbot, the sale of church offices (a practice known as simony) was an important source of income for leaders among the nobility, who themselves owned the land and by charity allowed the building of churches. The crisis began when a group within the church, members of the Gregorian Reform, decided to rebel against the rule of simony by forcefully taking the power of investiture from the ruling secular power, i.e. the Holy Roman Emperor and placing that power wholly within control of the church. The Gregorian reformers knew this would not be possible so long as the emperor maintained the ability to appoint the pope, so their first step was to forcibly gain the papacy from the control of the emperor. An opportunity came in 1056 when Henry IV became German king at six years of age. The reformers seized the opportunity to take the papacy by force while he was still a child and could not react. In 1059, a church council in Rome declared, with In Nomine Domini, that leaders of the nobility would have no part in the selection of popes and created the College of Cardinals as a body of electors made up entirely of church officials. Once Rome regained control of the election of the pope, it was ready to attack the practice of investiture and simony on a broad front. In 1075, Pope Gregory VII composed the Dictatus Papae. One clause asserted that the deposal of an emperor was under the sole power of the pope. It declared that the Roman church was founded by God alone – that the papal power (the auctoritas of Pope Gelasius) was the sole universal power; in particular, a council held in the Lateran Palace from 24 to 28 February the same year decreed that the pope alone could appoint or depose churchmen or move them from see to see. By this time, Henry IV was no longer a child, and he continued to appoint his own bishops. He reacted to this declaration by sending Gregory VII a letter in which he withdrew his imperial support of Gregory as pope in no uncertain terms: the letter was headed "Henry, king not through usurpation but through the holy ordination of God, to Hildebrand, at present not pope but false monk". It called for the election of a new pope. His letter ends, "I, Henry, king by the grace of God, with all of my Bishops, say to you, come down, come down!", and is often quoted with "and to be damned throughout the ages." which is a later addition. The situation was made even more dire when Henry IV installed his chaplain, Tedald, a Milanese priest, as Bishop of Milan, when another priest of Milan, Atto, had already been chosen in Rome by the pope for candidacy. In 1076 Gregory responded by excommunicating Henry, and deposed him as German king, releasing all Christians from their oath of allegiance. Enforcing these declarations was a different matter, but the advantage gradually came to be on the side of Gregory VII. German princes and the aristocracy were happy to hear of the king's deposition. They used religious reasons to continue the rebellion started at the First Battle of Langensalza in 1075, and for seizure of royal holdings. Aristocrats claimed local lordships over peasants and property, built forts, which had previously been outlawed, and built up localized fiefdoms to secure their autonomy from the empire. Thus, because of these combining factors, Henry IV had no choice but to back down, needing time to marshal his forces to fight the rebellion. In 1077, he traveled to Canossa in northern Italy to meet the pope and apologize in person. As penance for his sins, and echoing his own punishment of the Saxons after the First Battle of Langensalza, he dramatically wore a hair shirt and stood in the snow barefoot in the middle of winter in what has become known as the Walk to Canossa. Gregory lifted the excommunication, but the German aristocrats, whose rebellion became known as the Great Saxon Revolt, were not so willing to give up their opportunity. They elected a rival king, Rudolf von Rheinfeld. Three years later, Gregory declared his support for von Rheinfeld, and excommunicated Henry IV again. Henry IV then proclaimed Antipope Clement III to be pope. In 1080, Rudolf died, effectively ending the internal revolt against Henry. In 1081, Henry invaded Rome, for the first time, with the intent of forcibly removing Gregory VII and installing a more friendly pope. Gregory VII called on his allies, the Normans in southern Italy, and they rescued him from the Germans in 1085. The Normans sacked Rome in the process, and when the citizens of Rome rose up against Gregory, he was forced to flee south with the Normans. He died soon thereafter. The Investiture Controversy continued for several decades as each succeeding pope tried to diminish imperial power by stirring up revolt in Germany. These revolts were gradually successful. Henry IV was succeeded upon his death in 1106 by his son Henry V, who had rebelled against his father in favor of the papacy, and who had made his father renounce the legality of his antipopes before he died. Nevertheless, Henry V chose one more antipope, Gregory VIII. Later, he renounced some of the rights of investiture with the Concordat of Worms, abandoned Gregory, and was received back into communion and recognized as legitimate emperor as a result. English investiture controversy of 1102 to 1107 At the time of Henry IV's death, Henry I of England and the Gregorian papacy were also embroiled in a controversy over investiture, and its solution provided a model for the eventual solution of the issue in the empire. William the Conqueror had accepted a papal banner and the distant blessing of Pope Alexander II upon his invasion, but had successfully rebuffed the pope's assertion after the successful outcome, that he should come to Rome and pay homage for his fief, under the general provisions of the "Donation of Constantine". The ban on lay investiture in Dictatus Papae did not shake the loyalty of William's bishops and abbots. In the reign of Henry I, the heat of exchanges between Westminster and Rome induced Anselm, Archbishop of Canterbury, to give up mediating and retire to an abbey. Robert of Meulan, one of Henry's chief advisors, was excommunicated, but the threat of excommunicating the king remained unplayed. The papacy needed the support of English Henry while German Henry was still unbroken. A projected crusade also required English support. Henry I commissioned the Archbishop of York to collect and present all the relevant traditions of anointed kingship. "The resulting 'Anonymous of York' treaties are a delight to students of early-medieval political theory, but they in no way typify the outlook of the Anglo-Norman monarchy, which had substituted the secure foundation of administrative and legal bureaucracy for outmoded religious ideology" Concordat of London, 1107 According to René Metz, author of "What Is Canon Law?", a concordat is a convention concluded between the Holy See and the civil power of a country to define the relationship between the Catholic Church and the state in matters in which both are concerned. The concordat is one type of an international convention. Concordats began during the First Crusade's end in 1098. The Concordat of London (1107) suggested a compromise that was later taken up in the Concordat of Worms. In England, as in Germany, the king's chancery started to distinguish between the secular and ecclesiastical powers of the prelates. Employing this distinction, Henry gave up his right to invest his bishops and abbots while reserving the custom of requiring them to swear homage for the "temporalities" (the landed properties tied to the episcopate) directly from his hand, after the bishop had sworn homage and feudal vassalage in the commendation ceremony (commendatio), like any secular vassal. The system of vassalage was not divided among great local lords in England as it was in France, since the king was in control by right of the conquest. Concordat of Worms and its significance On the European mainland, after 50 years of fighting, the Concordat of Worms provided a similar, but longer lasting, compromise when signed on September 23, 1122. It eliminated lay investiture, while leaving secular leaders some room for unofficial but significant influence in the appointment process. While the monarchy was embroiled in the dispute with the Church, it declined in power and broke apart. Localized rights of lordship over peasants grew. This resulted in multiple effects: 1) increased serfdom that reduced human rights for the majority, 2) increased taxes and levies that royal coffers declined, and 3) localized rights of Justice where courts did not have to answer to royal authority. In the long term, the decline of imperial power would divide Germany until the 19th century. Similarly, in Italy, the investiture controversy weakened the emperor's authority and strengthened local separatist forces. The papacy grew stronger from the controversy. Marshalling for public opinion engaged lay people in religious affairs increasing lay piety, setting the stage for the Crusades and the great religious vitality of the 12th century. The dispute did not end with the Concordat of Worms. Future disputes between popes and Holy Roman Emperors continued until northern Italy was lost to the empire entirely. The church would Crusade against the Holy Roman Empire under Frederick II. According to Norman Cantor: The investiture controversy had shattered the early-medieval equilibrium and ended the interpenetration of ecclesia and mundus. Medieval kingship, which had been largely the creation of ecclesiastical ideals and personnel, was forced to develop new institutions and sanctions. The result during the late eleventh and early twelfth centuries, was the first instance of a secular bureaucratic state whose essential components appeared in the Anglo-Norman monarchy." - Rubenstein, Jay (2011), Armies of Heaven: The First Crusade and the Quest for Apocalypse, Basic Books, p. 18, ISBN 0-465-01929-3. - Cantor, Norman F (1958), Church. Kingship, and Lay Investiture in England: 1089-1135, Princeton University Press, pp. 8–9. - Blumenthal Investiture Controversy pp. 34–36 - Löffler, Klemens. "Conflict of Investitures." The Catholic Encyclopedia. Vol. 8. New York: Robert Appleton Company, 1910. 29 Jan. 2015 - Appleby, R. Scott. "How the pope got his political muscle." U.S. Catholic 64.9 (1999): 36. Academic Search Complete. EBSCO. Web. 5 June 2010. - Paravicini Bagliani, Agostino. "Sia fatta la mia volontà". Medioevo (143): 76. - Halsall, Paul. "Medieval Sourcebook: Henry IV: Letter to Gregory VII, 24 January 1076". Internet Medieval Source Book. 6/2/2010 <http://www.fordham.edu/halsall/source/henry4-to-g7a.html>. - Horst, Fuhrmann. Germany in the High Middle Ages c.1050-1200. Press Syndicate of the University of Cambridge. p. 64. ISBN 0 521 31980-3. - Shaff-Herzog. A Religious Encyclopedia: or Dictionary of Biblical, Historical, Doctrinal, and Practical Theology. II vols. New York, NY: Funk and Wagnalls Publishers, 1883. Pg. 911. 6/3/2010. - Halsall, Paul."Medieval Sourcebook: Gregory VII: First Deposition and Banning of Henry IV (22 February 1076)". Internet Medieval Source Book. 6/2/2010 <>. - Cantor, The Civilization of the Middle Ages, "The Entrenchment of Secular Leadership" p 286. - Metz, René (1960). "What Is Canon Law? p.137". The Twentieth Century Encyclopedia of Catholicism, Section VIII: The Organization of the Church. 80. New York: Hawthorn Books Inc. - H. Hearder, D. P. Waley, eds. A Short History of Italy: From Classical Times to the Present Day, 1963. - N. Cantor, The Civilization of the Middle Ages, "The Entrenchment of Secular Leadership", p 395. - Blumenthal, Uta-Renate (1988). The Investiture Controversy: Church and Monarchy from the Ninth to the Twelfth Century. University of Pennsylvania Press. - Cantor, Norman F. Church. Kingship, and Lay Investiture in England: 1089-1135 (Princeton University Press, 1958) - Cantor, Norman F. (1993). The Civilization of the Middle Ages. HarperCollins, PP 265–76, 284-88 - Cowdrey, H.E.J. (1998). Pope Gregory VII, 1073–1085. Oxford University Press. - Jolly, Karen Louise. (1997). Tradition & Diversity: Christianity in a World Context to 1500. ME Sharpe. - McCarthy, T. J. H. (2014). Chronicles of the Investiture Contest: Frutolf of Michelsberg and his continuators. Manchester: Manchester Medieval Sources. ISBN 9780719084706. - Metz, René. (1960). What Is Canon Law? Hawthorn Books. New York. - Morrison, Karl F., ed. The investiture controversy: issues, ideas, and results (Holt McDougal, 1971) excerpts from primary and secondary sources - Tellenbach, Gerd (1993). The Western Church from the Tenth to the Early Twelfth Century. Cambridge University Press. - Thompson, James Westfall, and Edgar Nathaniel Johnson. (1937) An introduction to medieval Europe, 300-1500 (1937) PP 380–90 - Slocum, Kenneth, ed. Sources in Medieval Culture and History (2010) pp 170–75 - "Investiture Controversy", from Encyclopædia Britannica Online. - "Canonical Investiture", from the Catholic Encyclopedia] - "Investiture", from the Columbia Encyclopedia. - "The Owl, The Cat, And The Investiture Controversy", from the Online Reference Book for Medieval Studies (ORB). - "Empire and Papacy", from the Internet Medieval Sourcebook. - Henry IV: Letter to Gregory VII, Jan 24 1076. - Gregory VII: First Deposition and Banning of Henry IV (Feb 22, 1076) - Gregory VII: Second Banning and Dethronement of Henry IV (March 7, 1080) - Gregory VII: Dictatus Papae 1090 - Ban on Lay Investitures, 1078 - The Concordat of Worms 1122 - The Canons of the First Lateran Council, 1123 - Avalon Project, Yale University: Documents relating to the War of the Investitures

https://en.wikipedia.org/wiki/Investiture_Controversy

Labeled cross section of the nasal cavities |Classification and external resources| Rhinorrhea or rhinorrhoea is a condition where the nasal cavity is filled with a significant amount of mucus fluid. The condition, commonly known as a runny nose, occurs relatively frequently. Rhinorrhea is a common symptom of allergies or certain diseases, such as the common cold or hay fever. It can be a side effect of crying, exposure to cold temperatures, cocaine abuse or withdrawal, such as from opioids like methadone. Treatment for rhinorrhea is not usually necessary, but there are a number of medical treatments and preventive techniques available. Signs and symptoms Rhinorrhea is characterized by an excess amount of mucus produced by the mucous membranes that line the nasal cavities. The membranes create mucus faster than it can be processed, causing a backup of mucus in the nasal cavities. As the cavity fills up, it blocks off the air passageway, causing difficulty breathing through the nose. Air caught in nasal cavities, namely the sinus cavities, cannot be released and the resulting pressure may cause a headache or facial pain. If the sinus passage remains blocked, there is a chance that sinusitis may result. If the mucus backs up through the Eustachian tube, it may result in ear pain or an ear infection. Excess mucus accumulating in the throat or back of the nose may cause a post-nasal drip, resulting in a sore throat or coughing. Additional symptoms include sneezing, nosebleeds, and nasal discharge. Rhinorrhea is especially common during winter months and certain low temperature seasons. Cold-induced rhinorrhea occurs due to a combination of thermodynamics and the body's natural reactions to cold weather stimuli. One of the purposes of nasal mucus is to warm inhaled air to body temperature as it enters the body. In order for this to happen, the nasal cavities must be constantly coated with liquid mucus. During cold, dry seasons, the mucus lining nasal passages tends to dry out, meaning that mucous membranes must work harder, producing more mucus to keep the cavity lined. As a result, the nasal cavity can fill up with mucus. At the same time, when air is exhaled, water vapor in breath condenses as the warm air meets the colder outside temperature near the nostrils. This causes an excess amount of water to build up inside nasal cavities. In these cases, the excess fluid usually spills out externally through the nostrils. Rhinorrhea can be a symptom of other diseases, such as the common cold or influenza. During these infections, the nasal mucous membranes produce excess mucus, filling the nasal cavities. This is to prevent infection from spreading to the lungs and respiratory tract, where it could cause far worse damage. It has also been suggested that rhinorrhea is a result of viral evolution, and may be a response that is not useful to the host, but which has evolved by the virus to maximise its own infectivity. Rhinorrhea caused by these infections usually occur on circadian rhythms. Over the course of a viral infection, sinusitis (the inflammation of the nasal tissue) may occur, causing the mucous membranes to release more mucus. Acute sinusitis consists of the nasal passages swelling during a viral infection. Chronic sinusitis occurs when one or more nasal polyps appear. This can be caused by a deviated septum as well as a viral infection. Rhinorrhea can also occur when individuals with allergies to certain substances, such as pollen, dust, latex, soy, shellfish, or animal dander, are exposed to these allergens. In people with sensitized immune systems, the inhalation of one of these substances triggers the production of the antibody immunoglobulin E (IgE), which binds to mast cells and basophils. IgE bound to mast cells are stimulated by pollen and dust, causing the release of inflammatory mediators such as histamine. In turn, this causes, among other things, inflammation and swelling of the tissue of the nasal cavities as well as increased mucus production. Particulate matter in polluted air and chemicals such as chlorine and detergents, which can normally be tolerated, can make the condition considerably worse. Rhinorrhea is also associated with shedding tears, whether from emotional events or from eye irritation. When excess tears are produced, the liquid drains through the inner corner of the eyelids, through the nasolacrimal duct, and into the nasal cavities. As more tears are shed, more liquid flows into the nasal cavities. The buildup of fluid is usually resolved via mucus expulsion through the nostrils. If caused by a head injury, rhinorrhea can be a much more serious condition. A basilar skull fracture can result in a rupture of the barrier between the sinonasal cavity and the anterior cranial fossae or the middle cranial fossae. This rupture can cause the nasal cavity to fill with cerebrospinal fluid. This condition, known as cerebrospinal fluid rhinorrhoea or CSF rhinorrhea, can lead to a number of serious complications and possibly death if not addressed properly. Rhinorrhea can occur as a symptom of opioid withdrawal accompanied by lachrymation. Other causes include cystic fibrosis, whooping cough, nasal tumors, hormonal changes, and cluster headaches. Rhinorrhea can also be the side effect of several genetic disorders, such as primary ciliary dyskinesia. In most cases treatment for rhinorrhea is not necessary since it will clear up on its own—especially if it is the symptom of an infection. For general cases, blowing one's nose can get rid of the mucus buildup. Though blowing may be a quick-fix solution, it would likely proliferate mucosal production in the sinuses, leading to frequent and higher mucus buildups in the nose. Alternatively, saline nasal sprays and vasoconstrictor nasal sprays may also be used, but may become counterproductive after several days of use, causing rhinitis medicamentosa. In recurring cases, such as those due to allergies, there are medicinal treatments available. For cases caused by histamine buildup, several types of antihistamines can be obtained relatively cheaply from drugstores. People who prefer to keep clear nasal passages, such as singers, who need a clear nasal passage to perform, may use a technique called "nasal irrigation" to prevent rhinorrhea. Nasal irrigation involves rinsing the nasal cavity regularly with salty water or store bought saline solutions. - "Palatal necrosis due to cocaine abuse". US National Library of Medicine. Retrieved 2012-09-21. - Eileen Trigoboff; Kneisl, Carol Ren; Wilson, Holly Skodol (2004). Contemporary psychiatric-mental health nursing. Upper Saddle River, N.J: Pearson/Prentice Hall. p. 274. ISBN 0-13-041582-0. - "Rhinorrhea". Online Etymology Dictionary. Retrieved 2011-09-21. - "Nasal discharge". Medline Plus. United States National Library of Medicine, National Institutes of Health. Retrieved 2007-11-01. - "Rhinorrhea Overview". FreeMd. Retrieved 2011-09-21. - "Why Does Cold Weather Cause Runny Noses?". NPR. Retrieved 2011-09-22. - "Why Does My Nose Run?". Kids Health. Retrieved 2011-09-22. - Smolensky MH, Reinberg A, Labrecque G (May 1995). "Twenty-four Hour Pattern in Symptom Intensity of Ciral and Allergic Rhinitis: Treatment Implications". The Journal of Allergy and Clinical Immunology 95 (5 Pt 2): 1084–96. doi:10.1016/s0091-6749(95)70212-1. PMID 7751526. - "Rhinorrhea – Definition, Symptoms, Causes, Diagnosis and Treatment". Prime Health Channel. 2011-08-30. Retrieved 2011-09-24. - Dipiro, J.T.; Talbert, R.L.; Yee, G.C. (2008). Pharmacotherapy: A Pathophysiologic Approach (7th ed.). New York, NY: The McGraw-Hill Companies, Inc. pp. 1565–1575. ISBN 978-0-07-147899-1. - Welch; et al. (2011-07-22). "CSF Rhinorrhea". Medscape. Retrieved 2011-09-22. - "Opioid Withdrawal Protocol" (PDF). Mental Health and Addiction Services. Retrieved 2011-09-24. - Aubrey, Allison (2007-02-22). "Got a Runny Nose? Flush it Out!". NPR. Retrieved 2011-09-21. - Runny Nose: A Guide for Parents at the Wayback Machine (archived February 25, 2012) from the Pennsylvania Medical Society - Cold and flu advice (NHS Direct) - How to Wipe Your Nose - How to Wipe Your Nose on Your Hands

https://en.wikipedia.org/wiki/Rhinorrhea

Drama Historical Context Teacher Resources Find Drama Historical Context educational ideas and activities Showing 1 - 20 of 115 resources The Diary of Anne Frank While designed to supplement a viewing of the PBS Masterpiece Classic The Diary of Anne Frank, this resource can also serve as an excellent informational text and activity source for your students on the historical context and timeline... 8th - 10th Social Studies & History CCSS: Adaptable Historical Context: African-American Oral Tradition For this African-American oral tradition worksheet, students read and learn about the vast and important history of the oral traditions that existed in the African-American culture. Students use this worksheet as a pre-reading text to... 10th - 11th English Language Arts To Kill a Mockingbird: Culture and History The second of 10 lessons in a unit study of To Kill a Mockingbird establishes the historical and cultural context of Harper Lee's novel. The class listens to second part of an audio that describes Lee's life experiences that parallel the... 9th - 12th English Language Arts CCSS: Adaptable Investigating the Harlem Renaissance The work of Langston Hughes opens the door to research into the origin and legacy of the Harlem Renaissance and how the literature of the period can be viewed as a commentary on race relations in America. In addition, groups are assigned... 11th - 12th English Language Arts Activities for Teaching “The Road Not Taken” by Robert Frost Use all of these exercises, assignments, and assessments or pick and choose your favorites for your study of "The Road Not Taken" by Robert Frost. In this resource you will find: an informational text to examine, vocabulary lists and... 7th - 9th English Language Arts Comprehension and Discussion Activities for the Movie Rabbit-Proof Fence Lead discussion and thoughtful analysis as pupils view Rabbit Proof Fence, a drama based on true story about three aboriginal girls who ran away from Western Australia to return to their Aboriginal families in 1931. Here you'll find... 6th - 8th Social Studies & History CCSS: Adaptable Creating Scrolls Based on the Illustrated Tale of Genji Now these are learning activities full of fun, art, and cultural exploration. Kids consider the art of storytelling through comic book images. They then look at the Tale of Genji as it was written in the 11th century. They discuss... 6th - 8th Visual & Performing Arts Comedy Across the Curriculum The New York Times Learning Network provides the resources that permit pupils to examine and then write and perform a fake news broadcast in the vein of “The Daily Show” or “Saturday Night Live” Weekend Update. The generated reports... 11th - 12th Visual & Performing Arts Hamlet and the Elizabethan Revenge Ethic in Text and Film Students research the social context of Elizabethan England for Shakespeare's "Hamlet". They identify cultural influences on the play focusing on the theme of revenge and then analyze and compare film interpretations of the play. 9th - 12th English Language Arts English Literature: An Overview Relate literary works and authors to the major themes of English literature from the Anglo-Saxon period through the 20th century. Working in groups, high schoolers will evaluate period philosophy, religion, and politics that influenced... 11th - 12th English Language Arts Figures of Speech: A Midsummer Night's Dream High school readers analyze figures of speech in Shakespeare's A Midsummer Night's Dream with support from a two-page worksheet. They respond to four multi-step questions regarding the use of metaphors, similes, hyperbole, and irony in... 8th - 12th English Language Arts The Play's the Thing: The Drama of Cyrano de Bergerac Students practice dramatic 'living' through various drama activities. In this drama lesson, students define drama, view examples of dramatic elements in Cyrano de Bergerac and Roxanne, define characterization within the dramas, study the... 8th English Language Arts To Kill a Mockingbird: Historical and Cultural Context As part of their study of the film adaptation of To Kill A Mockingbird, class members analyze how Robert Mulligan uses the film lens to depict the historical period and social issues presented in Harper Lee's novel. A superior resource... 7th - 10th Visual & Performing Arts CCSS: Adaptable Common Core Teaching and Learning Strategies New ReviewHere's a resource that deserves a place in your curriculum library, whether or not your school has adopted the Common Core. Designed for middle and high school language arts classes, the packet is packed with teaching tips, materials,... 6th - 12th English Language Arts CCSS: Designed Teaching the Holocaust through Literature Centered on the short story "The Tenth Man" by Polish Holocaust survivor Ida Fink, here is a solid one-day resource to support study of World War II or Nazi history, short stories, or to complement any ELA unit on The Diary of Anne Frank... 7th - 12th English Language Arts The Trial of Hamlet Hamlet, that is not a rat behind the curtain, it is Polonius, and now you’re on trial for his murder. Practice and develop close reading skills, discover how a trial works, and get the entire class involved in this trial. The class... 11th - 12th English Language Arts CCSS: Designed "Et tu, Brute?" - The Characters, Conflict and Historical context of Shakespeare's Julius Caesar Students analyze the Shakespearian play, "Julius Caesar" in this seven lesson unit. Through readings, hands-on projects, and the study of plot development, comparisons are made to the movie and the historical records available. 6th English Language Arts

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Real versus nominal value (economics) In economics, a nominal value is an economic value expressed in historical nominal monetary terms. By contrast, a real value is a value that has been adjusted from a nominal value to remove the effects of general price level changes over time and is thus measured in terms of the general price level in some reference year (the base year). For example, changes in the nominal value of some commodity bundle over time can happen because of a change in the quantities in the bundle or their associated prices, whereas changes in real values reflect only changes in quantities. The process of converting from nominal to real terms is known as inflation adjustment. Real values are a measure of purchasing power net of any price changes over time. For example, nominal income is often restated as real income, thus removing that part of income changes that merely reflect inflation (a general increase in prices). Similarly, for aggregate measures of output, such as gross domestic product (GDP), the nominal amount reflects production quantities and prices in that time period, whereas the differences between real amounts in different time periods reflect only changes in quantities. A series of real values over time, such as for real GDP, measures quantities over time expressed in prices of one year, called the base year (or more generally the base period). Real values in different years then express values of the bundles as if prices had been constant for all the years, with any differences due to differences in underlying quantities. The nominal/real value distinction can apply not only to time-series data, as above, but also to cross-section data varying by region. For example, the total sales value of a particular good produced in a particular region of a country is influenced by both the physical amount sold and the selling price, which may be different from that of the country as a whole; for purposes of comparing the economic activity of different regions, the nominal output of the good in that region can be adjusted into real terms by repricing the goods at national-average prices. The nominal value of a commodity bundle in a given year depends on both quantities and then-current prices, namely, as a sum of prices times quantities for the different commodities in the bundle. In turn nominal values are related to real values by the following arithmetic definition: - nominal value / real value = (P x Q) / Q = P. Here P is a price index, and Q is a quantity index of real value. In the equation, P is constructed to equal 1.00 in the base year. Alternatively, P can be constructed to equal 100 in the base year: - (nominal value / real value) x 100 = P. The base year can be any year, and comparisons of quantities are valid provided all values are adjusted to their values in the same base year. After a number of years have passed in which government statistics have been reported in terms of a particular base year, a new base year for comparisons is typically adopted; for the next several years all new data as well as all pre-existing data will be reported in terms of the new base year. The simple case of a bundle of commodities (goods) is one that has only one commodity. In that case, output or consumption may be measured either in terms of money value (nominal) or physical quantity (real). Let i designate that commodity and let: - Pi = the unit price of i, say, $5 - Qi = the quantity of good i, say, 10 units. The nominal value of the good would then be price times quantity: - nominal value of good i = Pi x Qi = $5 x 10 = $50. Given only the nominal value and price, derivation of a real value is immediate: - real value of good i = (Pi x Qi)/Pi = Qi = 50/5 = 10. The price "deflates" (divides) the nominal value to derive a real value, the quantity itself. Similar for a series of years, say five, given only nominal values of the good and prices in each year t, a real value can be derived for each of the five years: - real value in year t = (nominal value in year t) / (price relative to base year) = Qit. The following example generalizes from an individual good to a bundle of goods across different years for which P, a price index comparing the general price level across years, is available. Consider a nominal value (say of an hourly wage rate) in each different year t. To derive a real-value series from a series of nominal values in different years, one divides the nominal wage rate in each year by Pt, the price index in that year. By definition then: - real value in year t = (nominal value in year t) / Pt. If for years 1 and 2 (say 20 years apart) the nominal wage and price level P of goods are respectively then real wages using year 1 as the base year are respectively: The real wage so constructed in each different year indexes the amount of commodities in that year that could be purchased, for comparison to other years. Thus, in the example the price level increased by 33 percent, but the real wage rate still increased by 20 percent, permitting a 20 percent increase in the quantity of commodities the nominal wage could purchase. The above generalization to a commodity bundle from the previous sing-good illustration has practical use, because price indexes and the National Income and Product Accounts are constructed from such bundles of commodities and their respective prices. A sum of nominal values for each of the different commodities in the bundle is also called a nominal value. A bundle of n different commodities with corresponding prices and quantities for each year t defines: - nominal value of that bundle in year t = P1t x Q1t + . . . + Pnt x Qnt. From the above: - Pt = the value of a price index in year t. The nominal value of the bundle over a series of years and corresponding Pt define: - real value of the bundle in year t = Qt = (nominal value of the bundle in year t) / Pt. Alternatively, multiplying both sides by Pt: - nominal value of the bundle in year t = Pt x Qt. So, every nominal value can be dichotomized into a price-level part and a real part. The real part Qt is an index of the quantities in the bundle. Real values (such as real wages or real gross domestic product) can be derived by dividing the relevant nominal value (e.g., nominal wage rate or nominal GDP) by the appropriate price index. For consumers, a relevant bundle of goods is that used to compute the Consumer Price Index (CPI). So, for wage earners as consumers a relevant real wage is the nominal wage (after-tax) divided by the CPI. A relevant divisor of nominal GDP is the GDP price index. Real values represent the purchasing power of nominal values in a given period, such as wages or total production. In particular, price indexes are typically calculated relative to some base year. If for example the base year is 1992, real values are expressed in constant 1992 dollars, referenced as 1992=100, since the published index is usually normalized to have the price index equal 100 in the base year. To use the price index as a divisor for converting a nominal value into a real value, as in the previous section, the published index is first divided by the base-year price-index value of 100. In the U.S. National Income and Product Accounts, nominal GDP is called GDP in current dollars (that is, in prices current for each designated year), and real GDP is called GDP in [base-year] dollars (that is, in dollars that can purchase the same quantity of commodities as in the base year). In effect the price index of 100 for the base year is a numéraire for price-index values in other years. The terminology of classical economics used by Adam Smith used a unit of labour as the purchasing power unit, so monetary quantities were defined by the cost of an hours of labour required to produce or purchase a given quantity. Since interest rates are measured as percentages rather than in terms of units of some currency, real interest rates are measured as the difference between nominal interest rates and the rate of inflation. The expected real interest rate as of the starting time of a loan is the nominal interest rate minus the inflation rate expected over the term of the loan. The realized (ex post) real interest rate is computed by subtracting the actual inflation rate that ends up prevailing during the life of the loan from the nominal interest rate, and reflects what actually happened during the life of the loan. The relationship above is approximate only. The actual relationship is as follows: - IRN is the nominal interest rate, - IRR is the real interest rate, and - I is the inflation rate. - Aggregation problem - Classical dichotomy - Constant Item Purchasing Power Accounting - Cost-of-living index - Financial repression - Index (economics) - Inflation accounting - Money illusion - National accounts - Neutrality of money - Peppercorn (legal), a nominal fee paid to fulfill a contractual requirement - Real interest rate - Real prices and ideal prices - W.E. Diewert, "index numbers," ( 2008)The New Palgrave Dictionary of Economics, 2nd ed. Abstract. - R. O'Donnell (1987). "real and nominal quantities," The New Palgrave: A Dictionary of Economics, v. 4, pp. 97–98 (Adam Smith's early distinction vindicated) - Amartya Sen (1979). "The Welfare Basis of Real Income Comparisons: A Survey," Journal of Economic Literature, 17(1), p p. 1-45. - D. Usher (1987). "real income," The New Palgrave: A Dictionary of Economics, v. 4, pp. 104–05 - DataBasics: Deflating Nominal Values to Real Values from Federal Reserve Bank of Dallas - CPI Inflation Calculator from U.S. Bureau of Labor Statistics

https://en.wikipedia.org/wiki/Real_versus_nominal_value_(economics)

Communication Skills Teacher Resources Find Communication Skills educational ideas and activities Showing 1 - 20 of 15,512 resources Effective Communication for Successful Careers Having good written communication skills is a must in today's workplace. Foster these skills by engaging learners is a discussion on how good writing skills can improve communication in the workplace. Have them write a project proposal... 8th 21st Century Skills CCSS: Designed Communication, Day 1: Non-Verbal Communication Have your secondary special education class learn and practice effective communication skills. Both verbal and non-verbal communication is discussed and practiced. They communicate using body language, build listening skills, and discuss... 9th - 12th Health 101 Ways to Teach Children Social Skills Increasing pressure to improve student achievement has made it easy to overlook the social skills they also need to develop. With this collection of worksheets and activities, you'll be able to improve children's communication, teamwork,... K - 6th 21st Century Skills CCSS: Adaptable Conversation Visual Prompts Help learners understand the importance and proper placement of non-verbal communication using these visual prompts. This set includes graphics to remind students to listen (ear), keep personal space (ruler), use the right facial... K - 12th Special Education & Programs CCSS: Adaptable Shakespeare: Nonverbal and Verbal Communication Define nonverbal communication and view "The Shakespeare Sessions" for examples of nonverbal communication. Groups read through the dialogue of a scene and assign appropriate gestures, movements, and mannerisms to events and characters.... 7th - 12th English Language Arts Lesson: Communication, What's Valued, and the Written Word Upper graders compare their cell phones to a lacquer box from the Japanese Edo Period. They consider how each is a form of communication and how the very nature of each object communicates social norms, ideology, and beliefs. A really... 9th - 12th Visual & Performing Arts Ag Communications - One to One Communication Explore the many aspects of communication and conversation. There are definitions and examples given to identify, explain, and understand the terminology of non-verbal communication. Helping the class become aware of the skills involved... 7th - 12th Health Verbal Versus Nonverbal Communication Young scholars create a multimedia presentation. They will complete a verbal versus non-verbal communication chart to create a multimedia presentation which will include the different types of communication strategies. Then answer a... 6th - 9th English Language Arts BBC Learning English, Listening Comprehension In this specific listening comprehension worksheet, students listen to an audio file and then choose the best answer to 15 corresponding multiple choice questions. Students then respond to four questions about non-verbal communication in... Higher Ed English Language Arts

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Asperger syndrome now comes under the single umbrella term of autism spectrum disorder (ASD). It is classified as a developmental disorder that affects how the brain processes information. People with Asperger syndrome have a wide range of strengths, weaknesses, skills and difficulties. Common characteristics include difficulty in forming friendships, communication difficulties (such as a tendency to take things literally), and an inability to understand social rules and body language. Although Asperger syndrome cannot be cured, appropriate intervention and experience can help people to develop skills, use strategies to compensate and help build up coping skills. Social skills training, which teaches people how to behave in different social situations, is often considered to be of great value to those with Asperger syndrome. Counselling or psychological therapy, including cognitive behaviour therapy (CBT) can help people with Asperger syndrome to understand and manage their behavioural responses. New ASD classification system A new classification system for autism and Asperger syndrome, introduced in 2013 (in the fifth edition of the Diagnostic and Statistical Manual of Mental Disorders ), gives only one diagnosis of autism spectrum disorder. This is the result of much research that indicated there was not enough evidence to suggest that the conditions of autism and Asperger syndrome were distinct conditions, so now they all come under the single umbrella term, ASD. This means that a diagnosis of Asperger syndrome will no longer be given. The preferred term is now ASD, However, there are a number of people who have been diagnosed with Asperger’s in the past, and identify with this diagnosis. They will still be able to refer to their condition as having Asperger’s into the future, despite the fact that it is no longer a formal diagnosis. Symptoms of Asperger syndrome More males than females are diagnosed with Asperger syndrome or ASD. While every person who has the condition will experience different symptoms and severity of symptoms, some of the more common characteristics include: - average or above-average intelligence - difficulties with high-level language skills such as verbal reasoning, problem solving, making inferences and predictions - difficulties in empathising with others - problems with understanding another person’s point of view - difficulties engaging in social routines such as conversations and ‘small talk’ - problems with controlling feelings such as anger, depression and anxiety - a preference for routines and schedules which can result in stress or anxiety if a routine is disrupted - specialised fields of interest or hobbies. Emotions of other people A person with Asperger syndrome may have trouble understanding the emotions of other people, and the subtle messages sent by facial expression, eye contact and body language are often missed or misinterpreted. Because of this, people with Asperger syndrome might be mistakenly perceived as being egotistical, selfish or uncaring. These are unfair labels because the person concerned may be unable to understand other people’s emotional states. People with Asperger syndrome are usually surprised when told their actions were hurtful or inappropriate. Asperger syndrome and sexual codes of conduct Research into the sexual understanding of people with Asperger syndrome is in its infancy. Studies suggest that individuals with Asperger syndrome are as interested in sex as anyone else, but many struggle with the wide range of complex skills required to successfully have intimate relationships. People with Asperger syndrome can sometimes appear to have an ‘inappropriate’, ‘immature’ or ‘delayed’ understanding of sexual codes of conduct. They may not understand the boundaries of appropriate sexual behaviour and expression. This can sometimes result in sexually inappropriate behaviour. For example, an adult with Asperger syndrome might not understand the social rule that it is not considered socially appropriate to display sexualised behaviours in a public place. Even people who are high achieving and academically or vocationally successful can have trouble negotiating the ‘hidden rules’ of courtship. Issues for partners of people with Asperger syndrome or ASD Some people with Asperger syndrome can successfully maintain relationships and parent children. However, like most relationships, there are challenges. A common marital problem is unfair distribution of responsibilities. For example, the partner of a person with Asperger syndrome may be used to doing everything in the relationship when it is just the two of them. However, the partner may need practical and emotional support once children come along, something that the person with Asperger syndrome may not be fully able to provide. When the partner expresses frustration or becomes upset that they are given no help of any kind, the person with Asperger syndrome is typically baffled. Tension in the relationship often makes their symptoms worse. An adult’s diagnosis of Asperger syndrome often follows their child’s diagnosis of ASD. This ‘double whammy’ can be extremely distressing to the partner who has to cope simultaneously with both diagnoses. Counselling, or joining a support group where they can talk with other people who face the same challenges, can be helpful. Some common issues for partners of people with Asperger syndrome include: - feeling overly responsible for their partner - failure to have their own needs met by the relationship - lack of emotional support from family members and friends who do not fully understand or appreciate the extra strains placed on a relationship by Asperger syndrome - a sense of isolation, because the challenges of their relationship are unique and not easily understood by others - frustrations, since problems in the relationship do not seem to improve despite great efforts - doubting the integrity of the relationship, or frequently wondering about whether or not to end the relationship - difficulties in accepting that their partner will not ‘recover’ from Asperger syndrome - after accepting that their partner’s Asperger syndrome cannot be ‘cured’, partners can often experience emotions such as guilt, despair and disappointment. The workplace and Asperger syndrome A person with Asperger syndrome may find their job opportunities are limited by their disability. It may help to choose a vocation that takes into account their symptoms and capitalises on their strengths, rather than highlighting their weaknesses. Career suggestions for visual thinkers The following career suggestions are adapted from material written by Temple Grandin, who has high-functioning autism and is a professor at Colorado University, USA. Suggestions include: - computer programming - commercial art - equipment design - appliance repair - handcraft artisan - webpage designer - video game designer - building maintenance - building trades. Career suggestions for those good at mathematics or music - computer programming - journalist, copy editor - taxi driver - piano (or other musical instrument) tuner - filing positions - bank teller Where to get help - Your doctor - Aspergers Victoria Tel. (03) 9845 2766 - Amaze – Autism Victoria Tel. (03) 9657 1600 - Centre for Developmental Disability Health Victoria (CDDHV) Tel. (03) 9902 4467 Things to remember - A person with Asperger syndrome often experiences difficulties when trying to understand the emotions of other people. Subtle messages that are sent by facial expression, eye contact and body language are often missed. - Social skills training, which teaches people with Asperger syndrome how to behave in different social situations, is often considered to be of great value to people with this syndrome. This page has been produced in consultation with and approved by: Autism Victoria trading as amaze Page content currently being reviewed. Content on this website is provided for education and information purposes only. Information about a therapy, service, product or treatment does not imply endorsement and is not intended to replace advice from your doctor or other registered health professional. Content has been prepared for Victorian residents and wider Australian audiences, and was accurate at the time of publication. Readers should note that, over time, currency and completeness of the information may change. All users are urged to always seek advice from a registered health care professional for diagnosis and answers to their medical questions.

https://www.betterhealth.vic.gov.au/health/conditionsandtreatments/asperger-syndrome-and-adults

Scientists today measure the Earth's surface temperature using thermometers at weather stations and on ships and buoys all over the world. Such thermometer records cover a large fraction of the globe going back to the mid-19th century, allowing scientists to determine a global average temperature trend for the last 160 years. Before that time not many thermometer records are available, so scientists use indirect temperature measurements, supported by anecdotal evidence recorded by diarists, and the few thermometer records that do exist. Scientists must rely solely on indirect methods to look back further than recorded human history. Indirect ways of assessing past temperatures, using so-called temperature proxies, take measurements of responses to past temperature change that are preserved in natural archives such as ice, rocks and fossils. For example, ice sheets form as snow builds up, with each year's snowfall preserved as a single, visible layer. There are measurable chemical differences in snow formed at different temperatures, so ice cores provide a record of polar temperature going back around 250,000 years for Greenland and 800,000 years for Antarctica. Yearly banding is also found in fossilised corals and lake sediment deposits, and each band has a specific chemistry that reflects the temperature when it formed. Growth rings in tree trunks can be wider or thinner depending on the climate at the time of growth, so fossilised trees can reveal the length of growing seasons. And fossilised or frozen pollen grains allow scientists to determine what plants were growing in the past, which can give us a good idea of the climate at the time. Marine sediment cores provide temperature records spanning millions of years. They contain the fossilised shells of tiny marine creatures that preserve a chemical record of the sea temperature when they lived. To make their temperature reconstructions as accurate as possible scientists have calibrated each proxy by testing how it changes in response to changing temperature. However, the further back in time we look, the more sparse the proxy temperature records become. Therefore the most reliable way to work out past temperatures is to combine different proxies – and to use data from many locations to screen out local temperature fluctuations. • This article was written by Carbon Brief in conjunction with the Guardian and partners The ultimate climate change FAQ This editorial is free to reproduce under Creative Commons This post by The Guardian is licensed under a Creative Commons Attribution-No Derivative Works 2.0 UK: England & Wales License. Based on a work at theguardian.com

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The overconfidence effect is a well-established bias in which a person's subjective confidence in his or her judgments is reliably greater than the objective accuracy of those judgments, especially when confidence is relatively high. For example, Overconfidence is one example of a miscalibration of subjective probabilities. Throughout the research literature, overconfidence has been defined in three distinct ways: (1) overestimation of one's actual performance, (2) overplacement of one's performance relative to others, and (3) the excessive certainty regarding the accuracy of one's beliefs − called overprecision. The most common way in which overconfidence has been studied is by asking people how confident they are of specific beliefs they hold or answers they provide. The data show that confidence systematically exceeds accuracy, implying people are more sure that they are correct than they deserve to be. If human confidence had perfect calibration, judgments with 100% confidence would be correct 100% of the time, 90% confidence correct 90% of the time, and so on for the other levels of confidence. By contrast, the key finding is that confidence exceeds accuracy so long as the subject is answering hard questions about an unfamiliar topic. For example, in a spelling task, subjects were correct about 80% of the time, whereas they claimed to be 100% certain. Put another way, the error rate was 20% when subjects expected it to be 0%. In a series where subjects made true-or-false responses to general knowledge statements, they were overconfident at all levels. When they were 100% certain of their answer to a question, they were wrong 20% of the time. - 1 Overconfidence distinctions - 2 Practical implications - 3 Individual differences - 4 See also - 5 References - 6 Further reading One manifestation of the overconfidence effect is the tendency to overestimate one's standing on a dimension of judgement or performance. This subsection of overconfidence focuses on the certainty one feels in their own ability, performance, level of control or chance of success. This phenomenon is most likely to occur on hard tasks, hard items, when failure is likely or when the individual making the estimate is not especially skilled. Overestimation has been seen to occur across domains other than those pertaining to one's own performance. This includes the illusion of control, planning fallacy. Illusion of control Illusion of control describes the tendency for people to behave as if they might have some control when in fact they have none. However, evidence does not support the notion that people systematically overestimate how much control they have; when they have a great deal of control, people tend to underestimate how much control they have. The planning fallacy describes the tendency for people to overestimate their rate of work or to underestimate how long it will take them to get things done. It is strongest for long and complicated tasks, and disappears or reverses for simple tasks that are quick to complete. Wishful-thinking effects, in which people overestimate the likelihood of an event because of its desirability, are relatively rare. This may be in part because people engage in more defensive pessimism in advance of important outcomes, in an attempt to reduce the disappointment that follows overly optimistic predictions. Overprecision is the excessive confidence that one knows the truth. For reviews, see Harvey (1997) or Hoffrage (2004). Much of the evidence for overprecision comes from studies in which participants are asked about their confidence that individual items are correct. This paradigm, while useful, cannot distinguish overestimation from overprecision; they are one and the same in these item-confidence judgments. After making a series of item-confidence judgments, if people try to estimate the number of items they got right, they do not tend to systematically overestimate their scores. The average of their item-confidence judgments exceeds the count of items they claim to have gotten right. One possible explanation for this is that item-confidence judgments were inflated by overprecision, and that their judgments do not demonstrate systematic overestimation. The strongest evidence of overprecision comes from studies in which participants are asked to indicate how precise their knowledge is by specifying a 90% confidence interval around estimates of specific quantities. If people were perfectly calibrated, their 90% confidence intervals would include the correct answer 90% of the time. In fact, hit rates are often as low as 50%, suggesting people have drawn their confidence intervals too narrowly, implying that they think their knowledge is more accurate than it actually is. Overplacement is perhaps the most prominent manifestation of the overconfidence effect. Overplacement is a judgment of your performance compared to another. This subsection of overconfidence occurs when people believe themselves to be better than others, or "better-than-average". It is the act of placing yourself or rating yourself above others (superior to others). Overplacement more often occurs on simple tasks, ones we believe are easy to accomplish successfully. One explanation for this theory is its ability to self-enhance. Perhaps the most celebrated better-than-average finding is Svenson’s (1981) finding that 93% of American drivers rate themselves as better than the median. The frequency with which school systems claim their students outperform national averages has been dubbed the “Lake Wobegon” effect, after Garrison Keillor’s apocryphal town in which “all the children are above average.” Overplacement has likewise been documented in a wide variety of other circumstances. Kruger (1999), however, showed that this effect is limited to “easy” tasks in which success is common or in which people feel competent. For difficult tasks, the effect reverses itself and people believe they are worse than others. Some researchers have claimed that people think good things are more likely to happen to them than to others, whereas bad events were less likely to happen to them than to others. But others have pointed out that prior work tended to examine good outcomes that happened to be common (such as owning one’s own home) and bad outcomes that happened to be rare (such as being struck by lightning). Event frequency accounts for a proportion of prior findings of comparative optimism. People think common events (such as living past 70) are more likely to happen to them than to others, and rare events (such as living past 100) are less likely to happen to them than to others. Taylor and Brown (1988) have argued that people cling to overly positive beliefs about themselves, illusions of control, and beliefs in false superiority, because it helps them cope and thrive. Although there is some evidence that optimistic beliefs are correlated with better life outcomes, most of the research documenting such links is vulnerable to the alternative explanation that their forecasts are accurate. The cancer patients who are most optimistic about their survival chances are optimistic because they have good reason to be. Overconfidence has been called the most “pervasive and potentially catastrophic” of all the cognitive biases to which human beings fall victim. It has been blamed for lawsuits, strikes, wars, and stock market bubbles and crashes. Strikes, lawsuits, and wars could arise from overplacement. If plaintiffs and defendants were prone to believe that they were more deserving, fair, and righteous than their legal opponents, that could help account for the persistence of inefficient enduring legal disputes. If corporations and unions were prone to believe that they were stronger and more justified than the other side, that could contribute to their willingness to endure labor strikes. If nations were prone to believe that their militaries were stronger than were those of other nations, that could explain their willingness to go to war. Overprecision could have important implications for investing behavior and stock market trading. Because Bayesians cannot agree to disagree, classical finance theory has trouble explaining why, if stock market traders are fully rational Bayesians, there is so much trading in the stock market. Overprecision might be one answer. If market actors are too sure their estimates of an asset’s value is correct, they will be too willing to trade with others who have different information than they do. Oskamp (1965) tested groups of clinical psychologists and psychology students on a multiple-choice task in which they drew conclusions from a case study. Along with their answers, subjects gave a confidence rating in the form of a percentage likelihood of being correct. This allowed confidence to be compared against accuracy. As the subjects were given more information about the case study, their confidence increased from 33% to 53%. However their accuracy did not significantly improve, staying under 30%. Hence this experiment demonstrated overconfidence which increased as the subjects had more information to base their judgment on. Even if there is no general tendency toward overconfidence, social dynamics and adverse selection could conceivably promote it. For instance, those most likely to have the courage to start a new business are those who most overplace their abilities relative to those of other potential entrants. And if voters find confident leaders more credible, then contenders for leadership learn that they should express more confidence than their opponents in order to win election. Overconfidence can be beneficial to individual self-esteem as well as giving an individual the will to succeed in their desired goal. Just believing in oneself may give one the will to take one's endeavours further than those who do not. Very high levels of core self-evaluations, a stable personality trait composed of locus of control, neuroticism, self-efficacy, and self-esteem, may lead to the overconfidence effect. People who have high core self-evaluations will think positively of themselves and be confident in their own abilities, although extremely high levels of core self-evaluations may cause an individual to be more confident than is warranted. - Pallier, Gerry; Wilkinson, Rebecca; Danthiir, Vanessa; Kleitman, Sabina; Knezevic, Goran; Stankov, Lazar; Roberts, Richard D. (2002). "The Role of Individual Differences in the Accuracy of Confidence Judgments". The Journal of General Psychology 129 (3): 257–299. doi:10.1080/00221300209602099. - Moore, Don A.; Healy, Paul J. (2008). "The trouble with overconfidence.". Psychological Review 115 (2): 502–517. doi:10.1037/0033-295X.115.2.502. - Adams, P. A.; Adams, J. K. (1960). "Confidence in the recognition and reproduction of words difficult to spell". The American journal of psychology 73 (4): 544–552. doi:10.2307/1419942. PMID 13681411. - Lichtenstein, Sarah; Fischhoff, Baruch; Phillips, Lawrence D. (1982). "Calibration of probabilities: The state of the art to 1980". In Kahneman, Daniel; Slovic, Paul; Tversky, Amos. Judgment Under Uncertainty: Heuristics and Biases. Cambridge University Press. pp. 306–334. ISBN 978-0-521-28414-1. - Langer, Ellen J. (1975). "The illusion of control". Journal of Personality and Social Psychology 32 (2): 311–328. doi:10.1037/0022-35220.127.116.111. - Buehler, Roger; Griffin, Dale; Ross, Michael (1994). "Exploring the "planning fallacy": Why people underestimate their task completion times". Journal of Personality and Social Psychology 67 (3): 366–381. doi:10.1037/0022-3518.104.22.1686. - Krizan, Zlatan; Windschitl, Paul D. (2007). "The influence of outcome desirability on optimism" (PDF). Psychological Bulletin 133 (1): 95–121. doi:10.1037/0033-2909.133.1.95. PMID 17201572. - Norem, Julie K.; Cantor, Nancy (1986). "Defensive pessimism: Harnessing anxiety as motivation". Journal of Personality and Social Psychology 51 (6): 1208–1217. doi:10.1037/0022-3522.214.171.1248. - McGraw, A. Peter; Mellers, Barbara A.; Ritov, Ilana (2004). "The affective costs of overconfidence" (PDF). Journal of Behavioral Decision Making 17 (4): 281–295. doi:10.1002/bdm.472. - Harvey, Nigel (1997). "Confidence in judgment". Trends in Cognitive Sciences 1 (2): 78–82. doi:10.1016/S1364-6613(97)01014-0. - Hoffrage, Ulrich (2004). "Overconfidence". In Pohl, Rüdiger. Cognitive Illusions: a handbook on fallacies and biases in thinking, judgement and memory. Psychology Press. ISBN 978-1-84169-351-4. - Gigerenzer, Gerd (1993). "The bounded rationality of probabilistic mental models". In Manktelow, K. I.; Over, D. E. Rationality: Psychological and philosophical perspectives. London: Routledge. pp. 127–171. ISBN 9780415069557. - Alpert, Marc; Raiffa, Howard (1982). "A progress report on the training of probability assessors". In Kahneman, Daniel; Slovic, Paul; Tversky, Amos. Judgment Under Uncertainty: Heuristics and Biases. Cambridge University Press. pp. 294–305. ISBN 978-0-521-28414-1. - Svenson, Ola (1981). "Are we all less risky and more skillful than our fellow drivers?". Acta Psychologica 47 (2): 143–148. doi:10.1016/0001-6918(81)90005-6. - Cannell, John Jacob (1989). "How public educators cheat on standardized achievement tests: The "Lake Wobegon" report". Friends for Education (Albuquerque, NM). - Dunning, David (2005). Self-Insight: Roadblocks and Detours on the Path to Knowing Thyself. Psychology Press. ISBN 978-1841690742. - Kruger, Justin (1999). "Lake Wobegon be gone! The "below-average effect" and the egocentric nature of comparative ability judgments". Journal of Personality and Social Psychology 77 (2): 221–232. doi:10.1037/0022-35126.96.36.199. PMID 10474208. - Weinstein, Neil D. (1980). "Unrealistic optimism about future life events". Journal of Personality and Social Psychology 39 (5): 806–820. doi:10.1037/0022-35188.8.131.526. - Chambers, John R.; Windschitl, Paul D. (2004). "Biases in Social Comparative Judgments: The Role of Nonmotivated Factors in Above-Average and Comparative-Optimism Effects". Psychological Bulletin 130 (5): 813–838. doi:10.1037/0033-2909.130.5.813. - Chambers, John R.; Windschitl, Paul D.; Suls, Jerry (2003). "Egocentrism, Event Frequency, and Comparative Optimism: When what Happens Frequently is "More Likely to Happen to Me"". Personality and Social Psychology Bulletin 29 (11): 1343–1356. doi:10.1177/0146167203256870. - Kruger, Justin; Burrus, Jeremy (2004). "Egocentrism and focalism in unrealistic optimism (and pessimism)". Journal of Experimental Social Psychology 40 (3): 332–340. doi:10.1016/j.jesp.2003.06.002. - Taylor, Shelley E.; Brown, Jonathon D. (1988). "Illusion and well-being: A social psychological perspective on mental health". Psychological Bulletin 103 (2): 193–210. doi:10.1037/0033-2909.103.2.193. PMID 3283814. - Kahneman, Daniel (19 October 2011). "Don't Blink! The Hazards of Confidence". New York Times. Adapted from: Kahneman, Daniel (2011). Thinking, Fast and Slow. Farrar, Straus and Giroux. ISBN 978-1-4299-6935-2. - Plous, Scott (1993). The Psychology of Judgment and Decision Making. McGraw-Hill Education. ISBN 978-0-07-050477-6. - Thompson, Leigh; Loewenstein, George (1992). "Egocentric interpretations of fairness and interpersonal conflict" (PDF). Organizational Behavior and Human Decision Processes 51 (2): 176–197. doi:10.1016/0749-5978(92)90010-5. - Babcock, Linda C.; Olson, Craig A. (1992). "The Causes of Impasses in Labor Disputes". Industrial Relations 31 (2): 348–360. doi:10.1111/j.1468-232X.1992.tb00313.x. - Johnson, Dominic D. P. (2004). Overconfidence and War: The Havoc and Glory of Positive Illusions. Harvard University Press. ISBN 978-0-674-01576-0. - Aumann, Robert J. (1976). "Agreeing to Disagree". The Annals of Statistics 4 (6): 1236–1239. doi:10.1214/aos/1176343654. - Daniel, Kent; Hirshleifer, David; Subrahmanyam, Avanidhar (1998). "Investor Psychology and Security Market Under- and Overreactions". The Journal of Finance 53 (6): 1839–1885. doi:10.1111/0022-1082.00077. - Oskamp, Stuart (1965). "Overconfidence in case-study judgments" (PDF). Journal of Consulting Psychology 29 (3): 261–265. doi:10.1037/h0022125. Reprinted in Kahneman, Daniel; Slovic, Paul; Tversky, Amos, eds. (1982). Judgment Under Uncertainty: Heuristics and Biases. Cambridge University Press. pp. 287–293. ISBN 978-0-521-28414-1. - Radzevick, J. R.; Moore, D. A. (2009). "Competing To Be Certain (But Wrong): Social Pressure and Overprecision in Judgment" (PDF). Academy of Management Proceedings 2009 (1): 1–6. doi:10.5465/AMBPP.2009.44246308. - Fowler, James H.; Johnson, Dominic D. P. (2011-01-07). "On Overconfidence". Seed Magazine. ISSN 1499-0679. - Judge, Timothy A.; Locke, Edwin A.; Durham, Cathy C. (1997). "The dispositional causes of job satisfaction: A core evaluations approach". Research in Organizational Behavior 19. pp. 151–188. ISBN 978-0762301799. - Larrick, Richard P.; Burson, Katherine A.; Soll, Jack B. (2007). "Social comparison and confidence: When thinking you're better than average predicts overconfidence (and when it does not)". Organizational Behavior and Human Decision Processes 102 (1): 76–94. doi:10.1016/j.obhdp.2006.10.002. - Baron, Johnathan (1994). Thinking and Deciding. Cambridge University Press. pp. 219–224. ISBN 0-521-43732-6. - Gilovich, Thomas; Griffin, Dale; Kahneman, Daniel (2002). Heuristics and Biases: The Psychology of Intuitive Judgment. Cambridge University Press. ISBN 978-0-521-79679-8. - Sutherland, Stuart (2007). Irrationality. Pinter & Martin. pp. 172–178. ISBN 978-1-905177-07-3.

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Simile Teacher Resources Find Simile educational ideas and activities Showing 41 - 60 of 1,154 resources Digital Life 102 Catch your pupils' attention by starting class with a quiz about digital media! After going over their answers with a partner, individuals compose similes about the role of digital media in their lives and share these with the class.... 9th - 12th Social Studies & History CCSS: Designed Where the Red Fern Grows Chapter 1 Worksheet Break down Where the Red Fern Grows by Wilson Rawls into manageable chunks by focusing on plot points and literary elements in specific chapters. This resource is all about the first chapter, and asks pupils to use complete sentences to... 4th - 7th English Language Arts CCSS: Adaptable Pictures in Words: Poems of Tennyson and Noyes Young scholars examine how Tennyson and Noyes use words to paint vivid pictures. They read and analyze two poems, complete an online scavenger hunt, complete a worksheet, and write examples of alliteration, personification, metaphor,... 6th - 8th English Language Arts Simile and Metaphor- Poetry Toolbox Illustrate the connection between using figurative language (specifically similes and metaphors) and creating poetry. First this worksheet reviews the definition for each, and then writers create the endings to two examples provided.... 4th - 6th English Language Arts Play a figurative language game! Starting with a review of terms, this presentation quickly launches into a quiz game with hyper-linked answers. Simply click an answer to find out if it's wrong or right. The option to try again is always... 7th - 9th English Language Arts CCSS: Adaptable Setting the Tone with Figurative Language Explore figurative language with your secondary class. Extending a language arts unit, the lesson prompts middle schoolers to examine how an author's word choice establishes a story's tone, possibly using metaphors, similes,... 8th English Language Arts CCSS: Adaptable Like Water for Chocolate Guided Reading Worksheet Planning a novel unit on Like Water for Chocolate by Laura Esquivel? If so, this comprehensive guided reading worksheet might appeal to you. Steer readers through the entire novel with a routine of comprehension questions and analysis... 10th English Language Arts Here's an Instant Activity for September 19, 2005 Young writers study similes and then complete a writing activity for similes. They complete a teacher-led activity for similes and then work independently to write sentences using the given similes. A solid lesson plan! 1st - 2nd English Language Arts Narrative Practice: Similes, Metaphors, Personification Transform boring sentences with figurative language. Class members employ metaphor, simile, and personification to rewrite a series of provided sentences. Pupils can prepare for narrative essays and creative assignments by completing the... 7th - 12th English Language Arts Deciphering the Mechanics of Poetry After a review of poetic terms, groups are given an object and they create a poem using a simile, a metaphor, internal rhyme, end rhyme, alliteration, and personification. Groups then exchange objects and repeat the process. Consider... 7th - 9th English Language Arts

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Pneumonia may be suspected when the doctor examines the patient and hears coarse breathing or crackling sounds when listening to a portion of the chest with a stethoscope. There may be wheezing or the sounds of breathing may be faint in a particular area of the chest. A chest X-ray is usually ordered to confirm the diagnosis of pneumonia. The lungs have several segments referred to as lobes, usually two on the left and three on the right. When the pneumonia affects one of these lobes, it is often referred to as lobar pneumonia. Some pneumonias have a more patchy distribution that does not involve specific lobes. In the past, when both lungs were involved in the infection, the term "double pneumonia" was used. This term is rarely used today. Sputum samples can be collected and examined under the microscope. Pneumonia caused by bacteria or fungi can be detected by this examination. A sample of the sputum can be grown in special incubators, and the offending organism can be subsequently identified. It is important to understand that the sputum specimen must contain little saliva from the mouth and be delivered to the laboratory fairly quickly. Otherwise, overgrowth of noninfecting bacteria from the mouth may predominate. As we have used antibiotics in a broader uncontrolled fashion, more organisms are becoming resistant to the commonly used antibiotics. These types of cultures can help in directing more appropriate therapy. A blood test that measures white blood cell count (WBC) may be performed. An individual's white blood cell count can often give a hint as to the severity of the pneumonia and whether it is caused by bacteria or a virus. An increased number of neutrophils, one type of WBC, is seen in most bacterial infections, whereas an increase in lymphocytes, another type of WBC, is seen in viral infections, fungal infections, and some bacterial infections (like tuberculosis). Bronchoscopy is a procedure in which a thin, flexible, lighted viewing tube is inserted into the nose or mouth after a local anesthetic is administered. Using this device, the doctor can directly examine the breathing passages (trachea and bronchi). Simultaneously, samples of sputum or tissue from the infected part of the lung can be obtained. Sometimes, fluid collects in the pleural space around the lung as a result of the inflammation from pneumonia. This fluid is called a pleural effusion. If a significant amount of fluid develops, it can be removed. After numbing the skin with local anesthetic a needle is inserted into the chest cavity and fluid can be withdrawn and examined under the microscope. This procedure is called a thoracentesis. Often ultrasound is used to prevent complications from this procedure. In some cases, this fluid can become severely inflamed (parapneumonic effusion) or infected (empyema) and may need to be removed by more aggressive surgical procedures. Today, most often, this involves surgery through a tube or thoracoscope. This is referred to as video-assisted thoracoscopic surgery or VATS. This answer should not be considered medical advice...This answer should not be considered medical advice and should not take the place of a doctor’s visit. Please see the bottom of the page for more information or visit our Terms and Conditions. Archived: March 20, 2014 Thanks for your feedback. 361 of 402 found this helpful Read the Original Article: Pneumonia

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History of Manchuria Manchuria is a region in East Asia. Depending on the definition of its extent, Manchuria can either refer to a region falling entirely within China, or a larger region today divided between Northeast China and the Russian Far East. To differentiate between the two parts following the latter definition, the Russian part is also known as Outer Manchuria, while the Chinese part is known as Inner Manchuria. It is the homeland of the Manchu people, the designation introduced in 1636 for the Jurchen people, in origin a Tungusic people which took power in 17th century China, establishing the Qing dynasty that lasted until 1912. The population grew from about 1 million in 1750 to 5 million in 1850 and 14 million in 1900, largely because of the immigration of Chinese farmers. Lying at the juncture of the Chinese, Japanese and Russian spheres of influence, Manchuria has been a cockpit of conflict since the late 19th century. The Russian Empire established control over the northern part of Manchuria in 1860 (Beijing Treaty); it built a railway to consolidate its hold. Disputes over Manchuria and Korea led to the Russo-Japanese War of 1904 to 1905. The Japanese invaded Manchuria in 1931, setting up the puppet state of Manchukuo which became a centerpiece of the fast-growing Japanese Empire. The Soviet invasion of Manchuria in 1945 led to the overnight collapse of Japanese rule. Manchuria was a base of operations for the Mao Zedong's People's Liberation Army in the Chinese Civil War, leading to the formation of the People's Republic of China. In the Korean War, Chinese forces used Manchuria as a base to assist North Korea against the UN forces. During the Cold War era, Manchuria became a matter of contention, escalating to the Sino–Soviet border conflict in 1969. The Sino-Russian border dispute was resolved diplomatically only in 2004. In recent years there has been extensive scholarship on Manchuria in the 20th century, while the earlier period is less studied. Part of a series on the |History of Manchuria| - 1 Prehistory - 2 Early history - 3 History after 1860 - 4 Notes - 5 References - 6 Further reading |This section needs additional citations for verification. (May 2011)| At various times in antiquity, Han dynasty, Cao Wei dynasty, Western Jin dynasty, Tang dynasty and some other minor kingdoms of China had established control in parts of Manchuria. Various kingdoms of mixed proto-Korean and Tungusic ethnicity, such as Gojoseon, Buyeo, Goguryeo and Balhae were also established in parts of this area. Manchuria was the homeland of several Tungusic tribes, including the Ulchs and Nani. Various ethnic groups and their respective kingdoms, including the Sushen, Donghu, Xianbei, Wuhuan, Mohe and Khitan have risen to power in Manchuria. Finnish linguist Juha Janhunen believes that it was likely that a "Tungusic-speaking elite" ruled Goguryeo and Balhae, describing them as "protohistorical Manchurian states" and that part of their population was Tungusic, and that the area of southern Manchuria was the origin of Tungusic peoples and inhabited continuously by them since ancient times, and Janhunen rejected opposing theories of Goguryeo and Balhae's ethnic composition. From 698 to 926, the kingdom of Balhae occupied northern Korean peninsula and parts of Manchuria and Primorsky Krai, consisting of the people of the recently fallen Goguryeo kingdom of Korea as an aristocratic class, and the Nanai, the Udege, and the Evenks and descendants of the Tungus-speaking people as a lower class. Balhae was an early feudal medieval state of Eastern Asia, which developed its industry, agriculture, animal husbandry, and had its own cultural traditions and art. People of Balhae maintained political, economic and cultural contacts with the southern Chinese Tang dynasty, as well as Japan. Primorsky Krai settled at this moment by Northern Mohe tribes were incorporated to Balhae Kingdom under King Seon's reign (818–830) and put Balhae territory at its height. After subduing the Yulou Mohe (Hangul: 우루말갈 Hanja/Hanzi: 虞婁靺鞨 pinyin: Yúlóu Mòhé) first and the Yuexi Mohe (Hangul: 월희말갈 Hanja/Hanzi: 越喜靺鞨 pinyin: Yuèxǐ Mòhé) thereafter, King Seon administrated their territories by creating four prefectures : Solbin Prefecture, Jeongli Prefecture, Anbyeon Prefecture and Anwon Prefecture. Manchuria under the Liao and Jin With the Song dynasty to the south, the Khitan people of Western Manchuria, who probably spoke a language related to the Mongolic languages, created the Liao dynasty in the region, which went on to control adjacent parts of Northern China as well. In the early 12th century the Tungusic Jurchen people (the ancestors of the later Manchu people) originally lived in the forests in the eastern borderlands of the Liao Empire, and were Liao's tributaries, overthrew the Liao and formed the Jin dynasty (1115–1234). They went on to control parts of Northern China and Mongolia after a series of successful military campaigns. Most of the surviving Khitan either assimilated into the bulk of the Han Chinese and Jurchen population, or moved to Central Asia; however, it is thought that the Daur people, still living in northern Manchuria, are also descendants of the Khitans. The first Jin capital, Shangjing, located on the Ashi River not far from modern Harbin, was originally not much more than the city of tents, but in 1124 the second Jin emperor Wuqimai starting a major construction project, having his Chinese chief architect, Lu Yanlun, build a new city at this site, emulating, on a smaller scale, the Northern Song capital Bianjing (Kaifeng). When Bianjing fell to Jin troops in 1127, thousands of captured Song aristocrats (including the two Song emperors), scholars, craftsmen and entertainers, along with the treasures of the Song capital, were all taken to Shangjing (the Upper Capital) by the winners. Although the Jurchen ruler Wanyan Liang, spurred on by his aspirations to become the ruler of all China, moved the Jin capital from Shangjing to Yanjing (now Beijing) in 1153, and had the Shangjing palaces destroyed in 1157, the city regained a degree of significance under Wanyan Liang's successor, Emperor Shizong, who enjoyed visiting the region to get in touch with his Jurchen roots. The capital of the Jin, Zhongdu, was captured by the Mongols in 1215 at the Battle of Zhongdu. The Jin moved their capital Kaifeng, which fell to Mongols in 1233. In 1234, the Jin dynasty collapsed after the siege of Caizhou. The last emperor of the Jin, Emperor Modi, was killed while fighting the Mongols who had breached the walls of the city. Days earlier, his predecessor, Emperor Aizong, committed suicide because he was unable to escape the besieged city. Manchuria under the Mongols and the Yuan dynasty In 1211, after the conquest of Western Xia, Genghis Khan mobilized an army to conquer the Jin dynasty. His general Jebe and brother Qasar were ordered to reduce the Jurchen cities in Manchuria.[unreliable source] They successfully destroyed the Jin forts there. The Khitans under Yelü Liuge declared their allegiance to Genghis Khan and established nominally autonomous state in Manchuria in 1213. However, the Jin forces dispatched a punitive expedition against them. Jebe went there again and the Mongols pushed out the Jins. The Jin general, Puxian Wannu, rebelled against the Jin dynasty and founded the kingdom of Eastern Xia in Dongjing (Liaoyang) in 1215. He assumed the title Tianwang (天王; lit. Heavenly King) and the era name Tiantai (天泰). Puxian Wannu allied with the Mongols in order to secure his position. However, he revolted in 1222 after that and fled to an island while the Mongol army invaded Liaoxi, Liaodong, and Khorazm. As a result of an internal strife among the Khitans, they failed to accept Yelü Liuge's rule and revolted against the Mongol Empire. Fearing of the Mongol pressure, those Khitans fled to Goryeo without permission. But they were defeated by the Mongol-Korean alliance. Genghis Khan (1206–1227) gave his brothers and Muqali Chinese districts in Manchuria. Ögedei Khan's son Güyük crushed the Eastern Xia dynasty in 1233, pacifying southern Manchuria. Some time after 1234 Ögedei also subdued the Water Tatars in northern part of the region and began to receive falcons, harems and furs as taxation. The Mongols suppressed the Water Tatar rebellion in 1237. In Manchuria and Siberia, the Mongols used dogsled relays for their yam. The capital city Karakorum directly controlled Manchuria until the 1260s. During the Yuan dynasty (1271–1368), established by Kublai Khan by renaming his empire to "Great Yuan" in 1271, Manchuria was administered under the Liaoyang province. Descendants of Genghis Khan's brothers such as Belgutei and Hasar ruled the area under the Great Khans. The Mongols eagerly adopted new artillery and technologies. The world's earliest known cannon, dated 1282, was found in Mongol-held Manchuria. After the expulsion of the Mongols from China, the Jurchen clans remained loyal to Toghan Temür, the last Yuan emperor. In 1375, Naghachu, a Mongol commander of the Mongolia-based Northern Yuan dynasty in Liaoyang province invaded Liaodong with aims of restoring the Mongols to power. Although he continued to hold southern Manchuria, Naghachu finally surrendered to the Ming dynasty in 1387. In order to protect the northern border areas the Ming decided to "pacify" the Jurchens in order to deal with its problems with Yuan remnants along its northern border. The Ming solidified control only under Yongle Emperor (1402–1424). Manchuria during the Ming dynasty The Ming dynasty took control of Liaoning in 1371, just three years after the expulsion of the Mongols from Beijing. During the reign of the Yongle Emperor in the early 15th century, efforts were made to expand Chinese control throughout entire Manchuria by establishing the Nurgan Regional Military Commission. Mighty river fleets were built in Jilin City, and sailed several times between 1409 and ca. 1432, commanded by the eunuch Yishiha down the Songhua and the Amur all the way to the mouth of the Amur, getting the chieftains of the local tribes to swear allegiance to the Ming rulers. Soon after the death of the Yongle Emperor the expansion policy of the Ming was replaced with that of retrenchment in southern Manchuria (Liaodong). Around 1442, a defence wall was constructed to defend the northwestern frontier of Liaodong from a possible threat from the Jurched-Mongol Oriyanghan. In 1467–68 the wall was expanded to protect the region from the northeast as well, against attacks from Jianzhou Jurchens. Although similar in purpose to the Great Wall of China, this "Liaodong Wall" was of a simpler design. While stones and tiles were used in some parts, most of the wall was in fact simply an earthen dike with moats on both sides. Chinese cultural and religious influence such as Chinese New Year, the "Chinese god", Chinese motifs like the dragon, spirals, scrolls, and material goods like agriculture, husbandry, heating, iron cooking pots, silk, and cotton spread among the Amur natives like the Udeghes, Ulchis, and Nanais. Starting in the 1580s, a Jianzhou Jurchens chieftain Nurhaci (1558–1626), originally based in the Hurha River valley northeast of the Ming Liaodong Wall, started to unify Jurchen tribes of the region. Over the next several decades, the Jurchen (later to be called Manchu), took control over most of Manchuria, the cities of the Ming Liaodong falling to the Jurchen one after another. In 1616, Nurhaci declared himself a khan, and founded the Later Jin dynasty (which his successors renamed in 1636 to Qing dynasty). Manchuria during the Qing dynasty The process of unification of the Jurchen people completed by Nurhaci was followed by his son's, Hong Taiji, energetic expansion into Outer Manchuria. The conquest of the Amur basin people was completed after the defeat of the Evenk chief Bombogor, in 1640. In 1644, the Manchus took Beijing, overthrowing the Ming dynasty and soon established the Qing dynasty rule (1644–1912) over all of China. The Manchus ruled all of China, but they treated their homeland of Manchuria to a special status and ruled it separately. The "Banner" system that in China involved military units originated in Manchuria and was used as a form of government. During the Qing dynasty, the area of Manchuria was known as the "three eastern provinces" (東三省, dong san sheng) since 1683 when Jilin and Heilongjiang were separated even though it was not until 1907 that they were turned into actual provinces. The area of Manchuria was then converted into three provinces by the late Qing government in 1907. For decades the Manchu rulers tried to prevent large-scale immigration of Han Chinese, but they failed and the southern parts developed agricultural and social patterns similar to those of North China. Manchuria's population grew from about 1 million in 1750 to 5 million in 1850 and 14 million in 1900, largely because of the immigration of Chinese farmers. The Manchus became a small element in their homeland, although they retained political control until 1900. The region was separated from China proper by the Inner Willow Palisade, a ditch and embankment planted with willows intended to restrict the movement of the Han Chinese into Manchuria during the Qing dynasty, as the area was off-limits to the Han until the Qing started colonizing the area with them later on in the dynasty's rule. This movement of the Han Chinese to Manchuria is called Chuang Guandong. The Manchu area was still separated from modern-day Inner Mongolia by the Outer Willow Palisade, which kept the Manchu and the Mongols separate. However, the Qing rule saw a massive increase of Han Chinese settlement, both legal and illegal, in Manchuria. As Manchu landlords needed the Han peasants to rent their land and grow grain, most Han migrants were not evicted. During the 18th century, Han peasants farmed 500,000 hectares of privately owned land in Manchuria and 203,583 hectares of lands which were part of courier stations, noble estates, and banner lands, in garrisons and towns in Manchuria the Han Chinese made up 80% of the population. Han farmers were resettled from north China by the Qing to the area along the Liao River in order to restore the land to cultivation. To the north, the boundary with Russian Siberia was fixed by the Treaty of Nerchinsk (1689) as running along the watershed of the Stanovoy Mountains. South of the Stanovoy Mountains, the basin of the Amur and its tributaries belonged to the Qing Empire. North of the Stanovoy Mountains, the Uda Valley and Siberia belonged to the Russian Empire. In 1858, a weakening Qing Empire was forced to cede Manchuria north of the Amur to Russia under the Treaty of Aigun; however, Qing subjects were allowed to continue to reside, under the Qing authority, in a small region on the now-Russian side of the river, known as the Sixty-Four Villages East of the River. In 1860, at the Convention of Peking, the Russians managed to acquire a further large slice of Manchuria, east of the Ussuri River. As a result, Manchuria was divided into a Russian half known as "Outer Manchuria", and a remaining Chinese half known as "Inner Manchuria". In modern literature, "Manchuria" usually refers to Inner (Chinese) Manchuria. (cf. Inner and Outer Mongolia). As a result of the Treaties of Aigun and Peking, China lost access to the Sea of Japan. The Qing government began to actively encourage Han Chinese citizens to move into Manchuria since then. The Manza War in 1868 was the first attempt by Russia to expel Chinese from territory it controlled. Hostilities broke out around Vladivostok when the Russians tried to shut off gold mining operations and expel Chinese workers there. The Chinese resisted a Russian attempt to take Askold Island and in response, 2 Russian military stations and 3 Russian towns were attacked by the Chinese, and the Russians failed to oust the Chinese. However, the Russians finally managed it from them in 1892 History after 1860 By the 19th century, Manchu rule had become increasingly sinicized and, along with other borderlands of the Qing Empire such as Mongolia and Tibet, came under the influence of Japan and the European powers as the Qing dynasty grew weaker and weaker. Russian and Japanese encroachment Inner Manchuria also came under strong Russian influence with the building of the Chinese Eastern Railway through Harbin to Vladivostok. Some poor Korean farmers moved there. In Chuang Guandong many Han farmers, mostly from Shandong peninsula moved there, attracted by cheap farmland that was ideal for growing soybeans. During the Boxer Rebellion in 1899–1900, Russian soldiers killed ten-thousand Chinese (Manchu, Han Chinese and Daur people) living in Blagoveshchensk and Sixty-Four Villages East of the River. In revenge, the Chinese Honghuzi conducted guerilla warfare against the Russian occupation of Manchuria and sided with Japan against Russia during the Russo-Japanese War. Japan replaced Russian influence in the southern half of Inner Manchuria as a result of the Russo-Japanese War in 1904–1905. Most of the southern branch of the Chinese Eastern Railway (the section from Changchun to Port Arthur (Japanese: Ryojun)) was transferred from Russia to Japan, and became the South Manchurian Railway. Jiandao (in the region bordering Korea), was handed over to Qing dynasty as a compensation for the South Manchurian Railway. From 1911 to 1931 Manchuria was nominally part of the Republic of China. In practice it was controlled by Japan, which worked through local warlords. Japanese influence extended into Outer Manchuria in the wake of the Russian Revolution of 1917, but Outer Manchuria came under Soviet control by 1925. Japan took advantage of the disorder following the Russian Revolution to occupy Outer Manchuria, but Soviet successes and American economic pressure forced Japanese withdrawal. It was reported that among Banner people, both Manchu and Chinese (Hanjun) in Aihun, Heilongjiang in the 1920s, would seldom marry with Han civilians, but they (Manchu and Chinese Bannermen) would mostly intermarry with each other. Owen Lattimore reported that, during his January 1930 visit to Manchuria, he studied a community in Jilin (Kirin), where both Manchu and Chinese bannermen were settled at a town called Wulakai, and eventually the Chinese Bannermen there could not be differentiated from Manchus since they were effectively Manchufied. The Han civilian population was in the process of absorbing and mixing with them when Lattimore wrote his article. Manchuria was (and still is) an important region for its rich mineral and coal reserves, and its soil is perfect for soy and barley production. For Japan, Manchuria became an essential source of raw materials. 1931 Japanese invasion and Manchukuo Around the time of World War I, Zhang Zuolin, a former bandit (Honghuzi) established himself as a powerful warlord with influence over most of Manchuria. He was inclined to keep his Manchu army under his control and to keep Manchuria free of foreign influence. The Japanese tried and failed to assassinate him in 1916. They finally succeeded in June 1928. Following the Mukden Incident in 1931 and the subsequent Japanese invasion of Manchuria, Inner Manchuria was proclaimed to be Manchukuo, a puppet state under the control of the Japanese army. The last Manchu emperor, Puyi, was then placed on the throne to lead a Japanese puppet government in the Wei Huang Gong, better known as "Puppet Emperor's Palace". Inner Manchuria was thus detached from China by Japan to create a buffer zone to defend Japan from Russia's Southing Strategy and, with Japanese investment and rich natural resources, became an industrial domination. Under Japanese control Manchuria was one of the most brutally run regions in the world, with a systematic campaign of terror and intimidation against the local Russian and Chinese populations including arrests, organised riots and other forms of subjugation. The Japanese also began a campaign of emigration to Manchukuo; the Japanese population there rose from 240,000 in 1931 to 837,000 in 1939 (the Japanese had a plan to bring in 5 million Japanese settlers into Manchukuo). Hundreds of Manchu farmers were evicted and their farms given to Japanese immigrant families. Manchukuo was used as a base to invade the rest of China in 1937-40. At the end of the 1930s, Manchuria was a trouble spot with Japan, clashing twice with the Soviet Union. These clashes - at Lake Khasan in 1938 and at Khalkhin Gol one year later - resulted in many Japanese casualties. The Soviet Union won these two battles and a peace agreement was signed. However, the regional unrest endured.[clarification needed] After World War II After the atomic bombing of Hiroshima in August 1945, the Soviet Union invaded from Soviet Outer Manchuria as part of its declaration of war against Japan. From 1945 to 1948, Inner Manchuria was a base area for the Chinese People's Liberation Army in the Chinese Civil War. With the encouragement of the Soviet Union, Manchuria was used as a staging ground during the Chinese Civil War for the Communist Party of China, which emerged victorious in 1949. During the Korean War of the 1950s, 300,000 soldiers of the Chinese People's Liberation Army crossed the Sino-Korean border from Manchuria to repulse UN forces led by the United States from North Korea. In the 1960s, Manchuria's border with the Soviet Union became the site of the most serious tension between the Soviet Union and China. The treaties of 1858 and 1860, which ceded territory north of the Amur, were ambiguous as to which course of the river was the boundary. This ambiguity led to dispute over the political status of several islands. This led to armed conflict in 1969, called the Sino-Soviet border conflict. With the end of the Cold War, this boundary issue was discussed through negotiations. In 2004, Russia agreed to transfer Yinlong Island and one half of Heixiazi Island to China, ending an enduring border dispute. Both islands are found at the confluence of the Amur and Ussuri Rivers, and were until then administered by Russia and claimed by China. The event was meant to foster feelings of reconciliation and cooperation between the two countries by their leaders, but it has also provoked different degrees of dissent on both sides. Russians, especially Cossack farmers of Khabarovsk, who would lose their ploughlands on the islands, were unhappy about the apparent loss of territory. Meanwhile, some Chinese have criticised the treaty as an official acknowledgement of the legitimacy of Russian rule over Outer Manchuria, which was ceded by the Qing dynasty to Imperial Russia under a series of Unequal Treaties, which included the Treaty of Aigun in 1858 and the Convention of Peking in 1860, in order to exchange exclusive usage of Russia's rich oil resources. The transfer was carried out on October 14, 2008. - Janhunen (2006), p. 109. - Li (2001). - Tao (1976), pp. 28–32. - Tao (1976), p. 44. - Tao (1976), p. 78–79. - Franke (1994), p. 254. - Franke (1994), pp. 264–265. - Shanley (2008), p. 144. - Atwood (2004), pp. 341–342. - Berger (2003), p. 25. - Kamal (2003), p. 76. - Atwood (2004), p. 354. - Tsai (1996), pp. 129–130. - Edmonds (1985), pp. 38–40. - Forsyth (1994), p. 214. - Shao (2011), pp. 25-67. - Clausen & Thøgersen (1995), p. 7. - Isett (2007), p. 33. - Richards 2003, p. 141. - Anderson (2000), p. 504. - Lomanov (2005:89–90) Probably the first clash between the Russians and Chinese occurred in 1868. It was called the Manza War, Manzovskaia voina. "Manzy" was the Russian name for the Chinese population in those years. In 1868, the local Russian government decided to close down goldfields near Vladivostok, in the Gulf of Peter the Great, where 1,000 Chinese were employed. The Chinese decided that they did not want to go back, and resisted. The first clash occurred when the Chinese were removed from Askold Island, in the Gulf of Peter the Great. They organized themselves and raided three Russian villages and two military posts. For the first time, this attempt to drive the Chinese out was unsuccessful. - "俄军惨屠海兰泡华民五千余人(1900年)". News.163.com. Retrieved 2010-05-18. - (2008-10-15 16:41:01) (2008-10-15). "江东六十四屯". Blog.sina.com.cn. Retrieved 2010-05-18. - Riechers (2001). - Rhoads (2011), p. 263. - Lattimore (1933), p. 272. - Behr (1987), p. 202. - Behr (1987), p. 168. - Duara (2006). - Behr (1987), p. 204. - Battlefield – Manchuria - "Handover of Russian islands to China seen as effective diplomacy | Top Russian news and analysis online | 'RIA Novosti' newswire". En.rian.ru. 2008-10-14. 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"The rise of the Mongolian empire and Mongolian rule in north China". In Denis C. Twitchett; Herbert Franke; John King Fairbank. The Cambridge History of China: Volume 6, Alien Regimes and Border States, 710–1368. Cambridge University Press. pp. 321–413. ISBN 978-0-521-24331-5. - Crossley, Pamela Kyle. The Manchus (2002) excerpt and text search; review - Im, Kaye Soon. "The Development of the Eight Banner System and its Social Structure," Journal of Social Sciences & Humanities (1991), Issue 69, pp 59–93 - Lattimore, Owen. Manchuria: Cradle of Conflict (1932). - Matsusaka, Yoshihisa Tak. The Making of Japanese Manchuria, 1904-1932 (Harvard East Asian Monographs, 2003) - Mitter, Rana. The Manchurian Myth: Nationalism, Resistance, and Collaboration in Modern China (2000). - Sun, Kungtu C. The economic development of Manchuria in the first half of the twentieth century (Havard U.P. 1969, 1973), 123pp search text - Tamanoi, Mariko, ed. 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https://en.wikipedia.org/wiki/History_of_Manchuria

|This article needs additional citations for verification. (January 2007)| English number words include numerals and various words derived from them, as well as a large number of words borrowed from other languages. |2||two||12||twelve (a dozen)||20||twenty (a score)| |4||four||14||fourteen||40||forty (no "u")| |5||five||15||fifteen (note "f", not "v")||50||fifty (note "f", not "v")| |8||eight||18||eighteen (only one "t")||80||eighty (only one "t")| |9||nine||19||nineteen||90||ninety (note the "e")| If a number is in the range 21 to 99, and the second digit is not zero, one typically writes the number as two words separated by a hyphen. In English, the hundreds are perfectly regular, except that the word hundred remains in its singular form regardless of the number preceding it. So too are the thousands, with the number of thousands followed by the word "thousand" |10,000||ten thousand or (rarely used) a myriad| |100,000||one hundred thousand or one lakh (Indian English)| |999,000||nine hundred and ninety-nine thousand (inclusively British English, Irish English, Australian English, and New Zealand English) nine hundred ninety-nine thousand (American English) |10,000,000||ten million or one crore (Indian English)| In American usage, four-digit numbers with non-zero hundreds are often named using multiples of "hundred" and combined with tens and ones: "One thousand one", "Eleven hundred three", "Twelve hundred twenty-five", "Four thousand forty-two", or "Ninety-nine hundred ninety-nine." In British usage, this style is common for multiples of 100 between 1,000 and 2,000 (e.g. 1,500 as "fifteen hundred") but not for higher numbers. Americans may pronounce four-digit numbers with non-zero tens and ones as pairs of two-digit numbers without saying "hundred" and inserting "oh" for zero tens: "twenty-six fifty-nine" or "forty-one oh five". This usage probably evolved from the distinctive usage for years; "nineteen-eighty-one", or from four-digit numbers used in the American telephone numbering system which were originally two letters followed by a number followed by a four-digit number, later by a three-digit number followed by the four-digit number. It is avoided for numbers less than 2500 if the context may mean confusion with time of day: "ten ten" or "twelve oh four". Intermediate numbers are read differently depending on their use. Their typical naming occurs when the numbers are used for counting. Another way is for when they are used as labels. The second column method is used much more often in American English than British English. The third column is used in British English but rarely in American English (although the use of the second and third columns is not necessarily directly interchangeable between the two regional variants). In other words, British English and American English can seemingly agree, but it depends on a specific situation (in this example, bus numbers). |Common British vernacular||Common American vernacular||Common British vernacular| |"How many marbles do you have?"||"What is your house number?"||"Which bus goes to the high street?"| |101||"A hundred and one."||"One-oh-one." Here, "oh" is used for the digit zero. |109||"A hundred and nine."||"One-oh-nine."||"One-oh-nine."| |110||"A hundred and ten."||"One-ten."||"One-one-oh."| |117||"A hundred and seventeen."||"One-seventeen."||"One-one-seven."| |120||"A hundred and twenty."||"One-twenty."||"One-two-oh", "One-two-zero."| |152||"A hundred and fifty-two."||"One-fifty-two."||"One-five-two."| |208||"Two hundred and eight."||"Two-oh-eight."||"Two-oh-eight."| |334||"Three hundred and thirty-four."||"Three-thirty-four."||"Three-three-four."| Note: When writing a cheque (or check), the number 100 is always written "one hundred". It is never "a hundred". In American English, many students are taught not to use the word and anywhere in the whole part of a number, so it is not used before the tens and ones. It is instead used as a verbal delimiter when dealing with compound numbers. Thus, instead of "three hundred and seventy-three", one would say "three hundred seventy-three". Despite this rule, the and is used by some Americans in reading numbers containing tens and ones as an alternative variant. For details, see American and British English differences. For numbers above a million, there are two different systems for naming numbers in English (for the use of prefixes such as kilo- for a thousand, mega- for a million, milli- for a thousandth, etc. see SI units): - the long scale (decreasingly used in British English) designates a system of numeric names in which a thousand million is called a ‘‘milliard’’ (but the latter usage is now rare), and ‘‘billion’’ is used for a million million. - the short scale (always used in American English and increasingly in British English) designates a system of numeric names in which a thousand million is called a ‘‘billion’’, and the word ‘‘milliard’’ is not used. |Short scale||Long scale||Indian (or South Asian) English |1,000,000||106||one million||one million||ten lakh| a thousand million a thousand million |one hundred crore a thousand billion a million million |one lakh crore a thousand trillion a thousand billion |ten crore crore a thousand quadrillion a million billion |ten thousand crore crore a thousand quintillion a thousand trillion |one crore crore crore| The numbers past a trillion in short scale system, in ascending powers of ten, are as follows: quadrillion, quintillion, sextillion, septillion, octillion, nonillion, decillion, undecillion, duodecillion, tredecillion, quattuordecillion, and quindecillion (that's 10 to the 48th, or a one followed by 48 zeros). The highest number listed on Robert Munafo's table, is a milli-millillion. That's 10 to the 3000003rd. Although British English has traditionally followed the long-scale numbering system, the short-scale usage has become increasingly common in recent years. For example, the UK Government and BBC websites use the newer short-scale values exclusively. The terms arab, kharab, padm and shankh are more commonly found in old sections of Indian Mathematics. Here are some approximate composite large numbers in American English: |1,200,000||1.2 million||one point two million| |3,000,000||3 million||three million| |250,000,000||250 million||two hundred fifty million| |6,400,000,000||6.4 billion||six point four billion| |23,380,000,000||23.38 billion||twenty-three point three eight billion| Often, large numbers are written with (preferably non-breaking) half-spaces or thin spaces separating the thousands (and, sometimes, with normal spaces or apostrophes) instead of commas—to ensure that confusion is not caused in countries where a decimal comma is used. Thus, a million is often written 1 000 000. A few numbers have special names (in addition to their regular names): - 0: has several other names, depending on context: - zero: formal scientific usage - naught / nought: mostly British usage - aught: Mostly archaic but still occasionally used when a digit in mid-number is 0 (as in "thirty-aught-six", the .30-06 Springfield rifle cartridge and by association guns that fire it) - oh: used when spelling numbers (like telephone, bank account, bus line [British: bus route]) - nil: in general sport scores, British usage ("The score is two–nil.") - nothing: in general sport scores, American usage ("The score is two–nothing.") - null: used technically to refer to an object or idea related to nothingness. The 0th aleph number () is pronounced "aleph-null". - love: in tennis, badminton, squash and similar sports (origin disputed, often said to come from French l'œuf, "egg"; but the Oxford English Dictionary mentions the phrase for love, meaning nothing is at risk) - zilch, nada (from Spanish), zip: used informally when stressing nothingness; this is true especially in combination with one another ("You know nothing—zero, zip, nada, zilch!"); American usage - nix: also used as a verb; mostly American usage - cypher / cipher: archaic, from French chiffre, in turn from Arabic sifr, meaning zero - goose egg (informal) - duck (used in cricket when a batsman is dismissed without scoring) - blank the half of a domino tile with no pips - ace in certain sports and games, as in tennis or golf, indicating success with one stroke, and the face of a die, playing card or domino half with one pip - birdie in golf denotes one stroke less than par, and bogey, one stroke more than par - linear the degree of a polynomial is 1; also for explicitly denoting the first power of a unit: linear meter - unity in mathematics - protagonist first actor in theater of Ancient Greece - brace, from Old French "arms" (the plural of arm), as in "what can be held in two arms". - deuce the face of a die, playing card or domino half with two pips - eagle in golf denotes two strokes less than par - quadratic the degree of a polynomial is 2 - also square or squared for denoting the second power of a unit: square meter or meter squared - penutimate, second from the end - deuteragonist second actor in theater of Ancient Greece - trey the face of a die or playing card with three pips, a three-point field goal in basketball, nickname for the third carrier of the same personal name in a family - trips: three-of-a-kind in a poker hand. a player has three cards with the same numerical value - cubic the degree of a polynomial is 3 - also cube or cubed for denoting the third power of a unit: cubic meter or meter cubed - albatross in golf denotes three strokes less than par. Sometimes called double eagle - hat-trick or hat trick: achievement of three feats in sport or other contexts - antepenultimate third from the end - tritagonist third actor in theater of Ancient Greece - turkey in bowling, three consecutive strikes - cater: (rare) the face of a die or playing card with four pips - quartic or biquadratic the degree of a polynomial is 4 - quad (short for quadruple or the like) several specialized sets of four, such as four of a kind in poker, a carburetor with four inputs, etc., - condor in golf denotes four strokes less than par - preantepenultimate fourth from the end - cinque or cinq (rare) the face of a die or playing card with five pips - nickel (informal American, from the value of the five-cent US nickel, but applied in non-monetary references) - quintic the degree of a polynomial is 5 - quint (short for quintuplet or the like) several specialized sets of five, such as quintuplets, etc. - 11: a banker's dozen - 12: a dozen (first power of the duodecimal base), used mostly in commerce - 13: a baker's dozen - 20: a score (first power of the vigesimal base), nowadays archaic; famously used in the opening of the Gettysburg Address: "Four score and seven years ago..." The Number of the Beast in the King James Bible is rendered "Six hundred threescore and six". Also in The Book of Common Prayer, Psalm 90 as used in the Burial Service - "The days of our age are threescore years and ten; ...." - 50: half a century, literally half of a hundred, usually used in cricket scores. Normally referred to as a 'half-century' without the 'a'. - 60: a shock: historical commercial count, described as "three scores". - 110: eleventy (as 11 × 10) - A great hundred or long hundred (twelve tens; as opposed to the small hundred, i.e. 100 or ten tens), also called small gross (ten dozens), both archaic - Also sometimes referred to as duodecimal hundred, although that could literally also mean 144, which is twelve squared - Twelfty or twelvety (as 12 × 10) - 144: a gross (a dozen dozens, second power of the duodecimal base), used mostly in commerce - a grand, colloquially used especially when referring to money, also in fractions and multiples, e.g. half a grand, two grand, etc. Grand can also be shortened to "G" in many cases. - K, originally from the abbreviation of kilo-, e.g. "He only makes $20K a year." - 1728: a great gross (a dozen gross, third power of the duodecimal base), used mostly in commerce - 10,000: a myriad (a hundred hundred), commonly used in the sense of an indefinite very high number - 100,000: a lakh (a hundred thousand), loanword used mainly in Indian English - 10,000,000: a crore (a hundred lakh), loanword used mainly in Indian English - 10100: googol (1 followed by 100 zeros), used in mathematics; not to be confused with the name of the company Google (which was originally a misspelling of googol) - 10googol: googolplex (1 followed by a googol of zeros) - 10googolplex: googolplexplex (1 followed by a googolplex of zeros) Combinations of numbers in most sports scores are read as in the following examples: - 1–0 British English: one-nil; American English: one-nothing, one-zip, or one-zero - 0–0 British English: nil-nil, or more rarely nil all; American English: zero-zero or nothing-nothing, (occasionally scoreless or no score) - 2–2 two-two or two all; American English also twos, two to two, even at two, or two up. Naming conventions of Tennis scores (and related sports) are different from other sports. A few numbers have specialised multiplicative numbers expresses how many times some event happens (adverbs): Compare these specialist multiplicative numbers to express how many times some thing exists (adjectives): Other examples are given in the Specialist Numbers. The name of a negative number is the name of the corresponding positive number preceded by "minus" or (American English) "negative". Thus −5.2 is "minus five point two" or "negative five point two". For temperatures, North Americans colloquially say "below" — short for "below zero" — so a temperature of −5° is "five below" (in contrast, for example, to "two above" for 2°, occasionally used for emphasis when referring to several temperatures or ranges both positive and negative. This is particularly common in Canada where the use of Celsius in weather forecasting means that temperatures can regularly drift above and below zero at certain times of year.) |Look up Appendix:English ordinal numbers in Wiktionary, the free dictionary.| Ordinal numbers refer to a position in a series. Common ordinals include: |0th||zeroth or noughth (see below)||10th||tenth| |2nd||second||12th||twelfth (note "f", not "v")||20th||twentieth| |8th||eighth (only one "t")||18th||eighteenth||80th||eightieth| |9th||ninth (no "e")||19th||nineteenth||90th||ninetieth| Ordinal numbers such as 21st, 33rd, etc., are formed by combining a cardinal ten with an ordinal unit. Higher ordinals are not often written in words, unless they are round numbers (thousandth, millionth, billionth). They are written using digits and letters as described below. Here are some rules that should be borne in mind. - The suffixes -th, -st, -nd and -rd are occasionally written superscript above the number itself. - If the tens digit of a number is 1, then write "th" after the number. For example: 13th, 19th, 112th, 9,311th. - If the tens digit is not equal to 1, then use the following table: |If the units digit is:||0||1||2||3||4||5||6||7||8||9| |write this after the number||th||st||nd||rd||th||th||th||th||th||th| - For example: 2nd, 7th, 20th, 23rd, 52nd, 135th, 301st. These ordinal abbreviations are actually hybrid contractions of a numeral and a word. 1st is "1" + "st" from "first". Similarly, "nd" is used for "second" and "rd" for "third". In the legal field and in some older publications, the ordinal abbreviation for "second" and "third" is simply "d". - For example: 42d, 33d, 23d. NB: The practice of using "d" to denote "second" and "third" is still often followed in the numeric designations of units in the US armed forces, for example, 533d Squadron. Any ordinal name that doesn't end in "first", "second", or "third", ends in "th". There are a number of ways to read years. The following table offers a list of valid pronunciations and alternate pronunciations for any given year of the Gregorian calendar. |Year||Most common pronunciation method||Alternative methods| |1 BC||(The year) One Before Christ (BC)||1 before the Common era (BCE)| |1||(The year) One Anno Domini (AD)||of the Common era (CE) In the year of Our Lord 1 Two hundred (and) thirty-five Nine hundred (and) eleven Nine hundred (and) ninety-nine |1000||One thousand||Ten hundred |1004||One thousand (and) four||Ten oh-four| |1010||Ten ten||One thousand (and) ten| |1050||Ten fifty||One thousand (and) fifty| One thousand, two hundred (and) twenty-five |1900||Nineteen hundred||One thousand, nine hundred |1901||Nineteen oh-one||Nineteen hundred (and) one One thousand, nine hundred (and) one Nineteen aught one |1919||Nineteen nineteen||Nineteen hundred (and) nineteen One thousand, nine hundred (and) nineteen |1999||Nineteen ninety-nine||Nineteen hundred (and) ninety-nine One thousand, nine hundred (and) ninety-nine |2000||Two thousand||Twenty hundred |2001||Two thousand (and) one||Twenty oh-one Twenty hundred (and) one |2009||Two thousand (and) nine||Twenty oh-nine Twenty hundred (and) nine |2010||Two thousand (and) ten |Twenty hundred (and) ten Fractions and decimals In spoken English, ordinal numbers are also used to quantify the denominator of a fraction. Thus "fifth" can mean the element between fourth and sixth, or the fraction created by dividing the unit into five pieces. In this usage, the ordinal numbers can be pluralized: one seventh, two sevenths. The sole exception to this rule is division by two. The ordinal term "second" can only refer to location in a series; for fractions English speakers use the term 'half' (plural "halves"). |1/10 or 0.1||one tenth| |2/10 or 0.2||two tenths| |1/4||one quarter or (mainly American English) one fourth| |3/10 or 0.3||three tenths| |4/10 or 0.4||four tenths| |6/10 or 0.6||six tenths| |7/10 or 0.7||seven tenths| |3/4||three quarters or three fourths| |8/10 or 0.8||eight tenths| |9/10 or 0.9||nine tenths| Alternatively, and for greater numbers, one may say for 1/2 "one over two", for 5/8 "five over eight", and so on. This "over" form is also widely used in mathematics. Numbers with a decimal point may be read as a cardinal number, then "and", then another cardinal number followed by an indication of the significance of the second cardinal number (mainly U.S.); or as a cardinal number, followed by "point", and then by the digits of the fractional part. The indication of significance takes the form of the denominator of the fraction indicating division by the smallest power of ten larger than the second cardinal. This is modified when the first cardinal is zero, in which case neither the zero nor the "and" is pronounced, but the zero is optional in the "point" form of the fraction. - 0.002 is "two thousandths" (mainly U.S.); or "point zero zero two", "point oh oh two", "nought point zero zero two", etc. - 3.1416 is "three point one four one six" - 99.3 is "ninety-nine and three tenths" (mainly U.S.); or "ninety-nine point three". In English the decimal point was originally printed in the center of the line (0·002), but with the advent of the typewriter it was placed at the bottom of the line, so that a single key could be used as a full stop/period and as a decimal point. In many non-English languages a full-stop/period at the bottom of the line is used as a thousands separator with a comma being used as the decimal point. Fractions together with an integer are read as follows: - 1 1/2 is "one and a half" - 6 1/4 is "six and a quarter" - 7 5/8 is "seven and five eighths" A space is required between the whole number and the fraction; however, if a special fraction character is used like "½", then the space can be done without, e.g. - 9 1/2 Whether to use digits or words With very little deviation, most grammatical texts rule that the numbers zero to nine inclusive should be "written out" – meaning instead of "1" and "2", one would write "one" and "two". - Example: "I have two apples." (Preferred) - Example: "I have 2 apples." After "nine", one can head straight back into the 10, 11, 12, etc., although some write out the numbers until "twelve". - Example: "I have 28 grapes." (Preferred) - Example: "I have twenty-eight grapes." Another common usage is to write out any number that can be expressed as one or two words, and use figures otherwise. - "There are six million dogs." (Preferred) - "There are 6,000,000 dogs." - "That is one hundred and twenty-five oranges." (British English) - "That is one hundred twenty-five oranges." (US-American English) - "That is 125 oranges." (Preferred) Numbers at the beginning of a sentence should also be written out. The above rules are not always used. In literature, larger numbers might be spelled out. On the other hand, digits might be more commonly used in technical or financial articles, where many figures are discussed. In particular, the two different forms should not be used for figures that serve the same purpose; for example, it is inelegant to write, "Between day twelve and day 15 of the study, the population doubled." Colloquial English has a small vocabulary of empty numbers that can be employed when there is uncertainty as to the precise number to use, but it is desirable to define a general range: specifically, the terms "umpteen", "umpty", and "zillion". These are derived etymologically from the range affixes: - "-teen" (designating the range as being between 10 and 20) - "-ty" (designating the range as being in one of the decades between 20 and 100) - "-illion" (designating the range as being above 1,000,000; or, more generally, as being extremely large). The prefix "ump-" is added to the first two suffixes to produce the empty numbers "umpteen" and "umpty": it is of uncertain origin. There is a noticeable absence of an empty number in the hundreds range. Usage of empty numbers: - The word "umpteen" may be used as an adjective, as in "I had to go to umpteen stores to find shoes that fit." It can also be used to modify a larger number, usually "million", as in "Umpteen million people watched the show; but they still cancelled it." - "Umpty" is not in common usage. It can appear in the form "umpty-one" (parallelling the usage in such numbers as "twenty-one"), as in "There are umpty-one ways to do it wrong." "Umpty-ump" is also heard, though "ump" is never used by itself. - The word "zillion" may be used as an adjective, modifying a noun. The noun phrase normally contains the indefinite article "a", as in "There must be a zillion sites on the World Wide Web." - The plural "zillions" designates a number indefinitely larger than "millions" or "billions". In this case, the construction is parallel to the one for "millions" or "billions", with the number used as a plural count noun, followed by a prepositional phrase with "of", as in "Out in the countryside, the night sky is filled with zillions of stars." - Empty numbers are sometimes made up, with obvious meaning: "squillions" is obviously an empty, but very large, number; a "squintillionth" would be a very small number. - Some empty numbers may be modified by actual numbers, such as "four zillion", and are used for jest, exaggeration, or to relate abstractly to actual numbers. - Empty numbers are colloquial, and primarily used in oral speech or informal contexts. They are inappropriate in formal or scholarly usage. See also Placeholder name. - Indefinite and fictitious numbers - List of numbers - Long and short scales - Names of large numbers - Number prefixes and their derivatives - Natural number - "Hat trick, n.". Oxford English Dictionary. Oxford University Press. Retrieved 26 December 2014. - "Shock, n.2". Oxford English Dictionary. Oxford University Press. Retrieved 26 December 2014. - What is a partitive numeral? - Gary Blake and Robert W. Bly, The Elements of Technical Writing, pg. 22. New York: Macmillan Publishers, 1993. ISBN 0020130856 |Look up Appendix:English numerals in Wiktionary, the free dictionary.| - English Numbers - explanations, exercises and number generator (cardinal and ordinal numbers)

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South Carolina in the American Civil War |State of South Carolina| |Admission to Confederacy||February 4, 1861 (1st)| * 301,302 free * 402,406 slave |Forces supplied||23% of white population Total |Major garrisons/armories||Fort Sumter, Charleston Harbor| |Governor||Francis Pickens (1860-1862) |Senators||Robert Woodward Barnwell James Lawrence Orr |Restored to the Union||July 9, 1868| Part of a series on the |History of South Carolina| |South Carolina portal| American Civil War South Carolina was a site of a major political and military importance for the Confederacy during the American Civil War. The white population of the state strongly supported the institution of slavery long before the war, since the 18th century. Political leaders such as Democrats John Calhoun and Preston Brooks had inflamed regional and national passions in support of the institution, and for years before the eventual start of the Civil War in 1861, pro-slavery voices cried for secession. The Civil War began in South Carolina. On December 20, 1860, South Carolina, having the highest percentage of slaves of any U.S. state at 57% of its population enslaved and 46% of its families owning at least one slave, became the first state to declare that it had secessed from the Union. The first shots of the Civil War (January 9, 1861) were fired in Charleston by its Citadel cadets upon a U.S. civilian merchant ship, Star of the West, bringing supplies to the beleaguered U.S. garrison at Fort Sumter. The April 1861 bombardment of Fort Sumter by South Carolinian forces under the command of General Beauregard—the Confederacy did not yet have a functioning army—is commonly regarded as the beginning of the war. South Carolina was a source of troops for the Confederate army, and as the war progressed, also for the Union, as thousands of ex-slaves flocked to join the Union forces. The state also provided uniforms, textiles, food, and war material, as well as trained soldiers and leaders from The Citadel and other military schools. In contrast to most other Confederate states, South Carolina had a well-developed rail network linking all of its major cities without a break of gauge. Relatively free from Union occupation until the very end of the war, South Carolina hosted a number of prisoner of war camps. South Carolina also was the only Confederate state not to harbor pockets of anti-secessionist fervor strong enough to send large amounts of white men to fight for the Union, as every other state in the Confederacy did. Among the leading generals from the Palmetto State were Wade Hampton III, one of the Confederacy's leading cavalrymen, Maxcy Gregg, killed in action at Fredericksburg, Joseph B. Kershaw, whose South Carolina infantry brigade saw some of the hardest fighting of the Army of Northern Virginia and James Longstreet who served in that army under Robert E. Lee and in the Army of Tennessee under Gen. Braxton Bragg. For decades, South Carolinian political leaders had promoted regional passions with threats of nullification and secession in the name of southern states rights and protection of the interests of the slave power. Alfred P. Aldrich, a South Carolinian politician from Barnwell, stated that declaring secession would be necessary if a Republican candidate were to win the 1860 U.S. presidential election, stating that it was the only way for the state to preserve slavery and diminish the influence of the anti-slavery Republican Party, which, were its goals of abolition realized, would result in the "destruction of the South": If the Republican party with its platform of principles, the main feature of which is the abolition of slavery and, therefore, the destruction of the South, carries the country at the next Presidential election, shall we remain in the Union, or form a separate Confederacy? This is the great, grave issue. It is not who shall be President, it is not which party shall rule – it is a question of political and social existence.— Alfred P. Aldrich, In a January 1860 speech, South Carolinian congressman Laurence Massillon Keitt, summed up this view in an oratory condemning the Republican Party for its anti-slavery views, claiming that slavery was not morally wrong, but rather, justified: Later that year, in December, Keitt would state that South Carolina's declaring of secession was the direct result of slavery: On November 9, 1860 the South Carolina General Assembly passed a "Resolution to Call the Election of Abraham Lincoln as U.S. President a Hostile Act" and stated its intention to declare secession from the United States. In December 1860, amid the secession crisis, former South Carolinian congressman John McQueen wrote to a group of civic leaders in Richmond, Virginia, regarding the reasons as to why South Carolina was contemplating secession from the Union. In the letter, McQueen claimed that U.S. president-elect Abraham Lincoln supported equality and civil rights for African Americans as well as the abolition of slavery, and thus South Carolina, being opposed to such measures, was compelled to secede: I have never doubted what Virginia would do when the alternatives present themselves to her intelligent and gallant people, to choose between an association with her sisters and the dominion of a people, who have chosen their leader upon the single idea that the African is equal to the Anglo-Saxon, and with the purpose of placing our slaves on equality with ourselves and our friends of every condition! and if we of South Carolina have aided in your deliverance from tyranny and degradation, as you suppose, it will only the more assure us that we have performed our duty to ourselves and our sisters in taking the first decided step to preserve an inheritance left us by an ancestry whose spirit would forbid its being tarnished by assassins. We, of South Carolina, hope soon to great you in a Southern Confederacy, where white men shall rule our destinies, and from which we may transmit to our posterity the rights, privileges and honor left us by our ancestors. South Carolinian religious leader James Henley Thornwell also espoused a similar view to McQueen's, stating that slavery was justified under the Christian religion, and thus, those who viewed slavery as being immoral were opposed to Christianity: The parties in the conflict are not merely abolitionists and slaveholders. They are atheists, socialists, communists, red republicans, Jacobins on the one side, and friends of order and regulated freedom on the other. In one word, the world is the battleground – Christianity and Atheism the combatants; and the progress of humanity at stake. On November 10, 1860 the S.C. General Assembly called for a "Convention of the People of South Carolina" to consider secession. Delegates were to be elected on December 6. The secession convention convened in Columbia on December 17 and voted unanimously, 169-0, to declare secession from the United States. The convention then adjourned to Charleston to draft an ordinance of secession. When the ordinance was adopted on December 20, 1860, South Carolina became the first slave state in the south to declare that it had seceded from the United States. James Buchanan, the United States president, declared the ordinance illegal but did not act to stop it. A committee of the convention also drafted a Declaration of the Immediate Causes Which Induce and Justify the Secession of South Carolina which was adopted on December 24. The secession declaration stated the primary reasoning behind South Carolina's declaring of secession from the Union, which was described as: ...increasing hostility on the part of the non-slaveholding States to the Institution of Slavery ...— Declaration of the Immediate Causes Which Induce and Justify the Secession of South Carolina, (December 24, 1860). The declaration also claims that secession was declared as a result of the refusal of free states to enforce the Fugitive Slave Acts. Although the declaration does argue that secession is justified on the grounds of U.S. "encroachments upon the reserved rights of the States," the grievances that the declaration goes on to list are mainly concerned with the property of rights of slave holders. Broadly speaking, the declaration argues that the U.S. Constitution was framed to establish each State "as an equal" in the Union, with "separate control over its own institutions", such as "the right of property in slaves." We affirm that these ends for which this Government was instituted have been defeated, and the Government itself has been made destructive of them by the action of the non-slaveholding States. Those States have assumed the right of deciding upon the propriety of our domestic institutions; and have denied the rights of property established in fifteen of the States and recognized by the Constitution; they have denounced as sinful the institution of Slavery; they have permitted the open establishment among them of societies, whose avowed object is to disturb the peace and to eloign the property of the citizens of other States. They have encouraged and assisted thousands of our slaves to leave their homes; and those who remain, have been incited by emissaries, books and pictures to servile insurrection. A repeated concern is runaway slaves. The declaration argues that parts of the U.S. Constitution were specifically written to ensure the return of slaves who had escaped to other states, and quotes the 4th Article: "No person held to service or labor in one State, under the laws thereof, escaping into another, shall, in consequence of any law or regulation therein, be discharged from such service or labor, but shall be delivered up, on claim of the party to whom such service or labor may be due." The declaration goes on to state that this stipulation of the Constitution was so important to the original signers, "that without it that compact [the Constitution] would not have been made." Laws from the "General Government" upheld this stipulation "for many years," the declaration says, but "an increasing hostility on the part of the non-slaveholding States to the Institution of Slavery has led to a disregard of their obligations." Because the constitutional agreement had been "deliberately broken and disregarded by the non-slaveholding States," the consequence was that "South Carolina is released from her obligation" to be part of the Union. A further concern was Lincoln's recent election to the presidency, whom they claimed desired to see slavery on "the course of ultimate extinction": A geographical line has been drawn across the Union, and all the States north of that line have united in the election of a man to the high office of President of the United States whose opinions and purposes are hostile to slavery. He is to be entrusted with the administration of the Common Government, because he has declared that that "Government cannot endure permanently half slave, half free," and that the public mind must rest in the belief that Slavery is in the course of ultimate extinction. The South Carolinian secession declaration of December 1860 also channeled some elements from the U.S. Declaration of Independence from July 1776. However, the South Carolinian version omitted the phrases that "all men are created equal" and "that they are endowed by their Creator with certain unalienable Rights". Professor and historian Harry V. Jaffa noted these omissions as significant in his 2000 book, A New Birth of Freedom: Abraham Lincoln and the Coming of the Civil War: South Carolina cites, loosely, but with substantial accuracy, some of the language of the original Declaration. That Declaration does say that it is the right of the people to abolish any form of government that becomes destructive of the ends for which it was established. But South Carolina does not repeat the preceding language in the earlier document: 'We hold these truths to be self-evident, that all men are created equal'... The following day, on December 25, a South Carolinian convention delivered an "Address to the Slaveholding States": We prefer, however, our system of industry, by which labor and capital are identified in interest, and capital, therefore, protects labor–by which our population doubles every twenty years–by which starvation is unknown, and abundance crowns the land–by which order is preserved by unpaid police, and the most fertile regions of the world, where the white man cannot labor, are brought into usefulness by the labor of the African, and the whole world is blessed by our own productions. ... We ask you to join us, in forming a Confederacy of Slaveholding States.— Convention of South Carolina, Address of the people of South Carolina to the people of the Slaveholding States, (December 25, 1860) "Slavery, not states' rights, birthed the Civil War," argues sociologist James W. Loewen. Writing of South Carolina's Declaration of Secession, Loewen writes that South Carolina was further upset that New York no longer allowed "slavery transit." In the past, if Charleston gentry wanted to spend August in the Hamptons, they could bring their cook along. No longer — and South Carolina's delegates were outraged. In addition, they objected that New England states let black men vote and tolerated abolitionist societies. According to South Carolina, states should not have the right to let their citizens assemble and speak freely when what they said threatened slavery. Other seceding states echoed South Carolina. "Our position is thoroughly identified with the institution of slavery — the greatest material interest of the world," proclaimed Mississippi in its own secession declaration, passed Jan. 9, 1861. "Its labor supplies the product which constitutes by far the largest and most important portions of the commerce of the earth. . . . A blow at slavery is a blow at commerce and civilization." The state adopted the palmetto flag as its banner, a slightly modified version of which is used as its current state flag. South Carolina after secession was frequently called the "Palmetto Republic". After South Carolina declared its secession, former congressman James L. Petigru famously remarked, "South Carolina is too small for a republic and too large for an insane asylum." Soon afterwards, South Carolina began preparing for a presumed U.S. military response while working to convince other southern states to secede as well and join in a confederacy of southern states. On February 4, 1861, in Montgomery, Alabama, a convention consisting of delegates from South Carolina, Florida, Alabama, Mississippi, Georgia, and Louisiana met to form a new constitution and government modeled on that of the United States. On February 8, 1861, South Carolina officially joined the Confederacy. According to one South Carolinian newspaper editor: The South is now in the formation of a Slave Republic...— L.W. Spratt, The Philosophy of Secession: A Southern View, (February 13, 1861). South Carolina's declaring of secession was supported by the state's religious figures, who claimed that it was consistent with their religion: The triumphs of Christianity rest this very hour upon slavery; and slavery depends on the triumphs of the South... This war is the servant of slavery.— John T. Wightman, The Glory of God, the Defence of the South, (1861). American Civil War Six days after secession, on the day after Christmas, Major Robert Anderson, commander of the U.S. troops in Charleston, withdrew his men to the island fortress of Fort Sumter in Charleston Harbor. South Carolina militia swarmed over the abandoned mainland batteries and trained their guns on the island. Sumter was the key position for preventing a naval attack upon Charleston, so secessionists were determined not to allow U.S. forces to remain there indefinitely. More importantly, South Carolina's claim of independence would look empty if U.S. forces controlled its largest harbor. On January 9, 1861, the U.S. ship Star of the West approached to resupply the fort. Cadets from The Citadel, The Military College of South Carolina fired upon the Star of the West, striking the ship three times and causing it to retreat back to New York. Mississippi declared its secession several weeks after South Carolina, and five other states of the lower South soon followed. Both the outgoing Buchanan administration and President-elect Lincoln denied that any state had a right to secede. On February 4, a congress of the seven seceding states met in Montgomery, Alabama, and approved a new constitution for the Confederate States of America. South Carolina entered the Confederacy on February 8, 1861, fewer than six weeks after declaring itself the independent State of South Carolina. Upper Southern slave states such as Virginia and North Carolina, which had initially voted against secession, called a peace conference, to little effect. Meanwhile, Virginian orator Roger Pryor barreled into Charleston and proclaimed that the only way to get his state to join the Confederacy was for South Carolina to instigate war with the United States. The obvious place to start was right in the midst of Charleston Harbor. On April 10, the Mercury reprinted stories from New York papers that told of a naval expedition that had been sent southward toward Charleston. Lincoln advised the governor of South Carolina that the ships were sent to resupply the fort, not to reinforce it. The Carolinians could no longer wait if they hoped to take the fort before the U.S. Navy arrived. About 6,000 men were stationed around the rim of the harbor, ready to take on the 60 men in Fort Sumter. At 4:30 a.m. on April 12, after two days of intense negotiations, and with Union ships approaching the harbor, the firing began. Students from The Citadel were among those firing the first shots of the war, though Edmund Ruffin is usually credited with firing the first shot. Thirty-four hours later, Anderson's men raised the white flag and were allowed to leave the fort with colors flying and drums beating, saluting the U.S. flag with a 50-gun salute before taking it down. During this salute, one of the guns exploded, killing a young soldier—the only casualty of the bombardment and the first casualty of the war. In December 1861, South Carolina received $100,000 from Georgia after a disastrous fire in Charleston. The war ends The Confederacy was at a disadvantage in number, weaponry, and maritime skills, as few southerners were sailors before the war. Union ships sailed south and blocked off one port after another. As early as November, Union troops occupied the Sea Islands in the Beaufort area, establishing an important base for the men and ships who would obstruct the ports at Charleston and Savannah. When the plantation owners, many of which had already gone off with the Confederate army elsewhere, fled the area, the Sea Island slaves became the first "freedmen" of the war, and the Sea Islands became the laboratory for Union plans to educate the African Americans for their eventual role as full American citizens Despite South Carolina's important role in the start of the war, and a long unsuccessful attempt to take Charleston from 1863 onward, few military engagements occurred within the state's borders until 1865, when Sherman's Army, having already completed its March to the Sea in Savannah, marched to Columbia and leveled most of the town, as well as a number of towns along the way and afterward. South Carolina lost 12,922 men to the war, 23% of its male white population of fighting age, and the highest percentage of any state in the nation. Sherman's 1865 march through the Carolinas resulted in the burning of Columbia and numerous other towns. The destruction his troops wrought upon South Carolina was even worse than in Georgia, because many of his men bore a particular grudge against the state and its citizens, who they blamed for starting the war. One of Sherman's men declared, "Here is where treason began and, by God, here is where it shall end!" Poverty would mark the state for generations to come. In January 1865, the Charleston Courier newspaper condemned suggestions that the Confederacy abandon slavery were it to help in gaining independence, stating that such suggestions were "folly": To talk of maintaining our independence while we abolish slavery is simply to talk folly.— Courier, (January 24, 1865) On February 21, 1865, with the Confederate forces finally evacuated from Charleston, the black 54th Massachusetts Regiment marched through the city. At a ceremony at which the U.S. flag was once again raised over Fort Sumter, former fort commander Robert Anderson was joined on the platform by two men: African American Union hero Robert Smalls and the son of Denmark Vesey. Battles in South Carolina - Battle of Fort Sumter - Battle of Port Royal - Battle of Secessionville - Battle of Simmon's Bluff - First Battle of Charleston Harbor - Second Battle of Charleston Harbor - Second Battle of Fort Sumter - First Battle of Fort Wagner - Battle of Grimball's Landing - Second Battle of Fort Wagner (Morris Island) - Battle of Honey Hill - Battle of Tulifinny - Battle of Rivers' Bridge - Battle of Anderson County - Battle of Brattonsville - Battle of Broxton's Bridge - Battle of Cheraw - Battle of Gamble's Hotel (The Columns) - Battle of Aiken - Charleston, South Carolina in the American Civil War - Confederate States of America - animated map of state secession and confederacy - List of South Carolina Confederate Civil War units - List of South Carolina Union Civil War units - Military history of African Americans in the American Civil War - Origins of the American Civil War - Slaves and the American Civil War - "Results from the 1860 Census". 1860 United States Census. 1860. Retrieved June 4, 2004. - Hall, Andy (December 22, 2013). "Not Surprising, Part Deux". Dead Confederates: A Civil War Era Blog. The states with the largest proportions of slaves and slave-holders seceded earliest. - Channing, Steven. Crisis of Fear. pp. 141–142. Retrieved September 6, 2015. - Keitt, Lawrence M. (January 25, 1860). Congressman from South Carolina, in a speech to the House. Taken from a photocopy of the Congressional Globe, supplied by Steve Miller. The anti-slavery party contends that slavery is wrong in itself, and the Government is a consolidated national democracy. We of the South contend that slavery is right, and that this is a confederate Republic of sovereign States. - "The Charleston Courier". Charleston, South Carolina. December 22, 1860. Retrieved September 6, 2015. - "Resolution to Call the Election of Abraham Lincoln as U.S. President a Hostile Act and to Communicate to Other Southern States South Carolina's Desire to Secede from the Union." 9 November 1860. Resolutions of the General Assembly, 1779-1879. S165018. South Carolina Department of Archives and History, Columbia, S.C. - McQueen, John (December 24, 1860). "Correspondence to T. T. Cropper and J. R. Crenshaw". Washington, D.C. Retrieved March 25, 2015. - Rhea, Gordon (January 25, 2011). "Why Non-Slaveholding Southerners Fought". Civil War Trust. Civil War Trust. Retrieved March 21, 2011. - Cauthen, Charles Edward; Power, J. Tracy. South Carolina goes to war, 1860-1865. Columbia, SC: University of South Carolina Press, 2005. Originally published: Chapel Hill, NC: University of North Carolina Press, 1950. ISBN 978-1-57003-560-9. p. 60. - "'Declaration of the Immediate Causes Which Induce and Justify the Secession of South Carolina from the Federal Union,' 24 December 1860". Teaching American History in South Carolina Project. 2009. Retrieved November 18, 2012. - Jaffa, Harry V. (2000). A New Birth of Freedom: Abraham Lincoln and the Coming of the Civil War. Rowman & Littlefield Publishers. p. 231. - State of South Carolina (December 25, 1860). "Address of the people of South Carolina to the people of the Slaveholding States of the United States". Retrieved March 27, 2015. - Loewen, James (2011). "Five Myths About Why the South Seceded". Washington Post. - Edgar, Walter. South Carolina: A History, Columbia, SC: University of South Carolina Press:1998. ISBN 978-1-57003-255-4. p. 619 - Cauthen, Charles Edward; Power, J. Tracy. South Carolina goes to war, 1860-1865. Columbia, SC: University of South Carolina Press, 2005. Originally published: Chapel Hill, NC: University of North Carolina Press, 1950. ISBN 978-1-57003-560-9. p. 79. - Burger, Ken (February 13, 2010). "Too large to be an asylum". The Post and Courier (Charleston, South Carolina: Evening Post Publishing Co). Retrieved April 22, 2010. Paragraph 4 - Lee, Jr., Charles Robert. The Confederate Constitutions. Chapel Hill, NC: The University of North Carolina Press, 1963, 60. - Spratt, L.W. (February 13, 1861). "THE PHILOSOPHY OF SECESSION: A SOUTHERN VIEW". South Carolina. Retrieved September 13, 2015. Presented in a Letter addressed to the Hon. Mr. Perkins of Louisiana, in criticism on the Provisional Constitution adopted by the Southern Congress at Montgomery, Alabama, BY THE HON. L. W. SPRATT, Editor of the Charleston Mercury, 13th February, 1861. - Wightman, John T. (1861). "The Glory of God, the Defence of the South". Yorkville, South Carolina. Retrieved September 8, 2015. - McPherson, James M. This Mighty Scourge: Perspectives on the Civil War. Oxford University Press, 2009 - "Courier". Charleston. January 24, 1865. Retrieved September 8, 2015. - Burger, Ken (February 13, 2010). "Too large to be an asylum". The Post and Courier (Charleston, South Carolina: Evening Post Publishing Co). Retrieved April 22, 2010.. - Cauthen, Charles Edward; Power, J. Tracy. South Carolina goes to war, 1860-1865. Columbia, SC: University of South Carolina Press, 2005. Originally published: Chapel Hill, NC: University of North Carolina Press, 1950. ISBN 978-1-57003-560-9. - Edgar, Walter. South Carolina: A History, Columbia, SC: University of South Carolina Press:1998. ISBN 978-1-57003-255-4. - Rogers Jr. George C. and C. James Taylor. A South Carolina Chronology, 1497-1992 2nd Ed. (1994) - Wallace, David Duncan. South Carolina: A Short History, 1520-1948 (1951) standard scholarly history - WPA. South Carolina: A Guide to the Palmetto State (1941) - Wright, Louis B. South Carolina: A Bicentennial History' (1976) - Declaration of the Immediate Causes Which Induce and Justify the Secession of South Carolina from the Federal Union |Wikiquote has quotations related to: American Civil War| |Wikiquote has quotations related to: Confederate States of America|

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At a Glance - Quadratic Inequalities Remember back when we looked at linear inequalities, that we said, "we promise we'll try to make your brain hurt more later?" Well grab an ice pack and strap in, because now we're going to look at quadratic inequalities. Solve x2 – 5x < -4. When we solve an inequality, what we want are all of the values of x that make the statement true. So our answers won't be single values, but large, sweeping regions of number space. You can then fence off those regions and raise cows on them. We here at Shmoop love the equal sign. It's a good thing that the first step of solving an inequality is to pretend that the inequality is an equal sign. Set the equation "equal" to zero, and then solve to find the roots of the equation. They'll come in handy in a moment. x2 – 5x < -4 x2 – 5x + 4 < 0 (x – 1)(x – 4) < 0 Okay, our roots are x = 1 and x = 4. So what? Take a look at the graph of this equation. A parabola is a smooth, continuous curve. The only places that it can possibly change sign (from above zero to below, or vice versa) are at the roots. We'll use this to help us find our solutions. Hey, waitjustaminutehere! Couldn't we just graph the equation and solve it visually? We could, but there are two good reasons not to. First, it will often be just as or more difficult to graph the equation than it will be to solve it the other way (see: the next sample problem). Second, we can also use this technique to solve all kinds of polynomial inequalities, not just quadratic ones (see: the sample problem after that). Anyway, back to solving. We'll now set up our roots on a number line, like so. We now have three regions fenced off. We need to pick a point from each region to check what whether it is positive or negative within that region. Those regions that are negative will be our solutions. Afterwards we'll put our cows in the positive regions, to boost cow morale. All the values of x between 1 and 4 will cause the equation to be negative. So our solutions are 1 < x < 4. If you look back at the graph of the equation, you'll see that this is the region where it dips down below zero. Solve -2x2 ≤ 6x + 1. We again start off by getting all our stuff on one side of the equation, leaving a big fat zero on the other side of the inequality. 2x2 + 6x + 1 ≥ 0 Last time we had a nice, factorable equation to work with. Not this time, bucko. Now we need to use the quadratic formula to find our roots. Our calculator tells us that these are x = -0.177 and x = -2.823. Now let's set up the number line and check the signs of each region. We want the regions that are greater than zero, so the solutions are -∞ < x ≤ -2.823 and -0.177 ≤ x < ∞ So, is your brain starting to hurt? This next one is the last problem here, so stick with it a little longer. Solve (x + 3)2(3x2 – 6) < 0. This equation definitely isn't quadratic, but the method for finding the solutions is the same. It is even factored already, making things easier than they could have been. The roots are x = -3, , and . The number line looks like The equation is less than zero when . When working with polynomials larger than the quadratics there can be more than two roots, and we need to check the sign of every region. For every inequality, the sign won't necessarily follow a predictable pattern from one root to the next. It's as random as a corn syrup Huckleberry sauce. …How did you like that? Is your brain throbbing with knowledge? That doesn't sound pleasant, but at least you're a bit smarter from the experience. Solve the inequality 3x2 – 8x + 4 > 0. Solve the inequality -x2 – 4x > 3. Solve the inequality -2x2 + 8x + 8 ≤ 0.

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We live in a galaxy known as the Milky Way – a vast conglomeration of 300 billion stars, planets whizzing around them, and clouds of gas and dust floating in between. Though it has long been known that the Milky Way and its orbiting companion Andromeda are the dominant members of a small group of galaxies, the Local Group, which is about 3 million light years across, much less was known about our immediate neighbourhood in the universe. Now, a new paper by York University Physics & Astronomy Professor Marshall McCall, published today in the Monthly Notices of the Royal Astronomical Society, maps out bright galaxies within 35-million light years of the Earth, offering up an expanded picture of what lies beyond our doorstep. "All bright galaxies within 20 million light years, including us, are organized in a 'Local Sheet' 34-million light years across and only 1.5-million light years thick," says McCall. "The Milky Way and Andromeda are encircled by twelve large galaxies arranged in a ring about 24-million light years across – this 'Council of Giants' stands in gravitational judgment of the Local Group by restricting its range of influence." McCall says twelve of the fourteen giants in the Local Sheet, including the Milky Way and Andromeda, are "spiral galaxies" which have highly flattened disks in which stars are forming. The remaining two are more puffy "elliptical galaxies", whose stellar bulks were laid down long ago. Intriguingly, the two ellipticals sit on opposite sides of the Council. Winds expelled in the earliest phases of their development might have shepherded gas towards the Local Group, thereby helping to build the disks of the Milky Way and Andromeda. McCall also examined how galaxies in the Council are spinning. He comments: "Thinking of a galaxy as a screw in a piece of wood, the direction of spin can be described as the direction the screw would move (in or out) if it were turned the same way as the galaxy rotates. Unexpectedly, the spin directions of Council giants are arranged around a small circle on the sky. This unusual alignment might have been set up by gravitational torques imposed by the Milky Way and Andromeda when the universe was smaller." The boundary defined by the Council has led to insights about the conditions which led to the formation of the Milky Way. Most important, only a very small enhancement in the density of matter in the universe appears to have been required to produce the Local Group. To arrive at such an orderly arrangement as the Local Sheet and its Council, it seems that nearby galaxies must have developed within a pre-existing sheet-like foundation comprised primarily of dark matter. "Recent surveys of the more distant universe have revealed that galaxies lie in sheets and filaments with large regions of empty space called voids in between," says McCall. "The geometry is like that of a sponge. What the new map reveals is that structure akin to that seen on large scales extends down to the smallest." York University is helping to shape the global thinkers and thinking that will define tomorrow. York U's unwavering commitment to excellence reflects a rich diversity of perspectives and a strong sense of social responsibility that sets us apart. A York U degree empowers graduates to thrive in the world and achieve their life goals through a rigorous academic foundation balanced by real-world experiential education. As a globally recognized research centre, York U is fully engaged in the critical discussions that lead to innovative solutions to the most pressing local and global social challenges. York U's 11 faculties and 27 research centres are thinking bigger, broader and more globally, partnering with 288 leading universities worldwide. York U's community is strong − 55,000 students, 7,000 faculty and staff, and more than 250,000 alumni. Media Contact: Robin Heron, Media Relations, York University, 416 736 2100 x22097/ [email protected] Robin Heron | EurekAlert! LIGO confirms RIT's breakthrough prediction of gravitational waves 12.02.2016 | Rochester Institute of Technology Milestone in physics: gravitational waves detected with the laser system from LZH 12.02.2016 | Laser Zentrum Hannover e.V. Today, plants and microorganisms are heavily used for the production of medicinal products. The production of biopharmaceuticals in plants, also referred to as “Molecular Pharming”, represents a continuously growing field of plant biotechnology. Preferred host organisms include yeast and crop plants, such as maize and potato – plants with high demands. With the help of a special algal strain, the research team of Prof. Ralph Bock at the Max Planck Institute of Molecular Plant Physiology in Potsdam strives to develop a more efficient and resource-saving system for the production of medicines and vaccines. They tested its practicality by synthesizing a component of a potential AIDS vaccine. The use of plants and microorganisms to produce pharmaceuticals is nothing new. In 1982, bacteria were genetically modified to produce human insulin, a drug... Atomic clock experts from the Physikalisch-Technische Bundesanstalt (PTB) are the first research group in the world to have built an optical single-ion clock which attains an accuracy which had only been predicted theoretically so far. Their optical ytterbium clock achieved a relative systematic measurement uncertainty of 3 E-18. The results have been published in the current issue of the scientific journal "Physical Review Letters". Atomic clock experts from the Physikalisch-Technische Bundesanstalt (PTB) are the first research group in the world to have built an optical single-ion clock... The University of Würzburg has two new space projects in the pipeline which are concerned with the observation of planets and autonomous fault correction aboard satellites. The German Federal Ministry of Economic Affairs and Energy funds the projects with around 1.6 million euros. Detecting tornadoes that sweep across Mars. Discovering meteors that fall to Earth. Investigating strange lightning that flashes from Earth's atmosphere into... Physicists from Saarland University and the ESPCI in Paris have shown how liquids on solid surfaces can be made to slide over the surface a bit like a bobsleigh on ice. The key is to apply a coating at the boundary between the liquid and the surface that induces the liquid to slip. This results in an increase in the average flow velocity of the liquid and its throughput. This was demonstrated by studying the behaviour of droplets on surfaces with different coatings as they evolved into the equilibrium state. The results could prove useful in optimizing industrial processes, such as the extrusion of plastics. The study has been published in the respected academic journal PNAS (Proceedings of the National Academy of Sciences of the United States of America). Exceeding critical temperature limits in the Southern Ocean may cause the collapse of ice sheets and a sharp rise in sea levels A future warming of the Southern Ocean caused by rising greenhouse gas concentrations in the atmosphere may severely disrupt the stability of the West... 12.02.2016 | Event News 09.02.2016 | Event News 02.02.2016 | Event News 12.02.2016 | Physics and Astronomy 12.02.2016 | Life Sciences 12.02.2016 | Medical Engineering

http://www.innovations-report.com/html/reports/physics-astronomy/york-u-astronomer-maps-out-earth-s-place-in-the-universe-among-council-of-giants.html

The bands of color on a resistor are a code that indicates the magnitude of the resistance of the resistor. There are four color bands identified by letter: A, B, C, and D, with a gap between the C and D bands so that you know which end is A. This particular resistor has a red A band, blue B band, green C band, and gold D band, but the bands can be different colors on different resistors. Based on the colors of the bands, it is possible to identify the type of resistor. the A and B bands represent significant digits; red is 2 and blue is 6. The C band indicates the multiplier, and green indicates 105. These three together indicate that this particular resistor is a 26,000 Ohm resistor. Finally, the D band indicates the tolerance, in this case 5%, as shown by the gold band. These terms will be explained over the course of this lesson. Resistance and Ohm’s Law When a potential difference is placed across a metal wire, a large current will flow through the wire. If the same potential difference is placed across a glass rod, almost no current will flow. The property that determines how much current will flow is called the resistance. Resistance is measured by finding the ratio of potential difference, V, to current flow, I. When given in the form V=IR, this formula is known as Ohm's Law, after the man that discovered the relationship. The units of resistance can be determined using the units of the other terms in the equation, namely that the potential difference is in volts (J/C) and current in amperes (C/s): The units for resistance have been given the name ohms and the abbreviation is the Greek letter omega, Ω. 1.00 Ω is the resistance that will allow 1.00 ampere of current to flow through the resistor when the potential difference is 1.00 volt. Most conductors have a constant resistance regardless of the potential difference; these are said to obey Ohm's Law. There are two ways to control the current in a circuit. Since the current is directly proportional to the potential difference and inversely proportional to the resistance, you can increase the current in a circuit by increasing the potential or by decreasing the resistance. Example Problem: A 50.0 V battery maintains current through a 20.0 Ω resistor. What is the current through the resistor? Solution: I=VR=50.0 V20.0 Ω=2.50 amps - Resistance is the property that determines the amount of current flow through a particular material. V=IR is known as Ohm’s Law. - The unit for resistance is the ohm, and it has the abbreviation Ω. The following video covers Ohm's Law. Use this resource to answer the questions that follow. - What happens to current flow when voltage is increased? - What happens to current flow when resistance is increased? This website contains instruction and guided practice for Ohm’s Law. - If the potential stays the same and the resistance decreases, what happens to the current? - stay the same - If the resistance stays the same and the potential increases, what happens to the current? - stay the same - How much current can be pushed through a 30.0 Ω resistor by a 12.0 V battery? - What voltage is required to push 4.00 A of current through a 32.0 Ω resistor? - If a 6.00 volt battery will produce 0.300 A of current in a circuit, what is the resistance in the circuit?

http://www.ck12.org/book/CK-12-Physics-Concepts---Intermediate/r19/section/18.2/

|Part of a series on| A communist party is a political party that advocates the application of the social and economic principles of communism through state policy. The name originates from the 1848 tract Manifesto of the Communist Party by Karl Marx and Friedrich Engels. According to Leninist theory, a Communist party is the vanguard party of the working class (Proletariat), whether ruling or non-ruling, but when such a party is in power in a specific country, the party is said to be the highest authority of the dictatorship of the proletariat. Vladimir Lenin's theories on the role of a Communist party were developed as the early 20th-century Russian social democracy divided into Bolshevik (meaning "of the majority") and Menshevik (meaning "of the minority") factions. Lenin, leader of the Bolsheviks, argued that a revolutionary party should be a small vanguard party with a centralized political command and a strict cadre policy; the Menshevik faction, however, argued that the party should be a broad-based mass movement. The Bolshevik party, which eventually became the Communist Party of the Soviet Union, took power in Russia after the October Revolution in 1917. With the creation of the Communist International, the Leninist concept of party building was copied by emerging Communist parties worldwide. As the membership of a Communist party was to be limited to active cadres in Lenin's theory, there was a need for networks of separate organizations to mobilize mass support for the party. Typically, Communist parties have built up various front organizations whose membership is often open to non-Communists. In many countries the single most important front organization of the Communist parties has been its youth wing. During the time of the Communist International, the youth leagues were explicit Communist organizations, using the name 'Young Communist League'. Later the youth league concept was broadened in many countries, and names like 'Democratic Youth League' were adopted. Some trade unions, student, women's, grifters, peasant's and cultural organizations have been connected to Communist parties. Traditionally, these mass organizations were often politically subordinated to the political leadership of the party. However, in many contemporary cases mass organizations founded by communists have acquired a certain degree of independence. In some cases mass organizations have outlived the Communist parties in question. At the international level, the Communist International organized various international front organizations (linking national mass organizations with each other), such as the Young Communist International, Profintern, Krestintern, International Red Aid, Sportintern, etc.. These organizations were dissolved in the process of deconstruction of the Communist International. After the Second World War new international coordination bodies were created, such as the World Federation of Democratic Youth, International Union of Students, World Federation of Trade Unions, Women's International Democratic Federation and the World Peace Council. Historically, in countries where Communist Parties were struggling to attain state power, the formation of wartime alliances with non-Communist parties and wartime groups was enacted (such as the National Liberation Front of Albania). Upon attaining state power these Fronts were often transformed into nominal (and usually electoral) "National" or "Fatherland" Fronts in which non-communist parties and organizations were given token representation (a practice known as Blockpartei), the most popular examples of these being the National Front of East Germany (as a historical example) and the United Front of the People's Republic of China (as a modern-day example). Other times the formation of such Fronts were undertaken without the participation of other parties, such as the Socialist Alliance of Working People of Yugoslavia and the National Front of Afghanistan, though the purpose was the same: to promote the Communist Party line to generally non-communist audiences and to mobilize them to carry out tasks within the country under the aegis of the Front. A uniform naming scheme for Communist parties was adopted by the Communist International. All parties were required to use the name 'Communist Party of (name of country)', resulting in separate communist parties in some countries operating using (largely) homonymous party names (e.g. in India). Today, there are plenty of cases where the old sections of the Communist International have retained those names. In other cases names have been changed. Common causes for the shift in naming were either moves to avoid state repression or as measures to indicate a broader political acceptance. A typical example of the latter was the renaming of various East European Communist parties after the Second World War, as a result of mergers with the local Social Democratic parties. New names in the post-war era included "Socialist Party", "Socialist Unity Party", "Popular Party", "Workers' Party" and "Party of Labour". The naming conventions of Communist parties became more diverse as the international Communist movement was fragmented due to the Sino-Soviet split in the 1960s. Those who sided with China and Albania in their criticism of the Soviet leadership, often added words like 'Revolutionary' or 'Marxist-Leninist' to distinguish themselves from the pro-Soviet parties. |Wikimedia Commons has media related to Communist Parties.| - Harper, Douglas. "communism". Online Etymology Dictionary. Retrieved 2008-08-27. - "The Chinese Communist Party". Council on Foreign Relations. Retrieved 25 February 2015. - China's communist party members near 78 mln - Domeinnaam niet ingeschakeld - "Nieuws". PVDA. Retrieved 25 February 2015. - One such example is the Swiss Party of Labour, which was founded in 1944 to substitute the illegalized Communist Party of Switzerland. - Such mergers occurred in East Germany (Socialist Unity Party of Germany), Hungary (Hungarian Working People's Party), Poland (Polish United Workers Party) and Romania (Romanian Workers Party).

https://en.wikipedia.org/wiki/Communist_party

This sentence diagramming worksheet focuses on adjectives, adverbs and articles. Diagramming Sentences Worksheets A sentence diagram is a way to graphically represent the structure of a sentence, showing how words in a sentence function and relate to each other. The printable practice worksheets below provide supplemental help in learning the basic concepts of sentence diagramming. Feel free to print them off and duplicate for home or classroom use. It’s all about conjunctions in this diagramming sentences worksheet! Time to diagram sentences with direct and indirect objects! In this diagramming sentences worksheet, your student will practice with prepositional phrases. There are a lot of compounds in this sentence diagram worksheet! A helpful sentence diagramming guide for students to use at home or in the classroom. Now it’s time to practice diagramming sentences! Here’s a practice worksheet for your beginning sentence diagrammer that covers the subject and predicate. If you’re looking for a basic sentence diagramming worksheet, this is it! This worksheet focuses on diagramming complex sentences. Compound predicates are featured in this worksheet on diagramming sentences. Let’s diagram some compound sentences! In this worksheet your student will diagram sentences with compound subjects. This worksheet helps your student understand how to diagram helping verbs in a sentence. This activity provides students practice diagramming infinitives. This activity provides students practice diagramming intensive pronouns. This printable activity provides students practice diagramming interjections. What could be better than a sentence diagram worksheet on interrogatives? Object complements are the main attraction in this sentence diagramming worksheet. This activity provides students practice diagramming reflexive pronouns.

http://www.k12reader.com/subject/grammar/sentence-structure/diagramming-sentences/

- The Life of a Glacier - About Glaciers - Glacier Photo Gallery - Science and Data Resources - Further Reading - How to Cite What types of glaciers are there? These glaciers develop in high mountainous regions, often flowing out of icefields that span several peaks or even a mountain range. The largest mountain glaciers are found in Arctic Canada, Alaska, the Andes in South America, and the Himalaya in Asia. Commonly originating from mountain glaciers or icefields, these glaciers spill down valleys, looking much like giant tongues. Valley glaciers may be very long, often flowing down beyond the snow line, sometimes reaching sea level. As the name implies, these are valley glaciers that flow far enough to reach out into the sea. Tidewater glaciers are responsible for calving numerous small icebergs, which although not as imposing as Antarctic icebergs, can still pose problems for shipping lanes. Piedmont glaciers occur when steep valley glaciers spill into relatively flat plains, where they spread out into bulb-like lobes. Malaspina Glacier in Alaska is one of the most famous examples of this type of glacier, and is the largest piedmont glacier in the world. Spilling out of the Seward Icefield, Malaspina Glacier covers about 3,900 square kilometers (1,500 square miles) as it spreads across the coastal plain. When a major valley glacier system retreats and thins, sometimes the tributary glaciers are left in smaller valleys high above the shrunken central glacier surface. These are called hanging glaciers. If the entire system has melted and disappeared, the empty high valleys are called hanging valleys. These small, steep glaciers cling to high mountainsides. Like cirque glaciers, they are often wider than they are long. Ice aprons are common in the Alps and in New Zealand, where they often cause avalanches due to the steep inclines they occupy. Rock glaciers sometimes form when slow-moving glacial ice is covered by debris. They are often found in steep-sided valleys, where rocks and soil fall from the valley walls onto the ice. Rock glaciers may also form when frozen ground creeps downslope. Ice shelves occur when ice sheets extend over the sea and float on the water. They range from a few hundred meters to over 1 kilometer (0.62 mile) in thickness. Ice shelves surround most of the Antarctic continent. Ice caps are miniature ice sheets, covering less than 50,000 square kilometers (19,305 square miles). They form primarily in polar and sub-polar regions that are relatively flat and high in elevation. Ice streams are large ribbon-like glaciers set within an ice sheet—they are bordered by ice that is flowing more slowly, rather than by rock outcrop or mountain ranges. These huge masses of flowing ice are often very sensitive to changes such as the loss of ice shelves at their terminus or changing amounts of water flowing beneath them. The Antarctic ice sheet has many ice streams. Found now only in Antarctica and Greenland, ice sheets are enormous continental masses of glacial ice and snow expanding over 50,000 square kilometers (19,305 square miles). The ice sheet on Antarctica is over 4.7 kilometers (3 miles) thick in some areas, covering nearly all of the land features except the Transantarctic Mountains, which protrude above the ice. Another example is the Greenland Ice Sheet. In the past ice ages, huge ice sheets also covered most of Canada (the Laurentide Ice Sheet) and Scandinavia (the Scandinavian Ice Sheet), but these have now disappeared, leaving only a few ice caps and mountain glaciers behind. NSIDC's Glacier Glossary - Search and browse terms related to glaciers in NSIDC's comprehensive cryospheric glossary. NSIDC Glacier Photograph Collection - NSIDC archives a Glacier Photograph Collection of historical photos, which includes both aerial and terrestrial photos for the 1880s to 1975. The photos are primarily of Alaskan glaciers, but coverage also includes the Pacific Northwest and Europe.

https://nsidc.org/cryosphere/glaciers/questions/types.html

the The PhET Project and This PhET "Gold Star Winner" is an instructional unit on the topic of Waves, created by a high school teacher. It was designed to be used with interactive simulations developed by PhET, the Physics Education Technology project. Included are detailed lessons for integrating labs, simulations, demonstrations, and concept questions to introduce students to properties and behaviors of waves. Specific topics include frequency and wavelength, sound, the wave nature of light, geometric optics, resonance, wave interference, Doppler Effect, refraction, thin lenses, wave addition, and more. Activities are aligned to AAAS Benchmarks. Editor's Note:This could be a very useful resource for teachers in grades 8-12, allowing them to quickly customize a module which meets new content standards on Waves, outlined in the NextGen Science Framework. In its entirety, the unit is 4 weeks in duration. However, teachers of Grades 8-9 physical science could easily pull out 4-5 lessons addressing fundamental wave properties and the basics of refraction/reflection. All lessons include objectives and teaching tips, plus "clicker" or warm-up questions, worksheets, and unit tests with answer keys. Metadata instance created April 17, 2008 by Caroline Hall October 1, 2012 by Caroline Hall Last Update when Cataloged: March 31, 2008 AAAS Benchmark Alignments (2008 Version) 4. The Physical Setting 6-8: 4F/M2. Something can be "seen" when light waves emitted or reflected by it enter the eye—just as something can be "heard" when sound waves from it enter the ear. 6-8: 4F/M4. Vibrations in materials set up wavelike disturbances that spread away from the source. Sound and earthquake waves are examples. These and other waves move at different speeds in different materials. 6-8: 4F/M5. Human eyes respond to only a narrow range of wavelengths of electromagnetic waves-visible light. Differences of wavelength within that range are perceived as differences of color. 6-8: 4F/M6. Light acts like a wave in many ways. And waves can explain how light behaves. 6-8: 4F/M7. Wave behavior can be described in terms of how fast the disturbance spreads, and in terms of the distance between successive peaks of the disturbance (the wavelength). 9-12: 4F/H5ab. The observed wavelength of a wave depends upon the relative motion of the source and the observer. If either is moving toward the other, the observed wavelength is shorter; if either is moving away, the wavelength is longer. 9-12: 4F/H6ab. Waves can superpose on one another, bend around corners, reflect off surfaces, be absorbed by materials they enter, and change direction when entering a new material. All these effects vary with wavelength. 9-12: 4F/H6c. The energy of waves (like any form of energy) can be changed into other forms of energy. 11. Common Themes 6-8: 11B/M4. Simulations are often useful in modeling events and processes. 9-12: 11B/H3. The usefulness of a model can be tested by comparing its predictions to actual observations in the real world. But a close match does not necessarily mean that other models would not work equally well or better. 6-8: 11D/M3. Natural phenomena often involve sizes, durations, and speeds that are extremely small or extremely large. These phenomena may be difficult to appreciate because they involve magnitudes far outside human experience. Common Core State Standards for Mathematics Alignments High School — Functions (9-12) Interpreting Functions (9-12) F-IF.4 For a function that models a relationship between two quantities, interpret key features of graphs and tables in terms of the quantities, and sketch graphs showing key features given a verbal description of the relationship.? F-IF.5 Relate the domain of a function to its graph and, where applicable, to the quantitative relationship it describes.? F-IF.6 Calculate and interpret the average rate of change of a function (presented symbolically or as a table) over a specified interval. Estimate the rate of change from a graph. F-IF.7.a Graph linear and quadratic functions and show intercepts, maxima, and minima. Building Functions (9-12) F-BF.3 Identify the effect on the graph of replacing f(x) by f(x) + k, k f(x), f(kx), and f(x + k) for specific values of k (both positive and negative); find the value of k given the graphs. Experiment with cases and illustrate an explanation of the effects on the graph using technology. Include recognizing even and odd functions from their graphs and algebraic expressions for them. Trigonometric Functions (9-12) F-TF.5 Choose trigonometric functions to model periodic phenomena with specified amplitude, frequency, and midline.? This resource is part of 2 Physics Front Topical Units. Topic: Wave Energy Unit Title: Teaching About Waves and Wave Energy This is a unique, standards-based unit of instruction on Waves created by a high school teacher to be used with PhET interactive simulations on wave motion. It includes comprehensive lesson plans, lecture presentations, and assessments with answer keys. Be sure not to miss the "Clicker Questions" -- great introductory material. Topic: Wave Energy Unit Title: Wave Properties: Frequency, Amplitude, Period, Phase This exemplary unit of instruction was developed by a high school physics teacher to be used with PhET simulations. It includes six complete lesson plans that explore wave properties, the physics of sound, Fourier analysis, and wave phenomena such as reflection and superposition. Most of the lessons require that the simulation be open on a browser while students work. Don't miss the Clicker Questions, which can be readily downloaded for classroom use. Entire unit will take 2-3 weeks, but components may be pulled out separately. Can be used in a Physics First course, with teacher adaptation. %0 Electronic Source %A The PhET Project, %A Loeblein, Trish %D March 31, 2008 %T PhET Teacher Ideas & Activities: Wave Unit %V 2016 %N 6 February 2016 %8 March 31, 2008 %9 application/pdf %U http://phet.colorado.edu/en/contributions/view/3023 Disclaimer: ComPADRE offers citation styles as a guide only. We cannot offer interpretations about citations as this is an automated procedure. Please refer to the style manuals in the Citation Source Information area for clarifications. This is the full collection of teacher-created lesson plans and labs designed to be used with specific PhET simulations. Each resource has been approved by the PhET project, and may be freely downloaded.

http://www.compadre.org/Precollege/items/detail.cfm?ID=6883

A biological template ramps up electrode performance and scales down size. More than half the weight and size of today’s batteries comes from supporting materials that contribute nothing to storing energy. Now researchers have demonstrated that genetically engineered viruses can assemble active battery materials into a compact, regular structure, to make an ultra-thin, transparent battery electrode that stores nearly three times as much energy as those in today’s lithium-ion batteries. It is the first step toward high-capacity, self-assembling batteries. Applications could include high-energy batteries laminated invisibly to flat screens in cell phones and laptops or conformed to fit hearing aids. The same assembly technique could also lead to more effective catalysts and solar panels, according to the MIT researchers who developed the technology, by making it possible to finely control the positions of inorganic materials. “Most of it was done through genetic manipulation – giving an organism that wouldn’t normally make battery electrodes the information to make a battery electrode, and to assemble it into a device,” says Angela Belcher, a researcher on the project and an MIT professor of materials science and engineering and biological engineering. “My dream is to have a DNA sequence that codes for the synthesis of materials, and then out of a beaker to pull out a device. And I think this is a big step along that path.” The researchers, in work reported online this week in Science, used M13 viruses to make the positive electrode of a lithium-ion battery, which they tested with a conventional negative electrode. The virus is made of proteins, most of which coil to form a long, thin cylinder. By adding sequences of nucleotides to the virus’ DNA, the researchers directed these proteins to form with an additional amino acid that binds to cobalt ions. The viruses with these new proteins then coat themselves with cobalt ions in a solution, which eventually leads, after reactions with water, to cobalt oxide, an advanced battery material with much higher storage capacity than the carbon-based materials now used in lithium-ion batteries. To make an electrode, the researchers first dip a polymer electrolyte into a solution of engineered viruses. The viruses assemble into a uniform coating on the electrolyte. This coated electrolyte is then dipped into a solution containing battery materials. The viruses arrange these materials into an ordered crystal structure good for high-density batteries. [Click here for an illustration of the battery-forming process.] These electrodes proved to have twice the capacity of carbon-based ones. To improve this further, the researchers again turned to genetic engineering. While keeping the genetic code for the cobalt assembly, they added an additional strand of DNA that produces virus proteins that bind to gold. The viruses then assembled as nanowires composed of both cobalt oxide and gold particles – and the resulting electrodes stored 30 percent more energy. Using viruses to assemble inorganic materials has several advantages, says Daniel Morse, professor of molecular genetics and biochemistry at the University of California, Santa Barbara. First, the placement of the proteins, and the cobalt and gold that bind to them, is precise. The virus can also reproduce quickly, providing plenty of starting material, suggesting that this is manufacturing technique that could quickly scale up. And this assembly method does not require the costly processes now used to make battery materials. “You could do this at the industrial level really quickly,” says Brent Iverson, professor of organic chemistry and biochemistry at the University of Texas at Austin. “I can’t imagine a way to template or scaffold nanoparticles any cheaper.” Yet-Ming Chiang, materials science and engineering professor at MIT and one of Belcher’s collaborators, says that, while small batteries designed for specific applications could be made using this process within a couple of years, much work remains to be done. For example, cobalt oxide might not be the best material, so the researchers will be engineering viruses to bind to other materials. One of the ways they have done this in the past is using a process called “directed evolution.” They combine collections of viruses with millions of random variations in a vial containing a piece of the material they want the virus to bind to. Some of the viruses happen to have proteins that bind to the material. Isolating these viruses is a simple process of washing off the piece of material –only those viruses bound to the material remain. These can then be allowed to reproduce. After a few rounds of binding and washing, only viruses with the highest affinity for the material remain. The researchers also want to make viruses that assemble the negative electrode as well. They would then grow the positive and negative electrodes on opposite sides of a self-assembling polymer electrolyte developed by Paula Hammond*, another major contributor to the project. This would create self-assembled batteries, not just electrodes. Another goal is to make “interdigitated” batteries in which negative and positive electrode materials alternate, like the tines of two combs pushed together – this could pack in more energy and lead to batteries that deliver that energy in more powerful bursts. And batteries could be just the beginning. Since the viruses have different proteins at different locations – one protein in the center and others at the ends – the researchers can create viruses that bind to one material in the middle and different materials on the ends. Already, Belcher’s group has produced viruses that coat themselves with semiconductors and then attach themselves at the ends to gold electrodes, which could lead to working transistors. “If you can make batteries that truly are effective this way, it’s just mind-boggling what the applications could be,” Iverson says. *Correction: The virus-battery work was the result of a collaboration between researchers at MIT. The original article mentions Angela Belcher and Yet-Ming Chiang. An important part of this work was the development of a self-assembling polymer electrolyte by Paula Hammond, MIT chemical engineering professor. Home page image courtesy of Angela Belcher, MIT.

https://www.technologyreview.com/s/405635/virus-assembled-batteries/

Physics Demo -- Jumping Ring A solid metal ring is placed on an iron core whose base is wrapped in wire. When DC current is passed through the wire, a magnetic field is formed in the iron core. This sudden magnetic field induces a current in the metal ring, which in turn creates another magnetic field that opposes the original field. This causes the ring to briefly jump upwards. If there is a cut in the ring, it cannot form current inside it, and thus will not jump. When the ring is cooled in liquid nitrogen, the resistance of the metal is lowered, allowing more current to flow. This lets the ring jump higher. However, the magnetic field curves away at the top of the iron coil, meaning with DC power, the ring will never fly off the top. When AC current is passed through the wire, the ring flies off the top of the iron core. This is due to the fact that the current lags the emf by 90 degrees in inductors (which is what we have here). This yields forces on the ring that are always pointing upwards, even as the current oscillates. You need to log in, in order to post comments. Added over 7 years ago | 00:00:57 | 42244 views Added over 7 years ago | 00:00:53 | 39215 views Added almost 6 years ago | 00:02:02 | 37878 views Added over 6 years ago | 00:01:19 | 46821 views Added 7 years ago | 00:02:54 | 47274 views Added 7 years ago | 00:00:55 | 44012 views

http://techtv.mit.edu/collections/physicsdemos/videos/514-physics-demo-jumping-ring

Precession is a change in the orientation of the rotational axis of a rotating body. In an appropriate reference frame it can be defined as a change in the first Euler angle, whereas the third Euler angle defines the rotation itself. In other words, if the axis of rotation of a body is itself rotating about a second axis, that body is said to be precessing about the second axis. A motion in which the second Euler angle changes is called nutation. In physics, there are two types of precession: torque-free and torque-induced. In astronomy, "precession" refers to any of several slow changes in an astronomical body's rotational or orbital parameters, and especially to Earth's precession of the equinoxes. (See section Astronomy below.) Torque-free precession implies that no external moment (torque) is applied to the body. In torque-free precession, the angular momentum is a constant, but the angular velocity vector changes orientation with time. What makes this possible is a time-varying moment of inertia, or more precisely, a time-varying inertia matrix. The inertia matrix is composed of the moments of inertia of a body calculated with respect to separate coordinate axes (e.g. x, y, z). If an object is asymmetric about its principal axis of rotation, the moment of inertia with respect to each coordinate direction will change with time, while preserving angular momentum. The result is that the component of the angular velocities of the body about each axis will vary inversely with each axis' moment of inertia. The torque-free precession rate of an object with an axis of symmetry, such as a disk, spinning about an axis not aligned with that axis of symmetry can be calculated as follows: where is the precession rate, is the spin rate about the axis of symmetry, is the moment of inertia about the axis of symmetry, is moment of inertia about either of the other two equal perpendicular principal axes, and is the angle between the moment of inertia direction and the symmetry axis. For a generic solid object without any axis of symmetry, the evolution of the object's orientation, represented (for example) by a rotation matrix that transforms internal to external coordinates, may be numerically simulated. Given the object's fixed internal moment of inertia tensor and fixed external angular momentum , the instantaneous angular velocity is . Precession occurs by repeatedly recalculating and applying a small rotation vector for the short time ; e.g., for the skew-symmetric matrix . The errors induced by finite time steps tend to increase the rotational kinetic energy, ; this unphysical tendency can be counter-acted by repeatedly applying a small rotation vector perpendicular to both and , noting that . Another type of torque-free precession can occur when there are multiple reference frames at work. For example, Earth is subject to local torque induced precession due to the gravity of the sun and moon acting on Earth's axis, but at the same time the solar system is moving around the galactic center. As a consequence, an accurate measurement of Earth's axial reorientation relative to objects outside the frame of the moving galaxy (such as distant quasars commonly used as precession measurement reference points) must account for a minor amount of non-local torque-free precession, due to the solar system’s motion. Torque-induced precession (gyroscopic precession) is the phenomenon in which the axis of a spinning object (e.g.,a gyroscope) describes a cone in space when an external torque is applied to it. The phenomenon is commonly seen in a spinning toy top, but all rotating objects can undergo precession. If the speed of the rotation and the magnitude of the external torque are constant, the spin axis will move at right angles to the direction that would intuitively result from the external torque. In the case of a toy top, its weight is acting downwards from its center of mass and the normal force (reaction) of the ground is pushing up on it at the point of contact with the support. These two opposite forces produce a torque which causes the top to precess. The device depicted on the right is gimbal mounted. From inside to outside there are three axes of rotation: the hub of the wheel, the gimbal axis, and the vertical pivot. To distinguish between the two horizontal axes, rotation around the wheel hub will be called spinning, and rotation around the gimbal axis will be called pitching. Rotation around the vertical pivot axis is called rotation. First, imagine that the entire device is rotating around the (vertical) pivot axis. Then, spinning of the wheel (around the wheelhub) is added. Imagine the gimbal axis to be locked, so that the wheel cannot pitch. The gimbal axis has sensors, that measure whether there is a torque around the gimbal axis. In the picture, a section of the wheel has been named dm1. At the depicted moment in time, section dm1 is at the perimeter of the rotating motion around the (vertical) pivot axis. Section dm1, therefore, has a lot of angular rotating velocity with respect to the rotation around the pivot axis, and as dm1 is forced closer to the pivot axis of the rotation (by the wheel spinning further), because of the Coriolis effect, with respect to the vertical pivot axis, dm1 tends to move in the direction of the top-left arrow in the diagram (shown at 45°) in the direction of rotation around the pivot axis. Section dm2 of the wheel is moving away from the pivot axis, and so a force (again, a Coriolis force) acts in the same direction as in the case of dm1. Note that both arrows point in the same direction. The same reasoning applies for the bottom half of the wheel, but there the arrows point in the opposite direction to that of the top arrows. Combined over the entire wheel, there is a torque around the gimbal axis when some spinning is added to rotation around a vertical axis. It is important to note that the torque around the gimbal axis arises without any delay; the response is instantaneous. In the discussion above, the setup was kept unchanging by preventing pitching around the gimbal axis. In the case of a spinning toy top, when the spinning top starts tilting, gravity exerts a torque. However, instead of rolling over, the spinning top just pitches a little. This pitching motion reorients the spinning top with respect to the torque that is being exerted. The result is that the torque exerted by gravity – via the pitching motion – elicits gyroscopic precession (which in turn yields a counter torque against the gravity torque) rather than causing the spinning top to fall to its side. Precession is the result of the angular velocity of rotation and the angular velocity produced by the torque. It is an angular velocity about a line that makes an angle with the permanent rotation axis, and this angle lies in a plane at right angles to the plane of the couple producing the torque. The permanent axis must turn towards this line, because the body cannot continue to rotate about any line that is not a principal axis of maximum moment of inertia; that is, the permanent axis turns in a direction at right angles to that in which the torque might be expected to turn it. If the rotating body is symmetrical and its motion unconstrained, and, if the torque on the spin axis is at right angles to that axis, the axis of precession will be perpendicular to both the spin axis and torque axis. Under these circumstances the angular velocity of precession is given by: Where Is is the moment of inertia, is the angular velocity of spin about the spin axis, m is the mass, g is the acceleration due to gravity and r is the perpendicular distance of the spin axis about the axis of precession. The torque vector originates at the center of mass. Using = , we find that the period of precession is given by: There is a non-mathematical way of understanding the cause of gyroscopic precession. The behavior of spinning objects simply obeys the law of inertia by resisting any change in direction. If a force is applied to a spinning object to induce a change the direction of the spin axis, the object behaves as if that force was applied at a location exactly 90 degrees ahead, in the direction of rotation. This is why: A solid object can be thought of as an assembly of individual molecules. If the object is spinning, each molecule's direction of travel constantly changes as that molecule revolves around the object's spin axis. When a force is applied that is parallel to the axis, molecules are being forced to move in new directions at certain places during their path around the axis. These new changes in direction are resisted by inertia. Imagine the object to be a spinning bicycle wheel, held at both ends of its axle in the hands of a subject. The wheel is spinning clock-wise as seen from a viewer to the subject’s right. Clock positions on the wheel are given relative to this viewer. As the wheel spins, the molecules comprising it are travelling vertically downward the instant they pass the 3-o'clock position, horizontally to the left the instant they pass 6 o'clock, vertically upward at 9 o'clock, and horizontally to the right at 12 o'clock. Between these positions, each molecule travels components of these directions, which should be kept in mind as you read ahead. The viewer then applies a force to the wheel at the 3-o'clock position in a direction away from himself. The molecules at the 3-o'clock position are not being forced to change their direction when this happens; they still travel vertically downward. Actually, the force attempts to displace them some amount horizontally at that moment, but the ostensible component of that motion, attributed to the horizontal force, never occurs, as it would if the wheel was not spinning. Therefore, neither the horizontal nor downward components of travel are affected by the horizontally-applied force. The horizontal component started at zero and remains at zero, and the downward component is at its maximum and remains at maximum. The same holds true for the molecules located at 9 o'clock; they still travel vertically upward and not at all horizontally, thus are unaffected by the force that was applied. However, molecules at 6 and 12 o'clock are being forced to change direction. At 6 o'clock, molecules are forced to veer toward the viewer. At the same time, molecules that are passing 12 o'clock are being forced to veer away from the viewer. The inertia of those molecules resists this change in direction. The result is that they apply an equal and opposite reactive force in response. At 6 o'clock, molecules exert a push directly away from the viewer, while molecules at 12 o'clock push directly toward the viewer. This all happens instantaneously as the force is applied at 3 o'clock. Since no physical force was actually applied at 6 or 12 o’clock, there is nothing to oppose these reactive forces; therefore, the reaction is free to take place. This makes the wheel as a whole tilt toward the viewer. Thus, when the force was applied at 3 o'clock, the wheel behaved as if that force was applied at 6 o'clock, which is 90 degrees ahead in the direction of rotation. This principle is demonstrated in helicopters. Helicopter controls are rigged so that inputs to them are transmitted to the rotor blades at points 90 degrees prior to the point where the change in aircraft attitude is desired. Precession causes another phenomenon for spinning objects such as the bicycle wheel in this scenario. If the subject holding the wheel removes one hand from the end of the axle, the wheel will not topple over, but will remain upright, supported at just the other end of its axle. However, it will immediately take on an additional motion; it will begin to rotate about a vertical axis, pivoting at the point of support as it continues spinning. If the wheel was not spinning, it would topple over and fall when one hand is removed. The ostensible action of the wheel beginning to topple over is equivalent to applying a force to it at 12 o'clock in the direction of the unsupported side (or a force at 6 o’clock toward the supported side). When the wheel is spinning, the sudden lack of support at one end of its axle is equivalent to this same force. So, instead of toppling over, the wheel behaves as if a continuous force is being applied to it at 3 or 9 o’clock, depending on the direction of spin and which hand was removed. This causes the wheel to begin pivoting at the point of support while remaining upright. It should be noted that although it pivots at the point of support, it does so only because of the fact that it is supported there; the actual axis of precessional rotation is located vertically through the wheel, passing through its center of mass. Also, this explanation does not account for the effect of variation in the speed of the spinning object; it only describes how the spin axis behaves due to precession. More correctly, the object behaves according to the balance of all forces based on the magnitude of the applied force, mass and rotational speed of the object. The special and general theories of relativity give three types of corrections to the Newtonian precession, of a gyroscope near a large mass such as Earth, described above. They are: - Thomas precession a special relativistic correction accounting for the observer's being in a rotating non-inertial frame. - de Sitter precession a general relativistic correction accounting for the Schwarzschild metric of curved space near a large non-rotating mass. - Lense–Thirring precession a general relativistic correction accounting for the frame dragging by the Kerr metric of curved space near a large rotating mass. In astronomy, precession refers to any of several gravity-induced, slow and continuous changes in an astronomical body's rotational axis or orbital path. Precession of the equinoxes, perihelion precession, changes in the tilt of Earth's axis to its orbit, and the eccentricity of its orbit over tens of thousands of years are all important parts of the astronomical theory of ice ages. (See Milankovitch cycles.) Axial precession (precession of the equinoxes) Axial precession is the movement of the rotational axis of an astronomical body, whereby the axis slowly traces out a cone. In the case of Earth, this type of precession is also known as the precession of the equinoxes, lunisolar precession, or precession of the equator. Earth goes through one such complete precessional cycle in a period of approximately 26,000 years or 1° every 72 years, during which the positions of stars will slowly change in both equatorial coordinates and ecliptic longitude. Over this cycle, Earth's north axial pole moves from where it is now, within 1° of Polaris, in a circle around the ecliptic pole, with an angular radius of about 23.5 degrees. Hipparchus is the earliest known astronomer to recognize and assess the precession of the equinoxes at about 1° per century (which is not far from the actual value for antiquity, 1.38°). The precession of Earth's axis was later explained by Newtonian physics. Being an oblate spheroid, Earth has a non-spherical shape, bulging outward at the equator. The gravitational tidal forces of the Moon and Sun apply torque to the equator, attempting to pull the equatorial bulge into the plane of the ecliptic, but instead causing it to precess. The torque exerted by the planets, particularly Jupiter, also plays a role. The orbits of a planet around the Sun do not really follow an identical ellipse each time, but actually trace out a flower-petal shape because the major axis of each planet's elliptical orbit also precesses within its orbital plane, partly in response to perturbations in the form of the changing gravitational forces exerted by other planets. This is called perihelion precession or apsidal precession. In the adjunct image, the Earth apsidal precession is illustrated. As the Earth travels around the Sun, its elliptical orbit rotates gradually over time. The eccentricity of its ellipse and the precession rate of its orbit are exaggerated for visualization. Most orbits in the Solar System have a much smaller eccentricity and precess at a much slower rate, making them nearly circular and stationary. Discrepancies between the observed perihelion precession rate of the planet Mercury and that predicted by classical mechanics were prominent among the forms of experimental evidence leading to the acceptance of Einstein's Theory of Relativity (in particular, his General Theory of Relativity), which accurately predicted the anomalies. Deviating from Newton's law, Einstein's theory of gravitation predicts an extra term of A/r4, which accurately gives the observed excess turning rate of 43 arcseconds every 100 years. The gravitational force between the Sun and moon induces the precession in Earth's orbit, which is the major cause of the widely known climate oscillation of Earth that has a period of 19,000 to 23,000 years. It follows that changes in Earth's orbital parameters (e.g., orbital inclination, the angle between Earth's rotation axis and its plane of orbit) is important to the study of Earth's climate, in particular to the study of past ice ages. (See also nodal precession. For precession of the lunar orbit see lunar precession). |Wikimedia Commons has media related to Precession.| - Schaub, Hanspeter (2003), Analytical Mechanics of Space Systems, AIAA, pp. 149–150, ISBN 9781600860270, retrieved May 2014 - Boal, David (2001). "Lecture 26 – Torque-free rotation – body-fixed axes" (PDF). Retrieved 2008-09-17. - Teodorescu, Petre P (2002). Mechanical Systems, Classical Models. Springer. p. 420. - DIO 9.1 ‡3 - Bradt, Hale (2007). Astronomy Methods. Cambridge University Press. p. 66. ISBN 978 0 521 53551 9. - Max Born (1924), Einstein's Theory of Relativity (The 1962 Dover edition, page 348 lists a table documenting the observed and calculated values for the precession of the perihelion of Mercury, Venus, and Earth.) - An even larger value for a precession has been found, for a black hole in orbit around a much more massive black hole, amounting to 39 degrees each orbit. |Wikibooks has a book on the topic of: Rotational Motion| - Explanation and derivation of formula for precession of a top - Precession and the Milankovich theory from educational web site From Stargazers to Starships

https://en.wikipedia.org/wiki/Precession_of_the_equinox

This tutorial explains how to read, construct, and interpret basic kinematic graphs. Animated examples are accompanied by detailed discussion of how to understand the patterns produced by Position vs. Time, Velocity vs. Time, and Acceleration vs. Time graphs. The resource includes supplementary practice exercises, worksheets, and related problems for student exploration. Please note that this resource requires Microsoft Internet Explorer. acceleration, acceleration vs. time graph, average acceleration, displacement, graphing motion, instantaneous acceleration, kinematics, motion graphs, position vs. time graph, velocity, velocity vs. time graph Metadata instance created October 23, 2006 by Caroline Hall 6-8: 9B/M3. Graphs can show a variety of possible relationships between two variables. As one variable increases uniformly, the other may do one of the following: increase or decrease steadily, increase or decrease faster and faster, get closer and closer to some limiting value, reach some intermediate maximum or minimum, alternately increase and decrease, increase or decrease in steps, or do something different from any of these. 9-12: 9B/H4. Tables, graphs, and symbols are alternative ways of representing data and relationships that can be translated from one to another. Next Generation Science Standards Crosscutting Concepts (K-12) Graphs and charts can be used to identify patterns in data. (6-8) NGSS Science and Engineering Practices (K-12) Using Mathematics and Computational Thinking (5-12) Mathematical and computational thinking at the 9–12 level builds on K–8 and progresses to using algebraic thinking and analysis, a range of linear and nonlinear functions including trigonometric functions, exponentials and logarithms, and computational tools for statistical analysis to analyze, represent, and model data. Simple computational simulations are created and used based on mathematical models of basic assumptions. (9-12) Use mathematical representations of phenomena to describe explanations. (9-12) Common Core State Standards for Mathematics Alignments High School — Functions (9-12) Interpreting Functions (9-12) F-IF.5 Relate the domain of a function to its graph and, where applicable, to the quantitative relationship it describes.? F-IF.6 Calculate and interpret the average rate of change of a function (presented symbolically or as a table) over a specified interval. Estimate the rate of change from a graph. F-IF.7.a Graph linear and quadratic functions and show intercepts, maxima, and minima. F-IF.7.e Graph exponential and logarithmic functions, showing intercepts and end behavior, and trigonometric functions, showing period, midline, and amplitude. Linear, Quadratic, and Exponential Models? (9-12) F-LE.1.a Prove that linear functions grow by equal differences over equal intervals, and that exponential functions grow by equal factors over equal intervals. F-LE.3 Observe using graphs and tables that a quantity increasing exponentially eventually exceeds a quantity increasing linearly, quadratically, or (more generally) as a polynomial function. Common Core State Reading Standards for Literacy in Science and Technical Subjects 6—12 Range of Reading and Level of Text Complexity (6-12) RST.11-12.10 By the end of grade 12, read and comprehend science/technical texts in the grades 11—CCR text complexity band independently and proficiently. This resource is part of a Physics Front Topical Unit. Topic: Kinematics: The Physics of Motion Unit Title: Graphing A very well-organized tutorial on how to construct and interpret three basic kinematic graphs: P/T, V/T and A/T. It includes animated examples, links to five worksheets, and related problems for student exploration. %0 Electronic Source %A Elert, Glenn %D June 27, 2007 %T The Physics Hypertextbook: Graphs of Motion %V 2016 %N 8 February 2016 %8 June 27, 2007 %9 text/html %U http://physics.info/motion-graphs/ Disclaimer: ComPADRE offers citation styles as a guide only. We cannot offer interpretations about citations as this is an automated procedure. Please refer to the style manuals in the Citation Source Information area for clarifications.

http://www.thephysicsfront.org/items/detail.cfm?ID=4547

Updated at 6:05 p.m. ET Valleys on Mars were carved over long periods by recurring floods at a time when Mars might have had wet and dry seasons much like some of Earth's deserts, a new study suggests. The research contradicts other suggestions that the large valley networks on the red planet were the result of short-lived catastrophic flooding, lasting just hundreds to a few thousand years and perhaps triggered by asteroid impacts. The new modeling suggests wet periods lasted at least 10,000 years. "Precipitation on Mars lasted a long time – it wasn't a brief interval of massive deluges," said study leader Charles Barnhart, a graduate student in Earth and planetary sciences at the University of California, Santa Cruz. "Our results argue for liquid water being stable at the surface of Mars for prolonged periods in the past." NASA planetary scientist Jeffrey Moore and Alan Howard of the University of Virginia contributed to the research, which will be detailed in the Journal of Geophysical Research – Planets. In recent years, pictures of Mars have revealed a landscape clearly shaped by runoff. Most researchers now see water as a key player, though carbon dioxide has also been put forth as a possible culprit. A recent study of Box Canyon in Idaho concluded that similar features on Mars could have been sculpted by ancient megafloods. Whatever, Mars is bone-dry today, and it's not clear just how wet it was in the past. The new work, based on computer models, paints a picture of ancient Mars, more than 3.5 billion years ago, as looking somewhat like the deserts of the U.S. Southwest, sans cacti of course. The new thinking is based on the idea that asteroid impacts, as fuel for floods, would have created features that aren't found on Mars. "Our research finds that these catastrophic anomalies would be so humid and wet there would be breaching of the craters, which we don't see on Mars," he said. "The precipitation needs to be seasonal or periodic, so that there are periods of evaporation and infiltration. Otherwise the craters overflow." The computer models generated 70 different simulations, including one that yielded the best match to the observed topography of martian valleys. The sculpting was done in a semiarid to arid climate that persisted for tens of thousands to hundreds of thousands of years, the researchers say. Episodic flooding alternated with long dry periods when water could evaporate or soak into the ground. Rainfall may have been seasonal, or wet intervals may have occurred over longer cycles. But conditions that allowed for the presence of liquid water on the surface of Mars must have lasted for at least 10,000 years, Barnhart said. The study does not suggest what a typical day on Mars might have been like back then. Nor does the work pin down how long seasons might have lasted. Barnhart said the changes from dry to wet periods might have had to do with periods of greenhouse-gas outgassing associated with volcanic eruptions, large impacts, or a change in the tilt of Mars' rotation, though all that remains to be studied further. "Our results do suggest that river discharges were similar to flood stages in Earth-like desert environments like the Mojave desert or the Colorado Plateau," he told SPACE.com.

http://www.space.com/5816-long-wet-periods-sculpted-ancient-mars.html

Portuguese orthography is based on the Latin alphabet and makes use of the acute accent, the circumflex accent, the grave accent, the tilde, and the cedilla, to denote stress, vowel height, nasalization, and other sound changes. Accented letters and digraphs are not counted as separate characters for collation purposes. The spelling of Portuguese is largely phonemic, but some phonemes can be spelled in more than one way. In ambiguous cases, the correct spelling is determined through a combination of etymology with morphology and tradition so there is not a perfect one-to-one correspondence between sounds and letters or digraphs. Knowing the main inflectional paradigms of Portuguese and being acquainted with the orthography of other Western European languages can be helpful. A full list of sounds, diphthongs, and their main spellings, is given at Portuguese phonology. This article addresses the less trivial details of the spelling of Portuguese as well as other issues of orthography, such as accentuation. - 1 Letter names and pronunciations - 2 Digraphs - 3 Diacritics - 4 Consonants with more than one spelling - 5 Vowels - 6 Morphological considerations - 7 Etymological considerations - 8 Syllabification and collation - 9 Other symbols - 10 Brazilian vs. European spelling - 11 See also - 12 References - 13 External links Letter names and pronunciations Only the most frequent sounds appear below since a listing of all cases and exceptions would become cumbersome. Portuguese is a pluricentric language, and pronunciation of some of the letters differs in European Portuguese (EP) and in Brazilian Portuguese (BP). Apart from those variations, the pronunciation of most consonants is fairly straightforward. Only the consonants r, s, x, z, the digraphs ch, lh, nh, rr, and the vowels may require special attention from English speakers. Although many letters have more than one pronunciation, their phonetic value is often predictable from their position within a word; that is normally the case for the consonants (except x). Since only five letters are available to write the fourteen vowel sounds of Portuguese, vowels have a more complex orthography, but even then, pronunciation is somewhat predictable. Knowing the main inflectional paradigms of Portuguese can help. In the following table and in the remainder of this article, the phrase "at the end of a syllable" can be understood as "before a consonant, or at the end of a word". For the letter r, "at the start of a syllable" means "at the beginning of a word, or after l, n, s". For letters with more than one common pronunciation, their most common phonetic values are given on the left side of the semicolon; sounds after it occur only in a limited number of positions within a word. Sounds separated by "~" are allophones or dialectal variants. The names of the letters are masculine. Letter European Brazilian Phonemic Name Name (IPA) Name Name (IPA) Aa á /a/ á /a/ /a/, /ɐ/ Bb bê /be/ bê /be/ /b/ Cc cê /se/ cê /se/ /k/; /s/ nb 1 Dd dê /de/ dê /de/ /d/ ~ [dʒ] nb 2 Ee é /ɛ/ é or ê /ɛ/, /e/ /e/, /ɛ/, /i/ nb 3, /ɨ/, /ɐ/, /ɐi/ Ff efe /ˈɛfɨ/ efe /ˈɛfi/ /f/ Gg gê or guê /ʒe/, /ɡe/ gê /ʒe/ /ɡ/; /ʒ/ nb 1 Hh agá /ɐˈɡa/ agá /aˈɡa/ natively silent, /ʁ/ in loanwords nb 4 Ii i /i/ i /i/ /i/ nb 3 Jj jota /ˈʒɔtɐ/ jota /ˈʒɔta/ /ʒ/ Kk capa /ˈkapɐ/ cá /ka/ nb 5 Ll ele /ˈɛlɨ/ ele /ˈɛli/ /l/ ~ [ɫ ~ w] nb 6 Mm eme /ˈɛmɨ/ eme /ˈemi/ /m/ nb 7 Nn ene /ˈɛnɨ/ ene /ˈeni/ /n/ nb 7 Oo ó /ɔ/ ó or ô /ɔ/, /o/ /o/, /ɔ/, /u/ nb 3 Pp pê /pe/ pê /pe/ /p/ quê /ke/ quê /ke/ /k/ Rr erre or rê /ˈɛʁɨ/, /ˈʁe/ erre /ˈɛʁi/ /ɾ/, /ʁ/ nb 8 Ss esse /ˈɛsɨ/ esse /ˈɛsi/ /s/, /z/ nb 9, /ʃ/ nb 10 Tt tê /te/ tê /te/ /t/ ~ [tʃ] nb 2 Uu u /u/ u /u/ /u/ nb 3 Vv vê /ve/ vê /ve/ /v/ Ww dâblio or duplo vê /ˈdɐbliu/ dáblio or duplo vê /ˈdabliu/ nb 5 Xx xis /ʃiʃ/ xis /ʃis/ /ʃ/, /ks/, /z/, /s/ nb 10 nb 11 Yy ípsilon or i grego /ˈipsɨlɔn/ ípsilon /ˈipsilõ/ nb 5 Zz zê /ze/ zê /ze/ /z/, /s/, /ʃ/ nb 10 Listen to the alphabet recited by a native speaker from Brazil. The alphabet is spoken in a Brazilian dialect in which the 'E' is pronounced as 'É' |Problems playing this file? See media help.| - ^ Before the letters e, i, y, or with the cedilla. - ^ Allophonically affricated before the sound /i/ (spelled i, or sometimes e), in BP. - ^ May become an approximant as a form of vowel reduction when unstressed before or after another vowel. Words such as bóia and proa are pronounced [ˈbɔj.jɐ] and [ˈpɾow.wɐ]. - ^ Silent at the start or at the end of a word. Also part of the digraphs ch, lh, nh. See below. - ^ Not part of the official alphabet before 2009. Used only in foreign words, personal names, and hybrid words derived from them. The letters K, W and Y were included in the alphabet used in Brazil, East Timor, Macau, Portugal and five countries in Africa, when the 1990 Portuguese Language Orthographic Agreement went into legal effect, since January 1, 2009. However, they were used before 1911 (see the article on spelling reform in Portugal). - ^ Velarized to [ɫ] in EP and conservative registers of southern BP. Vocalized to [u̯], [ʊ̯], or seldom [o̯] (as influence from Spanish or Japanese), at the end of syllables in most of Brazil. - ^ Usually silent or voiceless at the end of syllables (word-final n is fully pronounced by some speakers in a few loaned words). See Nasalization section, below. - ^ At the start of syllables (in all dialects) or at the end of syllables (in some dialects of BP), a single r is pronounced /ʁ/ (see Portuguese phonology for variants of this sound). Elsewhere, it is pronounced /ɾ/. Word-final rhotics may also be silent when the last syllable is stressed, in casual and vernacular speech, especially in Brazil (pervasive nationwide, though not in educated and some colloquial registers) and in some African and Asian countries. - ^ A single s is pronounced voiced /z/ between vowels. - ^ The opposition between the four sibilants /s/, /z/, /ʃ/, /ʒ/ is neutralized at the end of syllables; see below for more information. - ^ The letter x may represent /ʃ/, /ks/, /z/, and /s/ (peixe, fixar, exemplo, próximo). It is always pronounced /ʃ/ at the beginning of words. Portuguese uses of digraphs, pairs of letters which represent a single sound different from the sum of their components. Digraphs are not included in the alphabet. The digraphs qu and gu, before e and i, may represent both plain or labialised sounds (quebra /ˈkebɾɐ/, cinquenta /sĩˈkʷẽtɐ/, guerra /ˈɡɛʁɐ/, sagui /saˈɡʷi/), but they are always labialised before a and o (quase, quociente, guaraná). Pronunciation divergences make some of those words be spelled in diffent forms (quatorze / catorze and quotidiano / cotidiano). The digraph ch is pronounced as an English sh by the overwhelmingly majority of speakers. The digraphs lh and nh, of Occitan origin, denote palatal consonants that do not exist in English. The digraphs rr and ss are used only between vowels. The pronunciation of the digraph rr varies with dialect (see the note on the phoneme /ʁ/, above). The acute accent and the circumflex accent indicate that a vowel is stressed and the quality of the accented vowel and, more precisely, its height: á, é, and ó are low vowels (except in nasal vowels); â, ê, and ô are high vowels. They also distinguish a few homographs: por "by" with pôr "to put", pode "[he/she/it] can" with pôde "[he/she/it] could". The tilde marks nasal vowels before glides such as in cãibra and nação, at the end of words, before final -s, and in some compounds: romãzeira "pomegranate tree", from romã "pomegranate", and vãmente "vainly", from vã "vain". It usually coincides with the stressed vowel unless there is an acute or circumflex accent elsewhere in the word or if the word is compound: órgão "organ", irmã + zinha ("sister" + diminutive suffix) = irmãzinha "little sister". The form õ is used only in the plurals of nouns ending in -ão (nação → nações) and in the second and third person singular forms of the verb pôr (pões, põe). The grave accent marks the contraction of two consecutive vowels in adjacent words (crasis), normally the preposition a and an article or a demonstrative pronoun: a + aquela = àquela "at that", a + a = à "at the". It does not indicate stress. The graphemes â, ê and ô typically represent oral vowels, but before m or n followed by another consonant, the vowels represented are nasal. Elsewhere, nasal vowels are indicated with a tilde (ã, õ). Below are the general rules for the use of the acute accent and the circumflex in Portuguese. Primary stress may fall on any of the three final syllables of a word but occurs usually on the last two. A word is called oxytone if it is stressed on its last syllable, paroxytone if stress falls on the syllable before the last (the penult), and proparoxytone if stress falls on the third syllable from the end (the antepenult). Most multisyllabic words are stressed on the penult. All words stressed on the antepenult take an accent mark. Words with two or more syllables, stressed on their last syllable, are not accented if they end with any consonant letter but -m and -s or -i, -is, -im, -u, -us, -um except in hiatuses as in açaí, but paroxytonic words may then be accented to differentiate them from oxytonic words, as in lápis. Monosyllables are typically not accented, but those whose last vowel is a, e, or o, possibly followed by final -s or final -m, may require an accent mark. - The verb pôr is accented to distinguish it from the preposition por. - Third-person plural forms of the verbs ter and vir, têm and vêm are accented to be distinguished from third-person singulars of the same verbs, tem, vem. Other monosyllables ending in -em are not accented. - Monosyllables ending in -o or -os with the vowel pronounced /u/ (as in English "do") or in -e or -es with the vowel pronounced /i/ (as in English "be") or /ɨ/ (approximately as in English "roses") are not accented. Otherwise, they are accented. - Monosyllables containing only the vowel a take an acute accent except for the contractions of the preposition a with the articles a, as, which take the grave accent, à, às, and for the following clitic articles, pronouns, prepositions, or contractions, which are not accented: a, da, la, lha, ma, na, ta; as, das, las, lhas, mas, nas, tas. Most of those words have a masculine equivalent ending in -o(s), also not accented: o(s), do(s), lo(s), lho(s), mo(s), no(s), to(s). - The endings -a, -e, -o, -as, -es, -os, -am, -em, -ens are unstressed. The stressed vowel of words with such endings is assumed to be the first one before the ending itself: bonita, bonitas, gente, viveram, seria, serias (verbs), seriam. If the word happens to be stressed elsewhere, it requires an accent mark: será, serás, até, séria, sérias (adjectives), Inácio, Amazônia/Amazónia. The endings -em and -ens take the acute accent when stressed (contém, convéns), except in third-person plural forms of verbs derived from ter and vir, which take the circumflex (contêm, convêm). Words with other endings are regarded as oxytone by default: viver, jardim, vivi, bambu, pensais, pensei, pensou. They require an accent when they are stressed on a syllable other than their last: táxi, fácil, amáveis. - Rising diphthongs (which may also be pronounced as hiatuses) containing stressed i or stressed u are accented so they will not be pronounced as falling diphthongs. Exceptions are those whose stressed vowel forms a syllable with a letter other than s. Thus, raízes (syllabified as ra-í-zes), incluído (u-í), and saíste (a-ís) are accented, but raiz (ra-iz), sairmos (a-ir) and saiu (a-iu) are not. (There are a few more exceptions, not discussed here.) - The stressed diphthongs ei, eu, oi take an acute accent on the first vowel whenever it is low. Aside from those cases, there are a few more words that take an accent, usually to disambiguate frequent homographs such as pode (present tense of the verb poder) and pôde (past tense of the same verb). Also, accentuation rules of Portuguese are somewhat different from those of Spanish (English "continuous" is Portuguese contínuo, Spanish continuo, and English "I continue" is Portuguese continuo, Spanish continúo, in both cases with the same syllable accented in Portuguese and Spanish). The use of diacritics in personal names is generally restricted to the combinations above, often also by the applicable Portuguese spelling rules. Portugal is more restrictive than Brazil in regard to given names. They must be Portuguese or adapted to the Portuguese orthography and sound and should also be easily discerned as either a masculine or feminine name by a Portuguese speaker. There are lists of previously accepted names, and names not included therein must be subject to consultation of the national director of registries. Brazilian birth registrars, on the other hand, are likely to accept names containing any (Latin) letters or diacritics and are limited only to the availability of such characters in their typesetting facility. Consonants with more than one spelling Most consonants have the same values as in the International Phonetic Alphabet, except for the palatals /ʎ/ and /ɲ/, which are spelled lh and nh, respectively, and the following velars, rhotics, and sibilants: |Phoneme||Default||Before e or i| The alveolar tap /ɾ/ is always spelled as a single r. The other rhotic phoneme of Portuguese, which may be pronounced as a trill [r] or as one of the fricatives [x], [ʁ], or [h], according to the idiolect of the speaker, is either written rr or r, as described below. |Phoneme||Start of syllable[rhotic note 1]||Between vowels||End of syllable[rhotic note 2]| |/ʁ/||r||rosa, tenro||rr||carro||r||sorte, mar| - only when it is the first sound in the syllable (in which case it is always followed by a vowel). For instance, a word like prato is pronounced with a tap, /ɾ/ - in some dialects; in the others, the r is usually a tap or approximant at the end of syllables For the following phonemes, the phrase "at the start of a syllable" can be understood as "at the start of a word, or between a consonant and a vowel, in that order". |Phoneme||Start of syllable1||Between vowels||End of syllable| |/s/||s, c3||sapo, psique, |ss, ç2, c3, x4||assado, passe, |s, x5, z6||isto, |z, s, x7||prazo, azeite, |s, x8, z8||turismo, |/ʃ/||ch, x||chuva, cherne, |ch, x||fecho, duche, |s, x5, z6||isto, |/ʒ/||j, g3||jogo, jipe, |j, g3||ajuda, pajem, |s, x8, z8||turismo, - 1 including consonant clusters that belong to a single syllable, like psique - 2 before a, o, u. Ç never starts or ends a word. - 3 before e, i - 4 only in a very small number of words derived from Latin, such as trouxe and próximo - 5 only in words derived from Latin or Greek, preceded by e and followed by one of the voiceless consonants c, p, s, t - 6 only at the end of words and in rare compounds - 7 only in a few words derived from Latin or Greek that begin with ex- or hex- followed by a vowel, and in compounds made from such words - 8 only in a few compound words Note that there are two main groups of accents in Portuguese, one in which the sibilants are alveolar at the end of syllables (/s/ or /z/), and another in which they are postalveolar (/ʃ/ or /ʒ/). In this position, the sibilants occur in complementary distribution, voiced before voiced consonants, and voiceless before voiceless consonants or at the end of utterances. The vowels in the pairs /a, ɐ/, /e, ɛ/, /o, ɔ/ only contrast in stressed syllables. In unstressed syllables, each element of the pair occurs in complementary distribution with the other. Stressed /ɐ/ appears mostly before the nasal consonants m, n, nh, followed by a vowel, and stressed /a/ appears mostly elsewhere although they have a limited number of minimal pairs in EP. In Brazilian Portuguese, both nasal and unstressed vowel phonemes that only contrast when stressed tend to a mid height though [a] may be often heard in unstressed position (especially when singing or speaking emphatically). In pre-20th-century European Portuguese, they tended to be raised to [ə], [i] (now [ɯ̽] except when close to another vowel) and [u]. It still is the case of most Brazilian dialects in which the word elogio may be variously pronounced as [iluˈʒiu], [e̞lo̞ˈʒiu], [e̞luˈʒiu], etc. Some dialects, such as those of Northeastern and Southern Brazil, tend to do less pre-vocalic vowel reduction and in general the unstressed vowel sounds adhere to that of one of the stressed vowel pair, namely [ɛ, ɔ] and [e, o] respectively. In the educated speech, vowel reduction is used less often than in colloquial and vernacular speech though still more than the more distant dialects, and in general, mid vowels are dominant over close-mid ones and especially open-mid ones in unstressed environments when those are in free variation (that is, sozinho is always [sɔˈzĩɲu], even in Portugal, while elogio is almost certainly [e̞lo̞ˈʒi.u]). Mid vowels are also used as choice for stressed nasal vowels in both Portugal and Rio de Janeiro though not in São Paulo and southern Brazil, but in Bahia, Sergipe and neighboring areas, mid nasal vowels supposedly are close-mid like those of French. Veneno can thus vary as EP [vɯ̽ˈne̞nu], RJ [vẽ̞ˈnẽ̞nu], SP [veˈnenʊ] and BA [vɛˈnɛ̃nu] according to the dialect. /ɐ̃/ also got significant dialectal variation, respectively in the same of the last sentence, banana [baˈnə̃nə], [bə̃ˈnə̃nə], [bəˈnənə] etc. Vowel reduction of unstressed nasal vowels is extremely pervasive nationwide in Brazil, in vernacular, colloquial and even most educated speech registers: então [ĩˈtɐ̃w], camondongo [kɐmũˈdõɡu]. It slightly more resisted but still present in Portugal. The pronunciation of the accented vowels is fairly stable except that they become nasal in certain conditions. See the section on Nasalization for further information about this regular phenomenon. In other cases, nasal vowels are marked with a tilde. The diacritic ` is used only in the letter A and is merely grammatical, meaning a crasis between two a such as adverb "to" and feminine pronoun "the" (vou a a cidade to vou à cidade "I'm going to the city"), not affecting pronunciation at all. The trema was official prior to the last orthographical reform and can still be found in older texts. It meant that the usually silent u between q or g and i or e is in fact pronounced: liqüido "liquid" and sangüíneo "related to blood". Some words have two acceptable pronunciations, varying largely by accents. The pronunciation of each diphthong is also fairly predictable, but one must know how to distinguish true diphthongs from adjacent vowels in hiatus, which belong to separate syllables. For example, in the word saio /ˈsaiu/ ([ˈsaj.ju]), the i forms a clearer diphthong with the previous vowel (but a slight yod also in the next syllable is generally present), but in saiu /sɐˈiu/ ([sɐˈiw]), it forms a diphthong with the next vowel. As in Spanish, a hiatus may be indicated with an acute accent, distinguishing homographs such as saia /ˈsaiɐ/ ([ˈsaj.jɐ]) and saía /sɐˈiɐ/. Oral Grapheme Pronunciation Grapheme Pronunciation ai, ái [ai ~ ɐi] au, áu [au ~ ɐu] ei, êi [ei ~ eː], [əi]1 eu, êu [eu] éi [ɛi], [əi]1 éu [ɛu] oi [oi] ou [ou ~ oː] ói [ɔi] óu [ɔu] ui [ui] iu [iu] Nasal Grapheme Pronunciation Grapheme Pronunciation ãe, ãi [ɐ̃ĩ] ão [ɐ̃ũ] õe [õĩ] - 1 In central Portugal. When a syllable ends with m or n, the consonant is not fully pronounced but merely indicates the nasalization of the vowel which precedes it. At the end of words, it sometimes produces a nasal diphthong. Monophthongs Diphthongs Grapheme Pronunciation Grapheme Pronunciation -an, -am, -ân, -âm1 /ɐ̃/ -am2 /ɐ̃ũ/ -en, -em, -ên, -êm1 /ẽ/ -em, -ém2 /ẽĩ/ ([ɐ̃ĩ]) -in, -im, -ín, -ím3 /ĩ/ -en-, -én-4 -on, -om, -ôn, -ôm3 /õ/ -êm2 /ẽĩ/ ([ɐ̃ĩ]) -un, -um, -ún, -úm3 /ũ/ 1 at the end of a syllable 2 at the end of a word 3 at the end of a syllable or word 4 before final s, for example in the words bens and parabéns The letter m is conventionally written before b or p or at the end of words (also in a few compound words such as comummente - comumente in Brazil), and n is written before other consonants. In the plural, the ending -m changes into -ns; for example bem, rim, bom, um → bens, rins, bons, uns. Some loaned words end with -n (which is usually pronounced in European Portuguese). Nasalization of u is left unmarked in the six words muito, muita, muitos, muitas, mui, ruim (the latter one only in Brazilian Portuguese). The word endings -am, -em, -en(+s), with or without an accent mark on the vowel, represent nasal diphthongs derived from various Latin endings, often -ant, -unt or -en(t)-. Final -am, which appears in polysyllabic verbs, is always unstressed. The grapheme -en- is also pronounced as a nasal diphthong in a few compound words, such as bendito (bem + dito), homenzinho (homem + zinho), and Benfica. Verbs whose infinitive ends in -jar have j in the whole conjugation: viagem "voyage" (noun) but viajem (third person plural of the present subjunctive of the verb viajar "to travel"). Verbs whose thematic vowel becomes a stressed i in one of their inflections are spelled with an i in the whole conjugation, as are other words of the same family: crio (I create) implies criar (to create) and criatura (creature). Verbs whose thematic vowel becomes a stressed ei in one of their inflections are spelled with an e in the whole conjugation, as are other words of the same family: nomeio (I nominate) implies nomear (to nominate) and nomeação (nomination). The majority of the Portuguese lexicon is derived from Latin, Greek, and some Arabic. In principle, that would require some knowledge of those languages. However, Greek words are Latinized before being incorporated into the language, and many words of Latin or Greek origin have easily recognizable cognates in English and other western European languages and are spelled according to similar principles. For instance, glória, "glory", glorioso, "glorious", herança "inheritance", real "real/royal". Some general guidelines for spelling are given below: - CU vs. QU: if u is pronounced syllabically, it is written with c, as in cueca [kuˈɛkɐ] (underwear), and if it represents a labialized velar plosive, it is written with q, as in quando [ˈkwɐ̃du] (when). - G vs. J: etymological g changes into j before a, o, u. - H: this letter is silent; it appears for etymology at the start of a word, in a few interjections, and as part of the digraphs ch, lh, nh. Latin or Greek ch, ph, rh, th, and y are usually converted into c/qu, f, r, t, and i, respectively. - O vs. OU: in many words, the variant oi normally corresponds to Latin and Arabic au or al, more rarely to Latin ap, oc. - S/SS vs. C/Ç: the letter s and the digraph ss correspond to Latin s, ss, or ns, and to Spanish s. The graphemes c (before e or i) and ç (before a, o, u) are usually derived from Latin c or t(i), or from s in non-European languages, such as Arabic and Amerindian languages. They correspond to Spanish z or c. At the beginning of words, however, s is written instead of etymological ç, by convention. - Z vs. S between vowels: the letter z corresponds to Latin c (+e, i) or t(i), to Greek or Arabic z. Intervocalic s corresponds to Latin s. - X vs. CH: the letter x derives from Latin x or s, or from Arabic sh and usually corresponds to Spanish j. The digraph ch (before vowels) derives from Latin cl, fl, pl or from French ch and corresponds to Spanish ll or ch. - S vs. X vs. Z at the end of syllables: s is the most common spelling for all sibilants. The letter x appears, preceded by e and followed by one of the voiceless consonants c, p, s, t, in some words derived from Latin or Greek. The letter z occurs only at the end of oxytone words and in compounds derived from them, corresponding to Latin x, c (+e, i) or to Arabic z. Loanwords with a /ʃ/ in their original languages receive the letter x to represent it when they are nativised: xampu (shampoo). While the pronunciations of ch and x merged long ago, some Galician-Portuguese dialects like the Galician language, the portunhol da pampa and the speech registers of northeastern Portugal still preserve the difference as ch /tʃ/ vs. x /ʃ/, as do other Iberian languages and Medieval Portuguese. When one wants to stress the sound difference in dialects in which it merged the convention is to use tch: tchau (ciao) and Brazilian Portuguese República Tcheca (Czech Republic). In most loanwords, it merges with /ʃ/ (or /t/ :moti for mochi), just as [dʒ] most often merges with /ʒ/. Alveolar affricates [ts] and [dz], though, are more likely to be preserved (pizza, Zeitgeist, tsunami, kudzu, adzuki, etc.) Syllabification and collation Portuguese syllabification rules require a syllable break between double letters: cc, cç, mm, nn, rr, ss, or other combinations of letters that may be pronounced as a single sound: fric-ci-o-nar, pro-ces-so, car-ro, ex-ce(p)-to, ex-su-dar. Only the digraphs ch, lh, nh, gu, qu, and ou are indivisible. All digraphs are however broken down into their constituent letters for the purposes of collation, spelling aloud, and in crossword puzzles. The apostrophe (') appears as part of certain phrases, usually to indicate the elision of a vowel in the contraction of a preposition with the word that follows it: de + água = d'água. It is used almost exclusively in poetry. The hyphen (-) is used to make compound words, especially animal names like papagaio-de-rabo-vermelho "red-tailed parrot". It is also extensively used to append clitic pronouns to the verb, as in quero-o "I want it" (enclisis), or even to embed them within the verb, as in levaria + vos + os = levar-vos-ia "I would take to you", "levar-vo-los-ia" = "I would take them to you" (mesoclisis). Proclitic pronouns are not connected graphically to the verb: não o quero "I do not want it". Each element in such compounds is treated as an individual word for accentuation purposes. In European Portuguese, as in many other European languages, angular quotation marks are used for general quotations in literature: - «Isto é um exemplo de como fazer uma citação em português europeu.» - “This is an example of how to make a quotation in European Portuguese.” Although American-style (“…”) or British-style (‘…’) quotation marks are sometimes used as well, especially in less formal types of writing (they are more easily produced in keyboards) or inside nested quotations, they are less common in careful writing. In Brazilian Portuguese, only American and British-style quote marks are used. - “Isto é um exemplo de como fazer uma citação em português brasileiro.” - “This is an example of how to make a quotation in Brazilian Portuguese.” In both varieties of the language, dashes are normally used for direct speech rather than quotation marks: - ― Aborreço-me tanto ― disse ela. - ― Não tenho culpa disso ― retorquiu ele. - “I’m so bored,” she said. - “That’s not my fault,” he shot back. Brazilian vs. European spelling |Portuguese-speaking countries except Brazil before the 1990 agreement||Brazil before the 1990 agreement||All countries after the 1990 agreement||translation| |anónimo||anônimo||Both forms remain||anonymous| |Vénus||Vênus||Both forms remain||Venus| |facto||fato||Both forms remain||fact| |Non-personal and non-geographical names| As of 2005[update], Portuguese has two orthographic standards: - The Brazilian orthography, official in Brazil. - The European orthography, official in Portugal, Macau, East Timor and the five African Lusophone countries. In East Timor, both orthographies are currently being taught in schools. The table to the right illustrates typical differences between the two orthographies. Some are due to different pronunciations, but others are merely graphic. The main ones are: - Presence or absence of certain consonants: The letters c and p appear in some words before c, ç or t in one orthography, but are absent from the other. Normally, the letter is written down in the European spelling, but not in the Brazilian spelling. - Different use of diacritics: the Brazilian spelling has a, ê or ô followed by m or n before a vowel, in several words where the European orthography has á, é or ó, due to different pronunciation. - Different usage of double letters: also due to different pronunciation, Brazilian spelling has only cc, rr and ss as double letters. So, Portuguese connosco becomes Brazilian conosco and words ended in m with suffix -mente added, (like ruimmente and comummente) become ruimente e comumente in Brazilian spelling. - Academia Brasileira de Letras - Differences between Spanish and Portuguese - Portuguese names - Portuguese phonology - Spelling reforms of Portuguese - The Vietnamese orthography, partly based on the orthography of Portuguese, through the work of 16th-century Catholic missionaries. - Accordo Ortográfico de 1990 - Wikipedia in Portuguese: Ortografia da língua portuguesa - (Portuguese) Delta: Documentation of studies on theoric and applied Linguistics – Problems in the tense variant of carioca speech. - Ministro da Cultura quer Acordo vigorando antes de janeiro de 2010 [Minister of Culture wants Agreement enforced before January 2010] (in Portuguese), Portugal: Sapo. In Brazil, the Orthographic Agreement went into legal effect from January 1, 2009. - catorze / quatorze [Ortografia / Fonética e Fonologia / Etimologia] - Portal do Cidadão (Portuguese) - (Italian) Accenti romanze: Portogallo e Brasile (portoghese) – The influence of foreign accents on Italian language acquisition. - Bergström, Magnus & Reis, Neves Prontuário Ortográfico Editorial Notícias, 2004. - Estrela, Edite A questão ortográfica — Reforma e acordos da língua portuguesa (1993) Editorial Notícias - Formulário Ortográfico (Orthographic Form) published by the Brazilian Academy of Letters in 1943 - the present day spelling rules in Brazil - Text of the decree of the Brazilian government, in 1971, amending the orthography adopted in 1943 - Orthographic Agreement of 1945 (in Portuguese) - the present day spelling rules in all Portuguese speaking countries - Orthographic Agreement of 1990 (PDF - in Portuguese) - to be adopted by all Portuguese speaking countries

https://en.wikipedia.org/wiki/Portuguese_orthography

Google has been working with its quantum computer for several years now, and finally has results that prove the D-Wave 2 can perform certain calculations up to a hundred million times faster than existing conventional computers. That kind of performance gain is only possible if the system is actually performing quantum computing. German researchers have devised a technique of creating self-assembled nanodiamond quantum bits (qubits) that could form the basis of quantum computers and storage devices that, unlike every other quantum tech that we’ve seen on ET, could operate at room temperature. It used to be that all you needed to do to talk about the frontiers of quantum research was regurgitate a few thought experiments about cats in boxes, but in the past few years the pure science has started to pay practical dividends. PHD comics is the latest to help the public try to wrap their head around some crazy but vital concepts. The new Quantum Artificial Intelligence Lab (QAIL), housed at NASA’s Ames Research Center in Silicon Valley and staffed by Google and NASA scientists, has become the second lab in the world to own a quantum computer. As the name suggests, the Google and NASA scientists will use the quantum computer to advance machine learning — a field of AI that deals with computers that autonomously optimize their behavior as they garner more experience. A computer scientist at Amherst College has performed the first ever head-to-head speed test between a conventional and quantum computer — and, you’ll be glad to hear that the quantum computer won. But only just. Researchers at the University of New South Wales in Australia have created the first quantum bit (qubit) based on the nuclear spin of an atom, within a silicon transistor. This breakthrough is significant for two reasons: The qubit produced by the researchers is highly stable — and it’s in silicon, meaning it can be wired up and controlled electronically, just like a conventional computer chip. A team of quantum engineers in Germany have created the first air-to-surface quantum network, between a base station and an airplane flying 20 kilometers (12.4 miles) above. This is a very tantalizing step towards a global quantum communications network. A team of international researchers have successfully teleported a quantum bit (qubit) over a record distance of 143 kilometers (89 miles), between the Canary Islands of La Palma and Tenerife. This distance is significant, as it is roughly the same distance to low Earth orbit (LEO) satellites — meaning it is now theoretically possible to build a satellite-based quantum communication network.

http://www.extremetech.com/tag/qubit

Educators often want to know how they can use PBL in their individual classroom. Project-based learning can be applied in any content area or any grade, but it may look very different across subjects. In the series of examples below, you will find descriptions of actual projects and exercises teachers have implemented in schools around the country, as well as links to the research reports on their outcomes. It should be noted that there are also many wonderful examples of cross-curricular projects, where teachers from two or more core subjects work together on a project. For an example, check out our Schools That Work package on an interdisciplinary project at Manor High School in Texas. Schools That Work: English teacher Mary Mobley (left) shared the PBL resources and tools she and her teaching partner created for their sophomore world studies project (right). Learn more about this project. Credit: Zachary Fink Urban students in grades 3-5 received inquiry-science instruction. Matched pre- and post-tests found substantial learning gains and a cumulative effect that lasted over several years (Lee, Buxton, Lewis, & LeRoy, 2006). Fourth graders learned science through PBL or through traditional methods with the same teacher. The PBL curriculum involved figuring out a way to create electricity during a blackout, as blackouts had commonly affected the school’s region. PBL students had fewer stereotypical images of scientists on a “draw-a-scientist” test and were able to generate more problem-solving strategies than students in the traditional group. Content knowledge learned was equivalent in both groups (Drake & Long, 2009). Urban middle school students engaged in a standards-based, inquiry-based science curriculum in ten middle schools showed higher levels of achievement on a curriculum-aligned test than students who received traditional instruction in a district-comparison group (Lynch, Kuipers, Pyke, & Szesze, 2005). Urban middle school students engaged in PBL showed increased academic performance in science and improved behavior ratings over a two-year period (Gordon, Rogers, Comfort, Gavula, & McGee, 2001). Urban students in grades seven and eight who were engaged in the LeTUS inquiry-based science curriculum demonstrated higher standardized test scores than students engaged in traditional instruction in a sample of 5,000 students. The LeTUs inquiry-science curriculum involves eight- to ten-week units addressing questions such as What Is the Quality of Air in My Community? or What Is the Water Like in My River? and is aligned with professional development, learning technology, and administrative support (Geier, Blumenfeld, Marx, Krajcik, Fishman, Soloway, & Clay-Chambers, 2008). Middle school students engaged in Learning by Design (LBD) consistently outperformed students engaged in traditional instruction on tests of collaboration and metacognitive skills, such as checking work, designing fair tests, and explaining evidence. LBD students also learned science content as well as or better than students engaged in traditional learning methods, with the largest gains among economically disadvantaged students (Kolodner, Camp, Crismond, Fasse, Gray, Holbrook, Puntambekar, & Ryan, 2003). Middle school students who received a computer-enhanced PBL unit had a better understanding of science concepts and felt more confident about being successful learners (Liu, Hsieh, Cho, & Schallert, 2006). Tenth-grade earth science students who engaged in PBL earned higher scores on an achievement test as compared to students who received traditional instruction (Chang, 2001). High school students engaged in PBL in biology, chemistry, and earth science classes outscored their peers on 44 percent of the items on the National Assessment of Educational Progress science test during their twelfth-grade year (Schneider, Krajcik, Marx, & Soloway, 2002). Students in grades five and up in 11 school districts learned math problems through videotaped problems over a three-week period (The Adventures of Jasper Woodbury series). The PBL students showed improved competence in solving basic math word problems and planning skills and more positive attitudes toward math (Cognition and Technology Group at Vanderbilt, 1992). PBL increased learning of macroeconomics at the high school level, as compared with traditional classes, in a sample of 252 students at 11 high schools (Maxwell, Mergendoller & Bellisimo, 2005). A randomized, controlled trial in Arizona and California in 2007-08 examined the effects of a project-based economics curriculum developed by the Buck Institute for Education on student learning and problem-solving skills in a sample of 7,000 twelfth graders in 66 high schools. Seventy-six teachers received 40 hours of professional development in teaching economics with PBL instead of their normal professional development activities. Students who received PBL scored significantly higher on problem-solving skills and in their ability to apply knowledge to real-world economic challenges than students taught economics using traditional methods. Economics teachers who used the PBL approach reported greater satisfaction with the materials and methods, and no significant differences were detected between intervention and control-group teachers (Finkelstein, Hanson, Huang, Hirschman, & Huang, 2010). Four veteran teachers taught macroeconomics using PBL in one or two courses and traditional instruction in another course. 246 twelfth-grade students in 11 classes completed pre- and post-tests in macroeconomics. Results showed that PBL was more effective than traditional instruction for teaching macroeconomics concepts (Mergendoller, Maxwell & Bellisimo, 2006). History and U.S. Government Second graders from low-income backgrounds participated in two project-based units which integrated literacy and social studies. The outcomes on standards-based social studies and content literacy assessments indicated that the project-based learning curriculum virtually erased the achievement gap between second graders of high and low-socioeconomic backgrounds (Halvorsen, Duke, Burgar, Block, Strachan, Berka, & Brown, 2012). Students in grades four and five collaboratively researched primary and secondary sources to discover themes and reasons for human migration in the local region. PBL students showed improved reasoning and collaboration skills and increased knowledge of local history and communities (Wieseman & Cadwell, 2005). Eighth-grade groups created mini-documentaries about their interpretation of a time period in the 1800s, using state standards as the content guide and presenting their completed work in a public event. PBL increased students’ content knowledge and historical-research skills (Hernandez-Ramos & De La Paz, 2009). High school students using PBL in American studies performed as well on multiple-choice tests as students who received a traditional model of instruction, and they showed a deeper understanding of content (Gallagher & Stepien, 1996). In the Knowledge in Action Research Project, high school students learned Advanced Placement U.S. Government with PBL or traditional instruction. The PBL course consisted of a public-policy action proposal and four role-playing projects: designing democracy, simulating legislation, a Supreme Court case, and an election. The PBL students showed improved performance on a complex scenario test, measuring strategies for monitoring and influencing public policy, and performed as well as or better than traditionally-taught students on the AP U.S. Government test (Boss et al., 2011; Parker et al., 2013; Parker et al., 2011). Learn more about the study's results to date and its course design. Continue to the next section of the PBL research review, Avoiding Pitfalls.

http://www.edutopia.org/pbl-research-practices-disciplines

- Development & Aid - Economy & Trade - Human Rights - Global Governance - Civil Society Sunday, February 14, 2016 - A black community in the southern Brazilian state of Rio de Janeiro is trying to maintain its cultural heritage on 287 hectares granted to it by the government in 1999 as part of reparations to the descendants of slaves. “I would define ‘quilombos’ as resistance by black people, as their essence,” Vagner do Nascimento, president of the Association of Residents of Campinho da Independencia, tells IPS. Situated to the southwest of the city of Rio de Janeiro, 20 km from the town of Paraty in the middle of lush jungle that forms part of the Mata Atlântica (Atlantic forest), Campinho – as it is referred to by the people who live here – is one of the 3,524 quilombos distributed throughout Brazil, according to the Culture Ministry’s Palmares Cultural Foundation. However, independent sources say there are 1,500 more quilombos – which were originally remote villages or collections of villages founded by runaway slaves, often hidden in the jungle. In the quilombos, escaped slaves kept alive the cultures and lifestyles brought over from Africa. They also became bastions in the struggle for freedom. After slavery was abolished on May 13, 1888, many quilombos became villages, where people depended on subsistence agriculture and small-scale trade. Zumbi dos Palmares was a famous colonial period quilombo located in the Serra da Barriga mountains in what is today the northern state of Alagoas. Palmares, which defended its freedom for over a century and at one point was home to as many as 50,000 escaped slaves, became a symbol in the fight against slavery in Brazil. “People there lived and worked together, and consolidated their own values. That’s why Palmares is a really strong reference point for us, because here in Campinho the land is collectively owned, and we have collective forms of producing, of generating culture, of working,” says Nascimento. Campinho has a unique history. Its 80 or so families descend from just three slave women: Antonia, Marcelina and Luiza. And according to the history that was orally transmitted from generation to generation, they weren’t just “ordinary” slaves, but came from the “Big House” and had culture and education. As the story goes, shortly after the abolition of slavery, the landowners of the three haciendas in the region distributed the property to their former slaves and left. Antonia, Marcelina and Luiza “gathered all the dispersed slaves together and brought them here with them,” says a member of the fifth generation of descendants, Laura María dos Santos, Campinho’s head of educational and cultural projects. Dos Santos, two young women named Daniele and Silvia, and the elderly Albertina do Nascimento head IPS’s welcoming committee – three generations of women who represent the strength of those who founded their community. “This heritage is transmitted to our girls, who become women who know what their role is,” says dos Santos. She goes on to tell an anecdote: “When a man made a sexist remark, his niece, a girl from our community, said ‘uncle, in a woman’s land, the woman never dies’.” Nor do the women in this community want their cultural heritage to die. The local residents association led by Nascimento is carrying out projects for the recuperation of the historical memory, craftsmaking, ethnic tourism, and the revival and sustainable production of traditional crops like cassava, rice, beans and corn. Campinho was the first quilombo in the state of Rio de Janeiro to obtain a formal collective land ownership title, on Mar. 21, 1999, after a struggle that began in the 1960s. On one hand, the creation of the Bocaína National Park kept them from hunting and collecting fruits in the forest – activities they depended on for a living. In addition, the construction of the Rio de Janeiro-to-Santos stretch of the BR-101 highway, between 1970 and 1973, drove up land prices and led to property speculation in the area. The entire Paraty region became the target of interest on the part of large tourism ventures, and many people began to sell and leave their lands. Those who stayed, like the community of Campinho, eventually won their battle. But other battles emerged, like the struggle for coexistence with a new touristy world of “rich people,” and the effort to preserve the local culture, says dos Santos. “It’s a question of continuing to fight for our ethnic and traditional identity, while at the same time incorporating the technology to which young people have a right,” she adds. On the other hand, as “Auntie Albertina” adds, there are other problems, like the fact that many of the local families have ended up working in nearby tourist condominiums. “Soon, no one is going to want to work the land any more,” laments the old woman, who says she would never give up her land, where she grows beans and rice, makes cachaça – a Brazilian liquor made from distilled sugar cane juice – and is her own boss. The haciendas in the region produced homemade cachaça since the time of the three original slave women. Referring to the women working outside the community, Silvia says “they see the rich women with their styles and want to imitate them. They have a nice house, but when they see their boss’s home, they start to suffer, because they want one just like it.” Another challenge is adopting new technologies, to which the children and young people in Campinho have access, generally in internet cafés, without forgetting their culture and traditional values. According to Silvia, “there’s no formula for how to do that,” other than raising awareness in the community about the importance of keeping the culture alive. The women see access to technologies like the internet, community radio stations and video cameras as important, to be able to record their history and culture. But the challenge is “to dominate technology, and not let technology dominate us.” “We teach our youngsters that the worst sin is to let themselves be enslaved by anything,” she says. Silvia moved away from the community when she was a little girl, and later became a community leader in a favela in Rio de Janeiro, where she went to live, like so many other quilombolas. Her fear is that the quilombos will go down the road followed by the favelas and become mere shantytowns without proper infrastructure or sanitation, due to a lack of space to expand. As the local families grow, and the land is increasingly subdivided, the houses are more and more crowded together. She recalls that many of the crowded favelas surrounding Rio and other cities today used to be areas similar to the quilombos, with gardens, open areas and natural water sources. That is why she wants the government to grant more land to the descendants of the original escaped slaves, who argue that they have a right to the property. “The houses have to be spaced out, so that if a woman fights with her husband, no one can hear,” Silvia jokes. Campinho is also involved in ecological projects. “Although some environmentalists say the opposite, the best-preserved areas are the ones where the quilombolas live,” she argues. The projects include the sustainable production of palm hearts and agroforestry. “First we feed the community, and later, we sell what’s left over, outside,” says Daniele. “Auntie Albertina”, who the younger people in the community listen to with respect, speaks up again. “On my plot of land I don’t let anyone kill even a little bird,” she says, recalling that in the past, her ancestors even ate toucans to survive, but pointing out that the birds have made a comeback in the surrounding jungle, which she calls “a great achievement.” The people of Campinho have other collective subsistence activities, such as the sale of traditional crafts in tourist areas of Paraty, and a restaurant that specialises in typical Afro-Brazilian dishes, like “feijoada”, a stew of black beans, pork and beef, which was traditionally made with what was left over after the master’s family was served. The community also has an ethnic tourism project, giving tours of the community that include visits to the old “senzala” or slave compound, the cassava flour production areas, and the community garden, as well as hikes in the jungle and cultural activities like traditional dance recitals. Nov. 20 is National Black Awareness Day in Brazil, in commemoration of the Nov. 20, 1695 death of the leader of black resistance in Brazil, Zumbi dos Palmares In some Brazilian states, including Rio de Janeiro and Sao Paulo, Nov. 20 is a holiday aimed at “reflection on the insertion of blacks in Brazilian society.” With blacks officially making up half of the population of Brazil, the country has the second largest black population in the world, after Nigeria. For Vagner do Nascimento, the leader of the community, there is no doubt: the question is to “survive with dignity. That is the essence of black people in Brazil today, whose ancestors were brought over from Africa.”

http://www.ipsnews.net/2009/11/brazil-quilombos-keep-black-cultural-identity-alive/

IN THIS ARTICLE How the Virus Is Spread Coxsackievirus is spread from person to person. The virus is present in the secretions and bodily fluids of infected people. The virus may be spread by coming into contact with respiratory secretions from infected patients. If infected people rub their runny noses and then touch a surface, that surface can harbor the virus and become a source of infection. The infection is spread when another person touches the contaminated surface and then touches his or her mouth or nose. People who have infected eyes (conjunctivitis) can spread the virus by touching their eyes and touching other people or touching a surface. Conjunctivitis may spread rapidly and appear within one day of exposure to the virus. Coxsackieviruses are also shed in stool, which may be a source of transmission among children. The virus can be spread if unwashed hands get contaminated with fecal matter and then touch the face. This is particularly important for spread within day-care centers or nurseries where diapers are handled. Medically Reviewed by a Doctor on 8/13/2015 Must Read Articles Related to Coxsackievirus Infectious Disease Resources - Early Care for Your Premature Baby - What to Eat When You Have Cancer - When to Take More Pain Medication

http://www.emedicinehealth.com/coxsackievirus/page2_em.htm

|This article does not cite any sources. (December 2009)| Carbonate rocks are a class of sedimentary rocks composed primarily of carbonate minerals. The two major types are limestone, which is composed of calcite or aragonite (different crystal forms of CaCO3) and dolostone, which is composed of the mineral dolomite (CaMg(CO3)2). Calcite can be either dissolved by groundwater or precipitated by groundwater, depending on several factors including the water temperature, pH, and dissolved ion concentrations. Calcite exhibits an unusual characteristic called retrograde solubility in which it becomes less soluble in water as the temperature increases. When conditions are right for precipitation, calcite forms mineral coatings that cement the existing rock grains together or it can fill fractures. Karst topography and caves develop in carbonate rocks because of their solubility in dilute acidic groundwater. Cooling groundwater or mixing of different groundwaters will also create conditions suitable for cave formation. |This article related to petrology is a stub. You can help Wikipedia by expanding it.|

https://en.wikipedia.org/wiki/Carbonate_rock

Skip to main content Wikispaces Classroom is now free, social, and easier than ever. Try it today. Pages and Files Ch 1 - Basics Ch 2 - Atoms, Molecules, and Ions Ch 3 - Mass Relationships in Chemical Reactions Ch 4 - Reactions in Aqueous Solutions Ch 5 - Gases Ch 6 - Thermochemistry Ch 7 - Quantum Theory Ch 8 - Periodic Relationships Ch 9 - Chemical Bonding I Ch_10 - Chemical Bonding II Ch_11 - Intermolecular Forces & Liquids and Solids Ch_12 - Physical Properties of Solutions Ch_13 - Chemical Kinetics Ch_14 - Chemical Equilibrium Ch_15 - Acids and Bases Ch_16 - Additional Equilibrium Topics Free Response Questions Add "All Pages" Ch_14 - Chemical Equilibrium Chapter 14 - Chemical Equilibrium In this chapter we will see how the kinetics of a reaction can lead to a state of chemical equilibrium where there is no observable change in the reaction. This means that the concentrations of the reactants and products are not change and therefore remain constant. We will learn how to show this state by writing equilibrium expressions and calculating the equilibrium constant. Related to the equilibrium expression is the reaction quotient which can tell us if a reaction is moving towards forming more reactants or products at any given time. Finally, we will focus on what is called an ICE table as the primary problem solving method for solving for concentration values at equilibrium. We will immediately look at some applications of equilibrium in Chapter 16, and then go back to Chapter 15. The Concept of Equilibrium and the Equilibrium Constant Writing Equilibrium Constant Expressions The Relationship Between Chemical Kinetics and Chemical Equilibrium What does the Equilibrium Constant Tell Us? Factors that Affect Chemical Equilibrium Along with the embedded videos, I have included some links to some tutorials, simulations and animations. Pay attention to my worked out Practice Problems and follow along with what we are doing in class! The listed problems from the book will be due the day of the Chapter 14/16 Test. This material will be covered very quickly so keep up. Practice Problems:14.1, 14.3, 14.6, 14.8, 14.9, 14.10, 14.11, 14.15, 14.16, 14.17, 14.18, 14.19, 14.22, 14.29, 14.30, 14.31, 14.32, 14.77, 14.33, 14.36, 14.37, 14.38, 14.40, 14.41, 14.43, 14.44, 14.48, 14.82, 14.51, 14.52, 14.53, 14.54, 14.55, 14.56, 14.58, 14.59, 14.98 14.1 - The Concept of Equilibrium and the Equilibrium Constant We have come across the term equilibrium several times already this year. Up to this point you probably have the basic thought that equilibrium means that two processes are occurring at the same time, and while this is a very generalized definition, we need something a lot more specific in order to clear up some common misconceptions. The first thing to forever remember about chemical equilibrium is that it means that the forward and reverse rates of a chemical reaction are exactly the same. Exactly. The same. The second is that the concentration of the reactants and products do not change at equilibrium and instead remain constant. Constant. This does NOT mean however the concentration of the reactants and products is exactly the same, absolutely not. Nope. Constant, not equal. Sear these two things into you brain forever. You can see what equilibrium looks like in the following graph of the formation of ammonia from hydrogen and nitrogen gases. It shows the disappearance of reactants (H2 & N2) and appearance of products (NH3) over time. The fact that the lines level off is what is meant by constant concentration. When this occurs, that is when equilibrium is achieved. So in the zone from A to B, the reaction is occurring and products are being made, but once you reach B and continue on through C and D, the concentrations are stable and a state of equilibrium exists. The equilibrium constant and expression are what govern a chemical reaction at equilibrium. Unlike writing rate laws, the stoichiometric coefficients in equilibrium reactions do play a part. They become the exponents on the reactants and products in the equilibrium expression. The equilibrium expression is always listed as the products over the reactants as you can see below. This is simply the mathematical expression of the law of mass action. Finally, the size of the equilibrium constant tells you a lot about the chemical reaction. If you have the exact same amount of reactants and products then mathematically the equilibrium expression must equal 1. If you have more products than reactants then the numerator will be larger and the equilibrium constant will be more than 1. If you have more reactants than products then the denominator will be larger and the equilibrium constant will be less than 1. Related Problems - 14.1, 14.3 14.2 - Writing Equilibrium Constant Expressions We will focus a lot of time now on writing equilibrium expressions as they are both fairly easy, and frequently asked for. The first thing to note is that there are reactions where everything is in the same phase (homogeneous) and those where the reactants and products are in different phases (heterogeneous). These cannot exactly be treated the same way, as they won’t have the same effect on the equilibrium. What we find is that gases and aqueous solutions greatly affect the equilibrium expression, but solid substances and pure liquids do not affect the equilibrium expression really at all. So when you look at a chemical reaction and need to write the equilibrium expression, you can ignore anything that is listed as a solid or a liquid. Writing Equilibrium Expressions So that leaves us with aqueous solutions and gases. Both of these can be written into an equilibrium expression based on concentration called Kc. These must have units in molarity for the concentration equilibrium expression. Unless specified otherwise, an equilibrium expression is always considered to be a concentration equilibrium expression, Kc. For reactions that are all expressed as gases, you can also write an equilibrium expression in units of pressure. So when you have units of atm or mm Hg then you have a pressure equilibrium expression called Kp. Now an equation that is all gases can be expressed as wither a Kc or a Kp depending on whether you use molarity or atm. Surprisingly, for the exact same reaction, the Kc and Kp will not be the same number so you must be certain which thing they are asking for. The two are related however and can be converted using equation 14.5 on pg 606. Ex 14.1 Prac Ex.png Ex 14.2 Prac Ex.png Ex 14.3 Prac Ex.png Ex 14.4 Prac Ex.png Ex 14.5 Prac Ex.png Ex 14.6 Prac Ex.png Converting between Kc and Kp Finally, we have seen over several chapters how adding chemical reactions together can lead to new ways to show information. Adding chemical reactions at equilibrium allows you to also combine the equilibrium expressions. To get the new equilibrium expression all you have to do is multiply the original two equilibrium expressions. Subtracting reactions means dividing equilibrium expressions. Ex 14.7 Prac Ex.png Related Problems - 14.6, 14.8, 14.9, 14.10, 14.11, 14.15, 14.16, 14.17, 14.18, 14.19, 14.22, 14.29, 14.30, 14.31, 14.32, 14.77 14.3 - The Relationship Between Chemical Kinetics and Chemical Equilibrium Remembering what we started off saying, the rate of the forward and reverse reactions at equilibrium are the same. So in theory you could set the two rate laws equal to each other. By rearranging the constants on to the same side, you can see that the ratio of the rate constants looks an awful lot like the equilibrium expressions we just finished writing. In fact, the equilibrium constant is simply the ratio of the two rate constants. Since the rate constant is dependent on temperature as we found in chapter 13, it makes sense that the equilibrium constant is also temperature dependent. So in order to write an equilibrium expression and calculate the equilibrium constant, the correctly balanced chemical equation and the temperature must be known. Related Problems - 14.33, 14.36 14.4 - What does the Equilibrium Constant Tell Us? We can use the equilibrium expression in several different ways. The most obvious one is to plug in concentration or pressure values and calculate the equilibrium constant. Make sure you take any exponents into account in this calculation. The second easiest thing to do is to calculate an unknown reactant or product concentration at equilibrium when the equilibrium constant and other concentrations are known. This should also be a simple calculation. As stated earlier, the value of the equilibrium constant can tell you if reactants or products are favored in the equilibrium. But how do you know if you are at equilibrium? As we saw in the graphs in the first section, it can take a while for a reaction to reach equilibrium. If you plug in the initial concentrations of reactants and products into the equilibrium expression you can calculate something called the reaction quotient, Q. The size of the reaction quotient in comparison to the equilibrium constant will tell you if you are at equilibrium, creating products, or creating reactants. If Q = K then you are at equilibrium. If Q < K then you haven’t reached equilibrium yet and still need to create more products, moving the reaction to the right. If Q > K then you have passed equilibrium and need to go back to reactants, moving the reaction to the left. Ex 14.8 Prac Ex.png The most common thing you will do with an equilibrium expression is to calculate the concentration values of reactants and products given the equilibrium constant and some initial concentration values. All problems of this nature are solved the same way using something typically called an “ICE” table. ICE stands for initial, change, and equilibrium and allows you to account for all of the possible changes in an equilibrium reaction. ICE Table example The common scenario is to know the initial concentration of reactants and assume that there are no products in existence. This means that any change in this condition will be to reduce the concentration of reactants and add to the concentration of the products. We always make this change equal x. We use the stoichiometric coefficients from the balanced equation as coefficients on the x. Then you can just use basic algebra to plug values into the equilibrium expression and solve for x. This is the calculation you will see most often on the free response portion of the exam and the one we will practice the most often. Ex 14.9 Prac Ex.png Related Problems - 14.37, 14.38, 14.40, 14.41, 14.43, 14.44, 14.48, 14.82 14.5 - Factors that Affect Chemical Equilibrium Work through this quick tutorial about a chicken (I'm not kidding) and stresses to its breathing to help you understand how "stress" affects the direction of equilibrium: A reaction at equilibrium is something dynamic and constant movement. It is not static. As such, it can change. But something must be done to it to make it change. We term these things “stresses” and a chemical reaction at equilibrium will react to stresses according to Le Chatlier’s principle. Le Chatlier’s Principle basically states that a chemical reaction at equilibrium will react to a stress in such a way to reduce that stress and establish a new equilibrium state. It is important to realize that the applied stress prevents a reaction from achieving the same equilibrium state, it will always be a new equilibrium state. Our most common types of stress include changes in concentration, pressure, volume, and temperature. Ex 14.11 Prac Ex.png For changes in concentration, the first thing to remember is that changes in concentration only matter if the thing that is being changed is represented in the equilibrium expression. So changes in concentration of solids and liquids DO NOT affect the equilibrium. For everything else, you need to see if the changed item is a reactant or product. Increasing a product will shift a reaction to the reactants so the extra product gets used up. The opposite is true if more reactant is used. The addition of something shifts the reaction in such a way to use it up. You can also take something out of a reaction in which case the reaction shifts to replace it. So you can keep a reaction continually producing products if you create a way to keep removing one of the products. Work through this LC Virtual Lab #1 to see these types of changes. Changing volume and pressure only affects gaseous reactants and products. In general we think of it as changes in pressure, but remember from chapter 5 that changing the volume will change the pressure. Increases in pressure will cause a reaction to shift towards the side with fewer moles of gas. If there are equal amounts of moles of gas, then there is no change with changes in pressure. Decreases in pressure will shift a reaction to the side with fewer moles of gas. Remember that a decrease in volume is an increase in pressure and vice versa. Please note that the addition of an inert gas such as helium is a favorite thing to ask about. Adding this inert gas at constant volume may increase the overall pressure but not the partial pressures of the individual gases so no change occurs. Ex 14.12 Prac Ex.png Work through this LC Virtual Lab #2 to see these types of changes. Finally, changes in temperature depend on whether a reaction is endo or exothermic. If a reaction is endothermic that means that energy is added to the reaction and we can think of heat/energy as a reactant. If a reaction is exothermic that means that energy is released from the reaction and we can think of heat/energy as a product. Then you can treat changes in temperature as simply changes in the concentration of heat/energy and it works the same way as originally discussed. If you don’t know if a reaction is endo or exothermic, then you cannot determine how changes in temperature will affect the equilibrium. But ANY change in temperature will change the value of the equilibrium constant. Ex 14.13 Prac Ex.png Work through this LC Virtual Lab #3 to see these types of changes. Finally, we need to explore how a catalyst affects equilibrium. As we found in chapter 14, the addition of a catalyst lowers the activation energy for a reaction making the rate of both the forward and reverse reactions faster. Since both the forward and reverse reactions speed up, there is actually no change in the equilibrium constant or the position of equilibrium. The only thing that does happen is that the state of equilibrium is achieved faster than would normally happen without the catalyst. Le Chatelier's Principle Related Problems - 14.51, 14.52, 14.53, 14.54, 14.55, 14.56, 14.58, 14.59, 14.98 help on how to format text Turn off "Getting Started"

http://baskinapchem.wikispaces.com/Ch_14%20-%20Chemical%20Equilibrium?responseToken=69e13420700c5925069bccb83dd9c8a4

Oxygen saturation (medicine) Oxygen saturation is a term referring to the fraction of oxygen-saturated hemoglobin relative to total hemoglobin (unsaturated + saturated) in the blood. The human body requires and regulates a very precise and specific balance of oxygen in the blood. Normal blood oxygen levels in humans are considered 95-100 percent. If the level is below 90 percent, it is considered low resulting in hypoxemia. Blood oxygen levels below 80 percent may compromise organ function, such as the brain and heart, and should be promptly addressed. Continued low oxygen levels may lead to respiratory or cardiac arrest. Oxygen therapy may be used to assist in raising blood oxygen levels. Oxygenation occurs when oxygen molecules (O 2) enter the tissues of the body. For example, blood is oxygenated in the lungs, where oxygen molecules travel from the air and into the blood. Oxygenation is commonly used to refer to medical oxygen saturation. In medicine, oxygen saturation (SO2), commonly referred to as "sats," measures the percentage of hemoglobin binding sites in the bloodstream occupied by oxygen. At low partial pressures of oxygen, most hemoglobin is deoxygenated. At around 90% (the value varies according to the clinical context) oxygen saturation increases according to an oxygen-hemoglobin dissociation curve and approaches 100% at partial oxygen pressures of >10 kPa. A pulse oximeter relies on the light absorption characteristics of saturated hemoglobin to give an indication of oxygen saturation. The body maintains a stable level of oxygen saturation for the most part by chemical processes of aerobic metabolism associated with breathing. Using the respiratory system, red blood cells, specifically the hemoglobin, gather oxygen in the lungs and distribute it to the rest of the body. The needs of the body's blood oxygen may fluctuate such as during exercise when more oxygen is required or when living at higher altitudes. A blood cell is said to be "saturated" when carrying a normal amount of oxygen. Both too high and too low levels can have adverse effects on the body. |This section does not cite any sources. (March 2015)| An SaO2 (arterial oxygen saturation) value below 90% causes hypoxemia (which can also be caused by anemia). Hypoxemia due to low SaO is indicated by cyanosis. Oxygen saturation can be measured in different tissues: - Venous oxygen saturation (SvO2) is measured to see how much oxygen the body consumes. Under clinical treatment, a SvO2 below 60% indicates that the body is in lack of oxygen, and ischemic diseases occur. This measurement is often used under treatment with a heart-lung machine (extracorporeal circulation), and can give the perfusionist an idea of how much flow the patient needs to stay healthy. - Tissue oxygen saturation (StO2) can be measured by near infrared spectroscopy. Although the measurements are still widely discussed, they give an idea of tissue oxygenation in various conditions. - Peripheral capillary oxygen saturation (SpO2) is an estimation of the oxygen saturation level usually measured with a pulse oximeter device. It can be calculated with the pulse oximetry according to the following formula: Pulse oximetry is a method used to measure the concentration of oxygen in the blood. A small device that clips to the body (typically a finger but may be other areas), called a pulse oximeter, uses infrared light to estimate the amount of oxygen in the blood. The clip attaches to a reading meter by a wire to collect the data. Oxygen levels may also be checked through an arterial blood gas test (ABG), where blood taken from an artery is analysed for oxygen level, carbon dioxide level and acidity. Oxygen saturation taken with a pulse oximeter is often designated SpO2. |85% and above||No evidence of impairment| |65% and less||Impaired mental function on average| |55% and less||Loss of consciousness on average| Healthy individuals at sea level usually exhibit oxygen saturation values between 96% and 99%, and should be above 94%. At 1600 meters altitude (about one mile high) oxygen saturation should be above 92%. An SaO2 (arterial oxygen saturation) value below 90% causes hypoxemia (which can also be caused by anemia). Hypoxemia due to low SaO2 is indicated by cyanosis, but oxygen saturation does not directly reflect tissue oxygenation. The affinity of hemoglobin to oxygen may impair or enhance oxygen release at the tissue level. Oxygen is more readily released to the tissues (i.e., hemoglobin has a lower affinity for oxygen) when pH is decreased, body temperature is increased, arterial partial pressure of carbon dioxide (PaCO2) is increased, and 2,3-DPG levels (a byproduct of glucose metabolism also found in stored blood products) are increased. When the hemoglobin has greater affinity for oxygen, less is available to the tissues. Conditions such as increased pH, decreased temperature, decreased PaCO2, and decreased 2,3-DPG will increase oxygen binding to the hemoglobin and limit its release to the tissue. - "Hypoxemia (low blood oxygen)". Mayo Clinic. mayoclinic.com. Retrieved 6 June 2013. - Kenneth D. McClatchey (2002). Clinical Laboratory Medicine. Philadelphia: Lippincott Williams & Wilkins. p. 370. - "Understanding Blood Oxygen Levels at Rest". fitday.com. fitday.com. Retrieved 6 June 2013. - Elllison, Bronwyn. "NORMAL RANGE OF BLOOD OXYGEN LEVEL". Livestrong.com. Livestrong.com. Retrieved 6 June 2013. - "Your Oxygen Level" (PDF). The Ohio State University Wexner Medical Center. Retrieved 6 June 2013. - "SPO2". TheFreeDictionary.com. 1998–2008. Retrieved 2014-01-28. - Oxymoron: Our Love-Hate Relationship with Oxygen, By Mike McEvoy at Albany Medical College, New York. 11/14/2012 - "Normal oxygen level". National Jewish Health. MedHelp. Feb 23, 2009. Retrieved 2014-01-28. - Schutz, Oxygen Saturation Monitoring by Pulse Oximetry, 2001

https://en.wikipedia.org/wiki/Oxygenation_(medical)

February 10, 2016, Hardening of the Arteries Coronary artery disease (CAD), also called heart disease or ischemic heart disease, results from a complex process known as atherosclerosis (commonly called "hardening of the arteries"). In atherosclerosis, fatty deposits (plaques) of cholesterol and other cellular waste products build up in the inner linings of the heart’s arteries. This causes blockage of arteries (ischemia) and prevents oxygen-rich blood from reaching the heart. There are many steps in the process leading to atherosclerosis, some not fully understood. Cholesterol and Lipoproteins. The atherosclerosis process begins with cholesterol and sphere-shaped bodies called lipoproteins that transport cholesterol. - Cholesterol is a substance found in all animal cells and animal-based foods. It is critical for many functions, but under certain conditions cholesterol can be harmful. - The lipoproteins that transport cholesterol are referred to by their size. The most commonly known are low-density lipoproteins (LDL) and high density lipoproteins (HDL). LDL is often referred to as "bad" cholesterol; HDL is often called "good" cholesterol. Oxidation. The damaging process called oxidation is an important trigger in the atherosclerosis story. - Oxidation is a chemical process in the body caused by the release of unstable particles known as oxygen-free radicals. It is one of the normal processes in the body, but under certain conditions (such as exposure to cigarette smoke or other environment stresses) these free radicals are overproduced. - In excess amounts, they can be very dangerous, causing damaging inflammation and even affecting genetic material in cells. - In heart disease, free radicals are released in artery linings and oxidize low-density lipoproteins (LDL). The oxidized LDL is the basis for cholesterol build-up on the artery walls and damage leading to heart disease. Inflammatory Response. For the arteries to harden there must be a persistent reaction in the body that causes ongoing harm. Researchers now believe that this reaction is an immune process known as the inflammatory response. There is growing evidence that the inflammatory response may be present not only in local plaques in single arteries but also throughout the arteries leading to the heart. Blockage in the Arteries. Eventually these calcified (hardened) arteries become narrower (a condition known as stenosis). - As this narrowing and hardening process continues, blood flow slows, preventing sufficient oxygen-rich blood from reaching the heart muscles. - Such oxygen deprivation in vital cells is called ischemia. When it affects the coronary arteries, it causes injury to the tissues of the heart. - These narrow and inelastic arteries not only slow down blood flow but also become vulnerable to injury and tears. The End Result: Heart Attack. A heart attack can occur as a result of one or two effects of atherosclerosis: - The artery becomes completely blocked and ischemia becomes so extensive that oxygen-bearing tissues around the heart die. - The plaque itself develops fissures or tears. Blood platelets stick to the site to seal off the plaque, and a blood clot (thrombus) forms. A heart attack can then occur if the blood clot completely blocks the passage of oxygen-rich blood to the heart. MOST POPULAR - HEALTH - Well: Why We Get Running Injuries (and How to Prevent Them) - Well: Simple Remedies for Constipation - Well: To Reduce the Risk of Alzheimer’s, Eat Fish - Study in Brazil Links Zika Virus to Eye Damage in Babies - Well: 12 Minutes of Yoga for Bone Health - Well: Ask Well: The Sugar in Fruit - Well: Opening Up About Depression - Johns Hopkins to Perform First H.I.V.-Positive Organ Transplants in U.S. - Well: A Diet and Exercise Plan to Lose Weight and Gain Muscle - How a Medical Mystery in Brazil Led Doctors to Zika

http://www.nytimes.com/health/guides/disease/atherosclerosis/background.html

First the bad news. Humans are driving species to extinction at around 1000 times the natural rate, at the top of the range of an earlier estimate. We also don’t know how many species we can afford to lose. Now the good news. Armed with your smartphone, you can help conservationists save them. Interactive map: “Where the threatened wild things are” The new estimate of the global rate of extinction comes from Stuart Pimm of Duke University in Durham, North Carolina, and colleagues. It updates a calculation Pimm’s team released in 1995, that human activities were driving species out of existence at 100 to 1000 times the background rate (Science, doi.org/fq2sfs). It turns out that Pimm’s earlier calculations both underestimated the rate at which species are now disappearing, and overestimated the background rate over the past 10 to 20 million years. Gone gone gone The Red List assessments of endangered species, conducted by the International Union for Conservation of Nature (IUCN), are key to Pimm’s analysis. They have evolved from patchy lists of threatened species into comprehensive surveys of animal groups and regions. “Twenty years ago we simply didn’t have the breadth of underlying data with 70,000 species assessments in hand,” says team member Thomas Brooks of the IUCN in Gland, Switzerland. By studying animals’ DNA, biologists have also created family trees for many groups of animals, allowing them to calculate when new species emerged. On average, it seems each vertebrate species gives rise to a new species once every 10 million years. It’s hard to measure the natural rate of extinction, but there is a workaround. Before we started destroying habitats, new species seem to have been appearing faster than old ones disappeared. That means the natural extinction rate cannot be higher than the rate at which they were forming, says Pimm. For the most part, the higher estimate of the modern extinction rate is not caused by any acceleration in extinctions since 1995. One exception is an increase in threats to amphibians, partly due to the global spread of the killer chytrid fungus. The big unknown is what the high current extinction rate means for the health of entire ecosystems. Some researchers have suggested “sustainable” targets for species’ loss, but there’s still no scientific way to predict at what point cumulative extinctions cause an ecosystem to collapse. “People who say that are pulling numbers out of the air,” says Pimm. Still, it seems unlikely that extinctions running at 1000 times the background rate can be sustained for long. “You can be sure that there will be a price to be paid,” says Brooks. Pimm’s team has also compiled detailed global maps of biodiversity, showing the numbers of threatened species and total species richness in a global grid consisting of squares 10 kilometres across. Such maps can help conservationists decide what to do. For instance, Pimm and his colleague Clinton Jenkins of the Institute for Ecological Research in Nazaré Paulista, Brazil, noticed high numbers of threatened species on Brazil’s Atlantic coast. Local forests were being cleared for cattle ranching. So they are working with a Brazilian group, the Golden Lion Tamarin Association, to buy land and reconnect isolated forest fragments. But conservationists need more data, and you can help, through projects like iNaturalist. Users share photos of the creatures they see via iPhone and Android apps, and experts identify them. “Right now, someone is posting an observation about every 30 seconds,” says co-director Scott Loarie of the California Academy of Sciences in San Francisco. Journal reference: Science, DOI: 10.1126/science.1246752 More on these topics:

https://www.newscientist.com/article/dn25645-we-are-killing-species-at-1000-times-the-natural-rate

September 7, 2005 Why Are Birds’ Eggs Speckled? Birds' eggs are unique in their diverse pigmentation. This diversity is greatest amongst perching birds (order Passeriformes: 60% of all bird species), which include many familiar species including tits and warblers. Despite intense interest, the purpose, in most species, of these patterns was unknown. Most passerines lay eggs speckled with reddish protoporphyrin spots forming a ring around the egg's blunt end, on an otherwise unpigmented shell. Evidence in a paper by Gosler, Higham & Reynolds soon to appear in Ecology Letters now suggests that rather than giving a visual signal, protoporphyrins strengthen the eggshell by compensating for reduced eggshell-thickness caused by calcium deficiency.Pigment spots on great tit eggs specifically marked thinner areas of shell, with darker spots marking yet thinner shell than paler spots, and females nesting on low-calcium soils, laid thinner-shelled, more-spotted eggs than those on high-calcium soils nearby. Pigmentation may offer a way to assess eggshell quality. On the World Wide Web:

http://www.redorbit.com/news/science/233301/why_are_birds_eggs_speckled/

ELA: KINDERGARTEN - GRADE 12 LITERACY: GRADES 6 - 12 College and Career Readiness Anchor Standards for Language Vocabulary Acquisition and Use 4. Determine or clarify the meaning of unknown and multiple-meaning words and phrases by using context clues, analyzing meaningful word parts, and consulting general and specialized reference materials, as appropriate. Toddlers and young children learn vocabulary without thinking about it—their brains simply absorb new words and new meanings for familiar words with little to no effort on their part. Would that we all were so lucky! As humans age, we lose the sponge-like ability to soak up new words and meanings—even though, in a globalized world, we need those skills more than ever. Luckily, we never completely lose the ability to learn vocabulary. Students preparing for college and/or a career should practice the skills they’ll need to decipher the forest of new words they’ll inevitably encounter. Specifically, the Anchor Standards recommend focusing on the following skills: Understanding words in context. The sentence or paragraph a new word lives in can provide plenty of information about what it means. Often, the sentence alone tells the reader enough to give him a good shot at guessing what the word means. Context is especially important when a single word has multiple meanings. For instance, the word “glass” may mean a container for a beverage, a flat pane of material used in a window or mirror, or the substance from which both of these are made. A sentence like “I raised my glass to toast their happy marriage,” however, immediately brings to mind the drink-holder type of “glass,” not the used-in-windows type of “glass.” (Similarly, “toast” in this context brings to mind a type of speech, not a piece of cooked bread.) Examining word parts. English may read like it picks the pockets of other languages for spare vocabulary, but many English words, especially those used in technical or professional fields, use recognizable word parts that give a clear view of what the word means. Perhaps the best-known example is the suffix “-ology,” which means “the study of.” Using reference materials. When all else fails, look it up! Most students preparing for college or a career are familiar with the basic reference trio of dictionary/thesaurus/encyclopedia, but most specialties have their own specific reference materials along with the tried-and-true favorites. P.S. If your students need to brush up on their spelling and grammar, send 'em over to our Grammar Learning Guides so they can hone their skills before conquering the Common Core. Sample Activities for Use in Class Reading Outside Your Sphere: This activity can be done in an hour, or it can serve as an ongoing semester project. Students will need access to magazines, books, and other reading materials. For a one-day assignment, the school library may suffice; for an entire semester, you may want to have students subscribe to a magazine, or use resources at a local public or university library, if available. Assign, or have each student choose, a topic or area of study about which the student knows little to nothing. More than one student may be assigned to each topic, if necessary. For the duration of the assignment, have the students read a magazine, newsletter, book, or other publication in their unknown topic. As they encounter words they don’t know, students should write down: 1. The word; 2. The sentence the word appears in; 3. Their best guess as to what the word means and what they base that guess on - the context, a definition or example given in the text, the parts of the word, etc. (one or two sentences will suffice); and 4. What the word actually means and where they learned that information (dictionary, website, asking a professional in the field, etc.) Reading topics they know nothing about will not only expose students to vocabulary they’ve never seen before, but also challenges them to decide what resources are best for finding out. For this activity, you’ll need a stack of general and specialized reference materials and a stack of cards containing vocabulary words, phrases and concepts that might be found in these reference works. For instance, if your pile of reference materials contains a medical dictionary and a copy of Gray’s Anatomy, your cards should include a few medical terms, such as the technical name for certain organs, diseases, or medical procedures. Divide the class into two or more teams of four to seven students each. Put the stack of reference materials on a table at the front of the room, and divide the cards into one stack per team and set them at the front of the room either with the reference materials or on their own table. (For scoring purposes, it may be easier to color-code the cards for each team.) Students should line up in single-file lines, one per team, facing the reference materials and stacks of cards. On “go!” (or some similar signal to begin), the first student in each line will race to the front of the room and grab the top card off her team’s stack. The student reads the card, then has to decide as quickly as possible which reference materials are most likely to have the definition of the word, phrase or concept on the card. The student should go through the reference material(s) until she finds the definition, then mark that page in the book with her card and race to the back of her team’s line, at which point the second student on the team runs forward and does the same thing. Once everyone on the team has stuck a card in a book, the entire team should raise its hands. A team has “won” if all its cards bookmark a page that defines the word on the card. Read the words and the definitions out loud, or have each student read his or hers out loud.

http://www.shmoop.com/common-core-standards/ccss-ela-literacy-ccra-l-4.html

By Steve Whitmoyer You need to be logged in to use this feature. Log In or Register Learners identify the parts of speech by following a certain order until each word in a sentence is labeled. In a variety of exercises, learners practice finding verbs, prepositional phrases, subjects, nouns, pronouns, adjectives, adverbs, and conjunctions. In this animated object, learners examine neutral fats, phospholipids, and cholesterol. The molecular formula and general function for each are shown. Students examine standard pressure in this interactive object. You'll practice converting between units of measure for volume in the English Measurement System. Spin to Win with conduction, convection, and radiation terminology!

https://www.wisc-online.com/learn/career-clusters/stem/eng3502/weight--volume-relationships-saturated-densit

This text was copied from Wikipedia on 29 November 2015 at 3:22PM. The slide rule, also known colloquially in the United States as a slipstick, is a mechanical analog computer. The slide rule is used primarily for multiplication and division, and also for functions such as roots, logarithms and trigonometry, but is not normally used for addition or subtraction. Though similar in name and appearance to a standard ruler, the slide rule is not ordinarily used for measuring length or drawing straight lines. Slide rules exist in a diverse range of styles and generally appear in a linear or circular form with a standardized set of markings (scales) essential to performing mathematical computations. Slide rules manufactured for specialized fields such as aviation or finance typically feature additional scales that aid in calculations common to those fields. The Reverend William Oughtred and others developed the slide rule in the 17th century based on the emerging work on logarithms by John Napier. Before the advent of the pocket calculator, it was the most commonly used calculation tool in science and engineering. The use of slide rules continued to grow through the 1950s and 1960s even as digital computing devices were being gradually introduced; but around 1974 the electronic scientific calculator made it largely obsolete and most suppliers left the business. - 1 Basic concepts - 2 Operation - 3 Physical design - 4 History - 5 Compared to electronic digital calculators - 6 The slide rule today - 7 See also - 8 Notes - 9 External links In its most basic form, the slide rule uses two logarithmic scales to allow rapid multiplication and division of numbers. These common operations can be time-consuming and error-prone when done on paper. More elaborate slide rules allow other calculations, such as square roots, exponentials, logarithms, and trigonometric functions. Scales may be grouped in decades, which are numbers ranging from 1 to 10 (i.e. 10n to 10n+1). Thus single decade scales C and D range from 1 to 10 across the entire width of the slide rule while double decade scales A and B range from 1 to 100 over the width of the slide rule. In general, mathematical calculations are performed by aligning a mark on the sliding central strip with a mark on one of the fixed strips, and then observing the relative positions of other marks on the strips. Numbers aligned with the marks give the approximate value of the product, quotient, or other calculated result. The user determines the location of the decimal point in the result, based on mental estimation. Scientific notation is used to track the decimal point in more formal calculations. Addition and subtraction steps in a calculation are generally done mentally or on paper, not on the slide rule. Most slide rules consist of three linear strips of the same length, aligned in parallel and interlocked so that the central strip can be moved lengthwise relative to the other two. The outer two strips are fixed so that their relative positions do not change. Some slide rules ("duplex" models) have scales on both sides of the rule and slide strip, others on one side of the outer strips and both sides of the slide strip (which can usually be pulled out, flipped over and reinserted for convenience), still others on one side only ("simplex" rules). A sliding cursor with a vertical alignment line is used to find corresponding points on scales that are not adjacent to each other or, in duplex models, are on the other side of the rule. The cursor can also record an intermediate result on any of the scales. A logarithm transforms the operations of multiplication and division to addition and subtraction according to the rules and . Moving the top scale to the right by a distance of , by matching the beginning of the top scale with the label on the bottom, aligns each number , at position on the top scale, with the number at position on the bottom scale. Because , this position on the bottom scale gives , the product of and . For example, to calculate 3×2, the 1 on the top scale is moved to the 2 on the bottom scale. The answer, 6, is read off the bottom scale where 3 is on the top scale. In general, the 1 on the top is moved to a factor on the bottom, and the answer is read off the bottom where the other factor is on the top. This works because the distances from the "1" are proportional to the logarithms of the marked values: Operations may go "off the scale;" for example, the diagram above shows that the slide rule has not positioned the 7 on the upper scale above any number on the lower scale, so it does not give any answer for 2×7. In such cases, the user may slide the upper scale to the left until its right index aligns with the 2, effectively dividing by 10 (by subtracting the full length of the C-scale) and then multiplying by 7, as in the illustration below: Here the user of the slide rule must remember to adjust the decimal point appropriately to correct the final answer. We wanted to find 2×7, but instead we calculated (2/10)×7=0.2×7=1.4. So the true answer is not 1.4 but 14. Resetting the slide is not the only way to handle multiplications that would result in off-scale results, such as 2×7; some other methods are: - Use the double-decade scales A and B. - Use the folded scales. In this example, set the left 1 of C opposite the 2 of D. Move the cursor to 7 on CF, and read the result from DF. - Use the CI inverted scale. Position the 7 on the CI scale above the 2 on the D scale, and then read the result off of the D scale below the 1 on the CI scale. Since 1 occurs in two places on the CI scale, one of them will always be on-scale. - Use both the CI inverted scale and the C scale. Line up the 2 of CI with the 1 of D, and read the result from D, below the 7 on the C scale. - Using a circular slide rule. Method 1 is easy to understand, but entails a loss of precision. Method 3 has the advantage that it only involves two scales. The illustration below demonstrates the computation of 5.5/2. The 2 on the top scale is placed over the 5.5 on the bottom scale. The 1 on the top scale lies above the quotient, 2.75. There is more than one method for doing division, but the method presented here has the advantage that the final result cannot be off-scale, because one has a choice of using the 1 at either end. In addition to the logarithmic scales, some slide rules have other mathematical functions encoded on other auxiliary scales. The most popular were trigonometric, usually sine and tangent, common logarithm (log10) (for taking the log of a value on a multiplier scale), natural logarithm (ln) and exponential (ex) scales. Some rules include a Pythagorean scale, to figure sides of triangles, and a scale to figure circles. Others feature scales for calculating hyperbolic functions. On linear rules, the scales and their labeling are highly standardized, with variation usually occurring only in terms of which scales are included and in what order: |A, B||two-decade logarithmic scales, used for finding square roots and squares of numbers| |C, D||single-decade logarithmic scales| |K||three-decade logarithmic scale, used for finding cube roots and cubes of numbers| |CF, DF||"folded" versions of the C and D scales that start from π rather than from unity; these are convenient in two cases. First when the user guesses a product will be close to 10 but is not sure whether it will be slightly less or slightly more than 10, the folded scales avoid the possibility of going off the scale. Second, by making the start π rather than the square root of 10, multiplying or dividing by π (as is common in science and engineering formulas) is simplified.| |CI, DI, CIF, DIF||"inverted" scales, running from right to left, used to simplify 1/x steps| |S||used for finding sines and cosines on the C (or D) scale| |T, T1, T2||used for finding tangents and cotangents on the C and CI (or D and DI) scales| |ST, SRT||used for sines and tangents of small angles and degree–radian conversion| |L||a linear scale, used along with the C and D scales for finding base-10 logarithms and powers of 10| |LLn||a set of log-log scales, used for finding logarithms and exponentials of numbers| |Ln||a linear scale, used along with the C and D scales for finding natural (base e) logarithms and| |The scales on the front and back of a Keuffel and Esser (K&E) 4081-3 slide rule.| Roots and powers There are single-decade (C and D), double-decade (A and B), and triple-decade (K) scales. To compute , for example, locate x on the D scale and read its square on the A scale. Inverting this process allows square roots to be found, and similarly for the powers 3, 1/3, 2/3, and 3/2. Care must be taken when the base, x, is found in more than one place on its scale. For instance, there are two nines on the A scale; to find the square root of nine, use the first one; the second one gives the square root of 90. For problems, use the LL scales. When several LL scales are present, use the one with x on it. First, align the leftmost 1 on the C scale with x on the LL scale. Then, find y on the C scale and go down to the LL scale with x on it. That scale will indicate the answer. If y is "off the scale," locate and square it using the A and B scales as described above. Alternatively, use the rightmost 1 on the C scale, and read the answer off the next higher LL scale. For example, aligning the rightmost 1 on the C scale with 2 on the LL2 scale, 3 on the C scale lines up with 8 on the LL3 scale. The S, T, and ST scales are used for trig functions and multiples of trig functions, for angles in degrees. For angles from around 5.7 up to 90 degrees, sines are found by comparing the S scale with C (or D) scale; though on many closed-body rules the S scale relates to the A scale instead, and what follows must be adjusted appropriately. The S scale has a second set of angles (sometimes in a different color), which run in the opposite direction, and are used for cosines. Tangents are found by comparing the T scale with the C (or D) scale for angles less than 45 degrees. For angles greater than 45 degrees the CI scale is used. Common forms such as can be read directly from x on the S scale to the result on the D scale, when the C-scale index is set at k. For angles below 5.7 degrees, sines, tangents, and radians are approximately equal, and are found on the ST or SRT (sines, radians, and tangents) scale, or simply divided by 57.3 degrees/radian. Inverse trigonometric functions are found by reversing the process. Many slide rules have S, T, and ST scales marked with degrees and minutes (e.g. some Keuffel and Esser models, late-model Teledyne-Post Mannheim-type rules). So-called decitrig models use decimal fractions of degrees instead. Logarithms and exponentials Base-10 logarithms and exponentials are found using the L scale, which is linear. Some slide rules have a Ln scale, which is for base e. Logarithms to any other base can be calculated by reversing the procedure for calculating powers of a number. For example, log2 values can be determined by lining up either leftmost or rightmost 1 on the C scale with 2 on the LL2 scale, finding the number whose logarithm is to be calculated on the corresponding LL scale, and reading the log2 value on the C scale. Addition and subtraction Slide rules are not typically used for addition and subtraction, but it is nevertheless possible to do so using two different techniques. The first method to perform addition and subtraction on the C and D (or any comparable scales) requires converting the problem into one of division. For addition, the quotient of the two variables plus one times the divisor equals their sum: For subtraction, the quotient of the two variables minus one times the divisor equals their difference: This method is similar to the addition/subtraction technique used for high-speed electronic circuits with the logarithmic number system in specialized computer applications like the Gravity Pipe (GRAPE) supercomputer and hidden Markov models. The second method utilizes a sliding linear L scale available on some models. Addition and subtraction are performed by sliding the cursor left (for subtraction) or right (for addition) then returning the slide to 0 to read the result. Standard linear rules The width of the slide rule is quoted in terms of the nominal width of the scales. Scales on the most common "10-inch" models are actually 25 cm, as they were made to metric standards, though some rules offer slightly extended scales to simplify manipulation when a result overflowed. Pocket rules are typically 5 inches. Models a couple of metres wide were sold to be hung in classrooms for teaching purposes. Typically the divisions mark a scale to a precision of two significant figures, and the user estimates the third figure. Some high-end slide rules have magnifier cursors that make the markings easier to see. Such cursors can effectively double the accuracy of readings, permitting a 10-inch slide rule to serve as well as a 20-inch. Various other conveniences have been developed. Trigonometric scales are sometimes dual-labeled, in black and red, with complementary angles, the so-called "Darmstadt" style. Duplex slide rules often duplicate some of the scales on the back. Scales are often "split" to get higher accuracy. Circular slide rules Circular slide rules come in two basic types, one with two cursors (left), and another with a free dish and one cursor (right). The dual cursor versions perform multiplication and division by holding a fast angle between the cursors as they are rotated around the dial. The onefold cursor version operates more like the standard slide rule through the appropriate alignment of the scales. The basic advantage of a circular slide rule is that the widest dimension of the tool was reduced by a factor of about 3 (i.e. by π). For example, a 10 cm circular would have a maximum precision approximately equal to a 31.4 cm ordinary slide rule. Circular slide rules also eliminate "off-scale" calculations, because the scales were designed to "wrap around"; they never have to be reoriented when results are near 1.0—the rule is always on scale. However, for non-cyclical non-spiral scales such as S, T, and LL's, the scale width is narrowed to make room for end margins. Circular slide rules are mechanically more rugged and smoother-moving, but their scale alignment precision is sensitive to the centering of a central pivot; a minute 0.1 mm off-centre of the pivot can result in a 0.2mm worst case alignment error. The pivot, however, does prevent scratching of the face and cursors. The highest accuracy scales are placed on the outer rings. Rather than "split" scales, high-end circular rules use spiral scales for more complex operations like log-of-log scales. One eight-inch premium circular rule had a 50-inch spiral log-log scale. The main disadvantages of circular slide rules are the difficulty in locating figures along a dish, and limited number of scales. Another drawback of circular slide rules is that less-important scales are closer to the center, and have lower precisions. Most students learned slide rule use on the linear slide rules, and did not find reason to switch. One slide rule remaining in daily use around the world is the E6B. This is a circular slide rule first created in the 1930s for aircraft pilots to help with dead reckoning. With the aid of scales printed on the frame it also helps with such miscellaneous tasks as converting time, distance, speed, and temperature values, compass errors, and calculating fuel use. The so-called "prayer wheel" is still available in flight shops, and remains widely used. While GPS has reduced the use of dead reckoning for aerial navigation, and handheld calculators have taken over many of its functions, the E6B remains widely used as a primary or backup device and the majority of flight schools demand that their students have some degree of proficiency in its use. Proportion wheels are simple circular slide rules used in graphic design to broaden or slim images and photographs. Lining up the desired values on the emmer and inner wheels (which correspond to the original and desired sizes) will display the proportion as a percentage in a small window. They are not as common since the advent of computerized layout, but are still made and used. In 1952, Swiss watch company Breitling introduced a pilot's wristwatch with an integrated circular slide rule specialized for flight calculations: the Breitling Navitimer. The Navitimer circular rule, referred to by Breitling as a "navigation computer", featured airspeed, rate/time of climb/descent, flight time, distance, and fuel consumption functions, as well as kilometer—nautical mile and gallon—liter fuel amount conversion functions. A Russian circular slide rule built like a pocket watch that works as single cursor slide rule since the two needles are ganged together. Breitling Navitimer wristwatch with circular slide rule Cylindrical slide rules There are two main types of cylindrical slide rules: those with helical scales such as the Fuller, the Otis King and the Bygrave slide rule, and those with bars, such as the Thacher and some Loga models. In either case, the advantage is a much longer scale, and hence potentially greater precision, than afforded by a straight or circular rule. In 1895, a Japanese firm, Hemmi, started to make slide rules from bamboo, which had the advantages of being dimensionally stable, strong and naturally self-lubricating. These bamboo slide rules were introduced in Sweden in September, 1933, and probably only a little earlier in Germany. Scales were made of celluloid, plastic, or painted aluminium. Later cursors were acrylics or polycarbonates sliding on Teflon bearings. All premium slide rules had numbers and scales engraved, and then filled with paint or other resin. Painted or imprinted slide rules were viewed as inferior because the markings could wear off. Nevertheless, Pickett, probably America's most successful slide rule company, made all printed scales. Premium slide rules included clever catches so the rule would not fall apart by accident, and bumpers to protect the scales and cursor from rubbing on tabletops. The recommended cleaning method for engraved markings is to scrub lightly with steel-wool. For painted slide rules use diluted commercial window-cleaning fluid and a soft cloth. The slide rule was invented around 1620–1630, shortly after John Napier's publication of the concept of the logarithm. Edmund Gunter of Oxford developed a calculating device with a single logarithmic scale; with additional measuring tools it could be used to multiply and divide. The first description of this scale was published in Paris in 1624 by Edmund Wingate (c.1593–1656), an English mathematician, in a book entitled L'usage de la reigle de proportion en l'arithmetique & geometrie. The book contains a double scale, logarithmic on one side, tabular on the other. In 1630, William Oughtred of Cambridge invented a circular slide rule, and in 1632 combined two handheld Gunter rules to make a device that is recognizably the modern slide rule. Like his contemporary at Cambridge, Isaac Newton, Oughtred taught his ideas privately to his students. Also like Newton, he became involved in a vitriolic controversy over priority, with his one-time student Richard Delamain and the prior claims of Wingate. Oughtred's ideas were only made public in publications of his student William Forster in 1632 and 1653. In 1722, Warner introduced the two- and three-decade scales, and in 1755 Everard included an inverted scale; a slide rule containing all of these scales is usually known as a "polyphase" rule. In 1815, Peter Mark Roget invented the log log slide rule, which included a scale displaying the logarithm of the logarithm. This allowed the user to directly perform calculations involving roots and exponents. This was especially useful for fractional powers. In 1821, Nathaniel Bowditch, described in the American Practical Navigator a "sliding rule" that contained scales trigonometric functions on the fixed part and a line of log-sines and log-tans on the slider used to solve navigation problems. A more modern form of slide rule was created in 1859 by French artillery lieutenant Amédée Mannheim, "who was fortunate in having his rule made by a firm of national reputation and in having it adopted by the French Artillery." It was around this time that engineering became a recognized profession, resulting in widespread slide rule use in Europe–but not in the United States. There Edwin Thacher's cylindrical rule took hold after 1881. The duplex rule was invented by William Cox in 1891, and was produced by Keuffel and Esser Co. of New York. Astronomical work also required fine computations, and in 19th-century Germany a steel slide rule about 2 meters long was used at one observatory. It had a microscope attached, giving it accuracy to six decimal places. Throughout the 1950s and 1960s the slide rule was the symbol of the engineer's profession in the same way the stethoscope is of the medical profession's. German rocket scientist Wernher von Braun brought two 1930s vintage Nestler slide rules with him when he moved to the U.S. after World War 2 to work on the American space effort. Throughout his life he never used any other pocket calculating device, even while heading the NASA program that landed a man on the moon in 1969. Aluminium Pickett-brand slide rules were carried on Project Apollo space missions. The model N600-ES owned by Buzz Aldrin that flew with him to the moon on Apollo 11 was sold at auction in 2007. The model N600-ES taken along on Apollo 13 in 1970 is owned by the National Air and Space Museum. Some engineering students and engineers carried ten-inch slide rules in belt holsters, a common sight on campuses even into the mid-1970s. Until the advent of the pocket digital calculator students also might keep a ten- or twenty-inch rule for precision work at home or the office while carrying a five-inch pocket slide rule around with them. In 2004, education researchers David B. Sher and Dean C. Nataro conceived a new type of slide rule based on prosthaphaeresis, an algorithm for rapidly computing products that predates logarithms. However, there has been little practical interest in constructing one beyond the initial prototype. Slide rules have often been specialized to varying degrees for their field of use, such as excise, proof calculation, engineering, navigation, etc., but some slide rules are extremely specialized for very narrow applications. For example, the John Rabone & Sons 1892 catalog lists a "Measuring Tape and Cattle Gauge", a device to estimate the weight of a cow from its measurements. There were many specialized slide rules for photographic applications; for example, the actinograph of Hurter and Driffield was a two-slide boxwood, brass, and cardboard device for estimating exposure from time of day, time of year, and latitude. Specialized slide rules were invented for various forms of engineering, business and banking. These often had common calculations directly expressed as special scales, for example loan calculations, optimal purchase quantities, or particular engineering equations. For example, the Fisher Controls company distributed a customized slide rule adapted to solving the equations used for selecting the proper size of industrial flow control valves. In World War II, bombardiers and navigators who required quick calculations often used specialized slide rules. One office of the U.S. Navy actually designed a generic slide rule "chassis" with an aluminium body and plastic cursor into which celluloid cards (printed on both sides) could be placed for special calculations. The process was invented to calculate range, fuel use and altitude for aircraft, and then adapted to many other purposes. The E6-B is a circular slide rule used by pilots & navigators. The importance of the slide rule began to diminish as electronic computers, a new but rare resource in the 1950s, became more widely available to technical workers during the 1960s. (See History of computing hardware (1960s–present).) Computers also changed the nature of calculation. With slide rules a great emphasis was put on working the algebra to get expressions into the most computable form. Users would simply approximate or drop small terms to simplify a calculation. FORTRAN allowed complicated formulas to be typed in from textbooks without the effort of reformulation. Numerical integration was often easier than trying to find closed-form solutions for difficult problems. The young engineer asking for computer time to solve a problem that could have been done by a few swipes on the slide rule became a humorous cliché. The availability of mainframe computing did not however significantly affect the ubiquitous use of the slide rule until cheap hand held electronic calculators for scientific and engineering purposes became available in the mid-1970s, at which point it rapidly declined. The first included the Wang Laboratories LOCI-2, introduced in 1965, which used logarithms for multiplication and division and the Hewlett-Packard HP-9100, introduced in 1968. The HP-9100 had trigonometric functions (sin, cos, tan) in addition to exponentials and logarithms. It used the CORDIC (coordinate rotation digital computer) algorithm, which allows for calculation of trigonometric functions using only shift and add operations. This method facilitated the development of ever smaller scientific calculators. As calculator price declined geometrically and functionality increased exponentially the slide rule's fate was sealed. The pocket-sized Hewlett-Packard HP-35 scientific calculator cost US$395 in 1972, too expensive for most students. By 1975 basic four-function electronic calculators could be purchased for less than $50, and by 1976 the TI-30 scientific calculator could be purchased for less than $25. Compared to electronic digital calculators Most people find slide rules difficult to learn and use. Even during their heyday, they never caught on with the general public. Addition and subtraction are not well-supported operations on slide rules and doing a calculation on a slide rule tends to be slower than on a calculator. This led engineers to take mathematical shortcuts favoring operations that were easy on a slide rule, creating inaccuracies and mistakes. On the other hand, the spatial, manual operation of slide rules cultivates in the user an intuition for numerical relationships and scale that people who have used only digital calculators often lack. A slide rule will also display all the terms of a calculation along with the result, thus eliminating uncertainty about what calculation was actually performed. A slide rule requires the user to separately compute the order of magnitude of the answer in order to position the decimal point in the results. For example, 1.5 × 30 (which equals 45) will show the same result as 1,500,000 × 0.03 (which equals 45,000). This separate calculation is less likely to lead to extreme calculation errors, but forces the user to keep track of magnitude in short-term memory (which is error-prone), keep notes (which is cumbersome) or reason about it in every step (which distracts from the other calculation requirements). The typical precision of a slide rule is about three significant digits, compared to many digits on digital calculators. As order of magnitude gets the greatest prominence when using a slide rule, users are less likely to make errors of false precision. When performing a sequence of multiplications or divisions by the same number, the answer can often be determined by merely glancing at the slide rule without any manipulation. This can be especially useful when calculating percentages (e.g. for test scores) or when comparing prices (e.g. in dollars per kilogram). Multiple speed-time-distance calculations can be performed hands-free at a glance with a slide rule. Other useful linear conversions such as pounds to kilograms can be easily marked on the rule and used directly in calculations. Being entirely mechanical, a slide rule does not depend on electricity or batteries. However, mechanical imprecision in slide rules that were poorly constructed or warped by heat or use will lead to errors. Many sailors keep slide rules as backups for navigation in case of electric failure or battery depletion on long route segments. Slide rules are still commonly used in aviation, particularly for smaller planes. They are being replaced only by integrated, special purpose and expensive flight computers, and not general-purpose calculators. The E6B circular slide rule used by pilots has been in continuous production and remains available in a variety of models. Some wrist watches designed for aviation use still feature slide rule scales to permit quick calculations. The Citizen Skyhawk AT is a notable example. The slide rule today ||This section possibly contains original research. (February 2015)| Even today some people prefer a slide rule over an electronic calculator as a practical computing device. Others keep their old slide rules out of a sense of nostalgia, or collect them as a hobby. A popular collectible model is the Keuffel & Esser Deci-Lon, a premium scientific and engineering slide rule available both in a ten-inch "regular" (Deci-Lon 10) and a five-inch "pocket" (Deci-Lon 5) variant. Another prized American model is the eight-inch Scientific Instruments circular rule. Of European rules, Faber-Castell's high-end models are the most popular among collectors. Although there is a large supply of slide rules circulating on the market, specimens in good condition tend to be expensive. Many rules found for sale on online auction sites are damaged or have missing parts, and the seller may not know enough to supply the relevant information. Replacement parts are scarce, expensive, and generally available only for separate purchase on individual collectors' web sites. The Keuffel and Esser rules from the period up to about 1950 are particularly problematic, because the end-pieces on the cursors, made of celluloid, tend to chemically break down over time. There are still a handful of sources for brand new slide rules. The Concise Company of Tokyo, which began as a manufacturer of circular slide rules in July 1954, continues to make and sell them today. In September 2009, on-line retailer ThinkGeek introduced its own brand of straight slide rules, described as "faithful replica[s]" that are "individually hand tooled". These are no longer available in 2012. In addition, Faber-Castell has a number of slide rules still in inventory, available for international purchase through their web store. Proportion wheels are still used in graphic design. Various slide rule simulator apps are available for Android and iOS-based smart phones and tablets. |Wikimedia Commons has media related to Slide rule.| - Lester V. Berrey and Melvin van den Bark (1953). American Thesaurus of Slang: A Complete Reference Book of Colloquial Speech. Crowell. - Roger R. Flynn (June 2002). Computer sciences 1. Macmillan. p. 175. ISBN 978-0-02-865567-3. Retrieved 30 March 2013. The slide rule is an example of a mechanical analog computer... - Swedin, Eric G.; Ferro, David L. (24 October 2007). Computers: The Life Story of a Technology. JHU Press. p. 26. ISBN 978-0-8018-8774-1. Retrieved 30 March 2013. Other analog mechanical computers included slide rules, the differential analyzer built by Vannevar E. Bush (1890–1974) at the ... - Peter Grego (2009). Astronomical cybersketching. Springer. p. 12. ISBN 978-0-387-85351-2. Retrieved 30 March 2013. It is astonishing to think that much of the routine mathematical work that put people into orbit around Earth and landed astronauts on the Moon in the 1960s was performed using an unassuming little mechanical analog computer – the 'humble' slide rule. - Ernst Bleuler; Robert Ozias Haxby (21 September 2011). Electronic Methods. Academic Press. p. 638. ISBN 978-0-08-085975-0. Retrieved 30 March 2013. For example, slide rules are mechanical analog computers, - Harry Henderson (1 January 2009). Encyclopedia of Computer Science and Technology, Revised Edition. Infobase Publishing. p. 13. ISBN 978-1-4381-1003-5. Retrieved 30 March 2013. Another analog computer, the slide rule, became the constant companion of scientists, engineers, and students until it was replaced ... logarithmic proportions, allowing for quick multiplication, division, the extraction of square roots, and sometimes the calculation of trigonometric functions. - Behrens, Lawrence; Rosen, Leonard J. (1982). Writing and reading across the curriculum. Little, Brown. p. 273. Then, just a decade ago, the invention of the pocket calculator made the slide rule obsolete almost overnight... - Maor, Eli (2009). e: The Story of a Number. Princeton University Press. p. 16. ISBN 978-0-691-14134-3. Then in the early 1970s the first electronic hand-held calculators appeared on the market, and within ten years the slide rule was obsolete. - Castleden, Rodney (2007). Inventions that Changed the World. Futura. p. 157. ISBN 978-0-7088-0786-6. With the invention of the calculator the slide rule became instantly obsolete. - Denning, Peter J.; Metcalfe, Robert M. (1998). Beyond calculation: the next fifty years of computing. Springer. p. xiv. ISBN 978-0-387-98588-6. The first hand calculator appeared in 1972 and made the slide rule obsolete overnight. - "instruction manual". sphere.bc.ca. pp. 7–8. Retrieved March 14, 2007. - "AntiQuark: Slide Rule Tricks". antiquark.com. - "Slide Rules". Tbullock.com. 2009-12-08. Retrieved 2010-02-20. - At least one circular rule, a 1931 Gilson model, sacrificed some of the scales usually found in slide rules in order to obtain additional resolution in multiplication and division. It functioned through the use of a spiral C scale, which was claimed to be 50 feet and readable to five significant figures. See http://www.sphere.bc.ca/test/gilson/gilson-manual2.jpg. A photo can be seen at http://www.hpmuseum.org/srcirc.htm. An instruction manual for the unit marketed by Dietzgen can be found at http://www.sliderulemuseum.com/SR_Library_General.htm. All retrieved March 14, 2007. - "336 (Teknisk Tidskrift / 1933. Allmänna avdelningen)". Runeberg.org. Retrieved 2010-02-20. - "Cameron's Nautical Slide Rule", The Practical Mechanic and Engineer's Magazine, April 1845, p187 and Plate XX-B - Kells, Lyman M.; Kern, Willis F.; Bland, James R. (1943). The Log-Log Duplex Decitrig Slide Rule No. 4081: A Manual. Keuffel & Esser. p. 92. Archived from the original on 14 February 2009. - The Polyphase Duplex Slide Rule, A Self-Teaching Manual, Breckenridge, 1922, p. 20. - "Lot 25368 Buzz Aldrin's Apollo 11 Slide Rule - Flown to the Moon. ... 2007 September Grand Format Air & Space Auction #669". Heritage Auctions. Retrieved 3 September 2013. - "Slide Rule, 5-inch, Pickett N600-ES, Apollo 13". Smithsonian National Air and Space Museum. Retrieved 3 September 2013. - Charles Overton Harris, Slide rule simplified, American Technical Society, 1961, p. 5. - "Prosthaphaeretic Slide Rule: A Mechanical Multiplication Device Based On Trigonometric Identities, The | Mathematics And Computer Education | Find Articles At Bnet". Findarticles.com. 2009-06-02. Retrieved 2010-02-20. - "Fisher sizing rules". natgasedu.com. Archived from the original on 6 January 2010. Retrieved 2009-10-06. - "The Wang LOCI-2". oldcalculatormuseum.com. - Wang Laboratories (December 1966). "Now you can determine Copolymer Composition in a few minutes at your desk". American Chemical Society 38 (13): 62A–63A. doi:10.1021/ac50155a005. Retrieved 2010-10-29. - "The HP 9100 Project". hp9825.com. - Volder, J. E. (2000). "The Birth of CORDIC". dx.doi.org (J. VLSI Signal Processing) 25 (2): 101. doi:10.1023/A:1008110704586. - Stoll, Cliff. "When Slide Rules Ruled," Scientific American, May 2006, pp. 80–87. "The difficulty of learning to use slide rules discouraged their use among the hoi polloi. Yes, the occasional grocery store manager figured discounts on a slipstick, and this author once caught his high school English teacher calculating stats for trifecta horse-race winners on a slide rule during study hall. But slide rules never made it into daily life because you could not do simple addition and subtraction with them, not to mention the difficulty of keeping track of the decimal point. Slide rules remained tools for techies." - Watson, George H. "Problem-based learning and the three C's of technology," The Power of Problem-Based Learning, Barbara Duch, Susan Groh, Deborah Allen, eds., Stylus Publishing, LLC, 2001. "Numerical computations in freshman physics and chemistry were excruciating; however, this did not seem to be the case for those students fortunate enough to already own a calculator. I vividly recall that at the end of 1974, the students who were still using slide rules were given an additional 15 minutes on the final examination to compensate for the computational advantage afforded by the calculator, hardly adequate compensation in the opinions of the remaining slide rule practitioners." - Stoll, Cliff. "When Slide Rules Ruled," Scientific American, May 2006, pp. 80–87. "With computation moving literally at a hand's pace and the lack of precision a given, mathematicians worked to simplify complex problems. Because linear equations were friendlier to slide rules than more complex functions were, scientists struggled to linearize mathematical relations, often sweeping high-order or less significant terms under the computational carpet. So a car designer might calculate gas consumption by looking mainly at an engine's power, while ignoring how air friction varies with speed. Engineers developed shortcuts and rules of thumb. At their best, these measures led to time savings, insight and understanding. On the downside, these approximations could hide mistakes and lead to gross errors." - Stoll, Cliff. "When Slide Rules Ruled", Scientific American, May 2006, pp. 80–87. "One effect was that users felt close to the numbers, aware of rounding-off errors and systematic inaccuracies, unlike users of today's computer-design programs. Chat with an engineer from the 1950s, and you will most likely hear a lament for the days when calculation went hand-in-hand with deeper comprehension. Instead of plugging numbers into a computer program, an engineer would understand the fine points of loads and stresses, voltages and currents, angles and distances. Numeric answers, crafted by hand, meant problem solving through knowledge and analysis rather than sheer number crunching." - "Citizen Watch Company – Citizen Eco-Drive / US, Canada, UK, IrelandCitizen Watch". citizenwatch.com. - "Greg's Slide Rules - Links to Slide Rule Collectors". Sliderule.ozmanor.com. 2004-07-29. Retrieved 2010-02-20. - "About CONCISE". Concise.co.jp. Archived from the original on 2012-03-12. Retrieved 2010-02-20. - "Slide Rule". ThinkGeek. Archived from the original on 2010-03-27. Retrieved 2015-04-08. - "Slide Rule". ThinkGeek. Archived from the original on 2012-10-12. Retrieved 2015-04-08. - "Rechenschieber". Faber-Castell. Retrieved 2012-01-17. - General information, history - International Slide Rule Museum - The history, theory and use of the engineering slide rule — By Dr James B. Calvert, University of Denver - United Kingdom Slide Rule Circle Home Page - Oughtred Society Slide Rule Home Page — Dedicated to the preservation and history of slide rules - Rod Lovett's Slide Rules - Comprehensive Aristo site with many search facilities - "Slide rule". New International Encyclopedia. 1905. - "Slide-rule". Encyclopedia Americana. 1920. - Reglas de Cálculo — A very big Faber Castell collection - Collection of slide rules — French Slide Rules (Graphoplex, Tavernier-Gravet and others) - Eric's Slide Rule Site — History and use

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Scientists studying the oceans depend on data from rivers to estimate how much fresh water and natural elements the continents are dumping into the oceans. But a new study in the Aug. 24 issue of Science finds that water quietly trickling along underground may double the amount of debris making its way into the seas. This study changes the equation for everything from global climate to understanding the ocean's basic chemistry. Since the late 1990s, Asish Basu, professor of earth and environmental sciences at the University of Rochester, has been sampling water and sediments from two of the world's largest rivers, the Ganges and the Brahmaputra of the Indian subcontinent, to understand a period in Earth's history called the Great Cool-Down. Forty million years ago, the global climate changed from the steamy world of the dinosaurs to the cooler world of today, largely because the amount of carbon dioxide, a greenhouse gas in the atmosphere, dropped significantly. Scientists have speculated that the cause of this cooling and the decline in atmospheric carbon dioxide was the result of the rise of the Himalayan mountains as the Indian and Asian continental plates pushed into one another. They believe the erosion of the new mountains increased the rate of removal of carbon dioxide from the atmosphere since the process of weathering silicate rocks such as those in the Himalayas absorbs carbon dioxide. This erosion may have depleted the atmosphere of a potent greenhouse gas and triggered the Great Cool-Down. Coinciding with the cooling period and Himalayan uplift 40 million years ago was a consistent change in the ratio of two isotopes of the element strontium in the oceans' water -- a change that continues to this day. Since strontium often comes from eroding silicates, it seemed obvious to scientists that the Ganges and Brahmaputra rivers were simply eroding the Himalayas into the ocean, but when they measured the amount of strontium in those rivers, they found it was far too low to account for the mysterious ratio change in the oceans, and thus too low to account for triggering the cool-down. To determine if enough silicate had eroded to spark the climate change, Basu and his colleagues analyzed both ground water and river water samples from the Bengal delta where the Ganges and Brahmaputra rivers empty. They found the missing strontium and confirmed the culprit that nudged down the thermostat. "Deep underground in the Bengal Basin, strontium concentration levels in the ground water are approximately 10 times higher than in the Ganges and Brahmaputra river waters," Basu explains. Knowing the speed the water is moving underground, Basu and his team calculated how much strontium could be leached out of the Bengal Basin and into the Indian Ocean. They calculated that about 1.4 times more strontium flows into the ocean through the groundwater than through the rivers above-easily enough to account for the 40 million-year rise. This study has other impacts in understanding ocean chemistry. "This means that we have to re-evaluate the residence times, the time a particular element remains in the ocean water before settling out, of various chemical elements and species," says Basu. "Most current studies on the ocean's chemistry are based on the supposition that the global rivers are the only carriers responsible for bringing in dissolved materials to the oceans. Our study changes that perception permanently." In addition, since the oceans are the biggest factor driving global weather, doubling the influx of fresh water will demand that global climate models must be restructured as well. Fresh water is lighter than salt water and so tends to float to the surface in the sea. This difference in density could move volumes of warm and cold water in ways that scientists gauging only the water's temperature would not normally predict. Working with Basu on the project were Stein Jacobsen of Harvard University, Robert Poreda and Carolyn Dowling of the University of Rochester, and Pradeep Agarwal of the International Atomic Energy Agency in Vienna, Austria. The research was partially supported by grants from the National Science Foundation. University of Rochester Subscribe To SpaceDaily Express Tropical Glaciers Formed While Earth Was Giant Snowball Boston - May 29, 2001 Glacial deposits that formed on tropical land areas during snowball Earth episodes around 600 million years ago, lead to questions about how the glaciers that left the deposits were created. Now, Penn State geoscientists believe that these glaciers could only have formed after the Earth's oceans were entirely covered by thick sea ice. New Research Documents Extremely High Atmospheric Carbon 14 During Last Ice Age by Lori Stiles Tucson - May 14, 2001 A team of American and British scientists report that radiocarbon levels in Earth's atmosphere during the last Ice Age were more than twice as high as today, higher even than the nuclear weapons tests of nearly half a century ago. They also reported in the May 11 issue of the journal Science of having extended the record for atmospheric radiocarbon more than 45,000 years. Climate Wobble Linked To Rare Anomaly In Earth's Orbit Santa Cruz - April 12, 2001 About 23 million years ago, a huge ice sheet spread over Antarctica, temporarily reversing a general trend of global warming and decreasing ice volume. Now a team of researchers has discovered that this climatic blip at the boundary between the Oligocene and Miocene epochs corresponded with a rare combination of events in the pattern of Earth's orbit around the Sun. |The content herein, unless otherwise known to be public domain, are Copyright 1995-2016 - Space Media Network. All websites are published in Australia and are solely subject to Australian law and governed by Fair Use principals for news reporting and research purposes. AFP, UPI and IANS news wire stories are copyright Agence France-Presse, United Press International and Indo-Asia News Service. ESA news reports are copyright European Space Agency. All NASA sourced material is public domain. Additional copyrights may apply in whole or part to other bona fide parties. Advertising does not imply endorsement, agreement or approval of any opinions, statements or information provided by Space Media Network on any Web page published or hosted by Space Media Network. Privacy Statement All images and articles appearing on Space Media Network have been edited or digitally altered in some way. Any requests to remove copyright material will be acted upon in a timely and appropriate manner. Any attempt to extort money from Space Media Network will be ignored and reported to Australian Law Enforcement Agencies as a potential case of financial fraud involving the use of a telephonic carriage device or postal service.|

http://www.spacedaily.com/news/iceage-01d.html

Since the observations made by English naturalist Charles Darwin on the Galapagos Islands, researchers have been interested in how physical barriers, such as isolation on a particular island, can lead to the formation of new species through the process of natural selection. Natural selection is a process whereby heritable traits that enhance survival become more common in successive generations, while unfavorable heritable traits become less common. Over time, animals and plants that have morphologies or other attributes that enhance their suitability to a particular environment become more common and more adapted to that specific environment. Researchers today are intimately familiar with how physical barriers and reproduction isolation can lead to the formation of new species on land, especially among plants and animals with short generation times such as insects and annual plants. Michael E. Hellberg, associate professor in the Department of Biological Sciences at LSU, however, is interested in a more obscure form of speciation: the speciation of animals in the ocean. "Marine plants and animals can drift around in the ocean extremely long distances," Hellberg said. "So how do they specialize?" In a recent publication in the Proceedings of the National Academy of Sciences, or PNAS, Hellberg and his graduate student Carlos Prada investigate how corals specialize to particular environments in the ocean. Corals, animals that form coral reefs and some of the most diverse ecosystems in the world, start their lifecycle with a free floating larval stage. Coral larvae can disperse vast distances in open water. Different coral species share similar geographical locations, with different species often existing only yards apart. As Prada and Hellberg propose in their recent publication, the large dispersal potential of coral larvae in open water and the proximity of different species on the ocean floor creates a mystery for researchers who study speciation. Hellberg and Prada ask, "How can new marine species emerge without obvious geographic isolation?" When it comes to corals within the relatively small confines of the Caribbean, which spans approximately 3 million square kilometers, the key to the puzzle appears to be habitat depth in the ocean. In others words, natural selection has led to the formation of different coral species according to how deep in the ocean these different corals grow. Prada and Hellberg study candelabrum corals of the genus Eunicea, generally known as "sea fans," for which sister species have been shown to be segregated by ocean depth. One sister species survives better in shallow waters, while the other is better adapted to deep waters. These corals, like other corals, are very slow-growing animals. In fact, sea fan corals don't reach reproduction age until they are 15-30 years old, and can continue reproducing until they are 60 or more years old. So while candelabrum coral larvae can disperse large distances from their parents, landing and beginning to grow in either shallow or deep water habitats, small differences in survival rates at different depths between the two species and long generation times can combine to produce segregation. "When these coral larvae first settle out after dispersal, they are all mixed up," Hellberg said. "But long larvae-to-reproduction times can compound small differences in survival at different depths. By the time these corals get to reproduction age, a lot has changed." The shallow water sea fan coral even has a different morphology than its deep water sister. The shallow water coral fans out into a wide network of branches, while the deep water coral grows tall and spindly. According to Hellberg, these differences in morphology might well be genetic, with the different corals having different protein structures and levels of expression that are better adapted to their specific water depth environment. Hellberg hopes in future research to investigate the genetic basis of these different morphologies. In other interesting results, Prada explained how transplanting the shallow coral species to deep water environments, and vice versa, can cause the coral to take on a morphology more like that of its sister species. "Their morphologies are not super fixed," Prada said. "But they can't change all the way to a different morphology." Prada observed that while shallow water sea fans can become taller and more spindly when transplanted in deep water environments, they don't seem to be able to make a complete transition to the morphology of the deep water sea fan. This suggests that the two corals, while they likely had a common ancestor, have adapted genetically and biochemically to their respective water depths. Prada did ocean dives in the Bahamas, Panama, Puerto Rico and Curaao to sample candelabrum coral colonies. Back in the lab, he performed tests on the coral samples' genes to determine how shallow and deep corals become genetically different. "Normally, organisms are differentiated by geography," Prada said. "But these corals are differentiated by depth." Prada and Hellberg's research provides new insights into how new species form in the ocean, a topic of relatively limited research as opposed to speciation of terrestrial organisms. |Contact: Ashley Berthelot| Louisiana State University

http://www.bio-medicine.org/biology-news-1/LSU-professor-discovers-how-new-corals-species-form-in-the-ocean-28752-1/

Birthplace: Saukenuk, Ill. In the late 18th century, the Indians of the upper Mississippi Valley witnessed the replacement of the relatively sympathetic French, Spanish, and British with the aggressive Americans pushing westward. The Sauk warrior, Black Hawk, resisted the pioneers and fought to have Indians retain their lands and traditions. Black Hawk was especially incensed by an 1804 treaty between the Sauk and Fox tribes and the United States that ceded all tribal lands east of the Mississippi. The treaty had never been ratified by the tribe, and Black Hawk repeatedly condemned it as spurious. Black Hawk and his band of warriors fought on the side of the British during the War of 1812, hoping to halt the American westward expansion. While Black Hawk was fighting the United States, the young Keokuk, who was friendlier towards the Americans, became leader of the Sauk and Fox. By 1814, Black Hawk and his forces defeated the Americans, who were under the command of Gen. Zachary Taylor. Treaties were signed in 1816 and an uneasy peace reigned until 1832 when the Black Hawk War broke out. The Sauk, Fox, and other tribes refused to move westward to accommodate the increasingly large population of American pioneers. President Andrew Jackson sent troops and a massacre at Black Axe River ended the war and resulted in Black Hawk's imprisonment for several months. After Black Hawk's defeat, Keokuk, who had maintained good relations with the U.S. government, was granted a tract of land in Iowa for his people. Black Hawk joined what remained of his tribe in Iowa and died there in 1838.Died: 1838

http://www.infoplease.com/ipa/A0909619.html

X-ray crystallography, the study of crystal structures through X-ray diffraction techniques. When an X-ray beam bombards a crystalline lattice in a given orientation, the beam is scattered in a definite manner characterized by the atomic structure of the lattice. This phenomenon, known as X-ray diffraction, occurs when the wavelength of X-rays and the interatomic distances in the lattice have the same order of magnitude. In 1912, the German scientist Max von Laue predicted that crystals exhibit diffraction qualities. Concurrently, W. Friedrich and P. Knipping created the first photographic diffraction patterns. A year later Lawrence Bragg successfully analyzed the crystalline structures of potassium chloride and sodium chloride using X-ray crystallography, and developed a rudimentary treatment for X-ray/crystal interaction (Bragg's Law). Bragg's research provided a method to determine a number of simple crystal structures for the next 50 years. In the 1960s, the capabilities of X-ray crystallography were greatly improved by the incorporation of computer technology. Modern X-ray crystallography provides the most powerful and accurate method for determining single-crystal structures. Structures containing 100–200 atoms now can be analyzed on the order of 1–2 days, whereas before the 1960s a 20-atom structure required 1–2 years for analysis. Through X-ray crystallography the chemical structure of thousands of organic, inorganic, organometallic, and biological compounds are determined every year. See M. Buerger, X-Ray Crystallography (1980). The Columbia Electronic Encyclopedia, 6th ed. Copyright © 2012, Columbia University Press. All rights reserved.

http://www.factmonster.com/encyclopedia/science/x-ray-crystallography.html

1789 to Present Article I, section 7 of the Constitution grants the President the authority to veto legislation passed by Congress. This authority is one of the most significant tools the President can employ to prevent the passage of legislation. Even the threat of a veto can bring about changes in the content of legislation long before the bill is ever presented to the President. The Constitution provides the President 10 days (excluding Sundays) to act on legislation or the legislation automatically becomes law. There are two types of vetoes: the “regular veto” and the “pocket veto.” The regular veto is a qualified negative veto. The President returns the unsigned legislation to the originating house of Congress within a 10 day period usually with a memorandum of disapproval or a “veto message.” Congress can override the President’s decision if it musters the necessary two–thirds vote of each house. President George Washington issued the first regular veto on April 5, 1792. The first successful congressional override occurred on March 3, 1845, when Congress overrode President John Tyler’s veto of S. 66. The pocket veto is an absolute veto that cannot be overridden. The veto becomes effective when the President fails to sign a bill after Congress has adjourned and is unable to override the veto. The authority of the pocket veto is derived from the Constitution’s Article I, section 7, “the Congress by their adjournment prevent its return, in which case, it shall not be law.” Over time, Congress and the President have clashed over the use of the pocket veto, debating the term “adjournment.” The President has attempted to use the pocket veto during intra- and inter- session adjournments and Congress has denied this use of the veto. The Legislative Branch, backed by modern court rulings, asserts that the Executive Branch may only pocket veto legislation when Congress has adjourned sine die from a session. President James Madison was the first President to use the pocket veto in 1812. |Congresses||President||Regular Vetoes||Pocket Vetoes||Total Vetoes||Vetoes Overriden| |19th–20th||John Quincy Adams||.....||.....||.....||.....| |25th–26th||Martin Van Buren||.....||1||1||.....| |27th||William Henry Harrison||.....||.....||.....||.....| |29th–30th||James K. Polk||2||1||3||.....| |41st–44th||Ulysses S. Grant||45||48||93||4| |45th–46th||Rutherford B. Hayes||12||1||13||1| |47th||James A. Garfield||.....||.....||.....||.....| |47th–48th||Chester A. Arthur||4||8||12||1| |61st–62nd||William H. Taft||30||9||39||1| |67th||Warren G. Harding||5||1||6||.....| |71st–72nd||Herbert C. Hoover||21||16||37||3| |73rd–79th||Franklin D. Roosevelt||372||263||635||9| |79th–82nd||Harry S. Truman||180||70||250||12| |83rd–86th||Dwight D. Eisenhower||73||108||181||2| |87th–88th||John F. Kennedy||12||9||21||.....| |88th–90st||Lyndon B. Johnson||16||14||30||.....| |91st–93rd||Richard M. Nixon||26||17||43||7| |93rd–94th||Gerald R. Ford||48||18||66||12| |95th–96th||James Earl Carter||13||18||31||2| |101st–102nd||George H. W. Bush1||29||15||44||1| |103rd–106th||William J. Clinton2||36||1||37||2| |107th–110th||George W. Bush3||12||.....||12||4| |111th–114th||Barack H. Obama4||8||.....||8||.....| 1President George H. W. Bush withheld his signature from two measures during intrasession recess periods (H.J. Res. 390, 101st Congress, 1st sess. and S. 1176, 102nd Congress, 1st sess.). See, “Permission to Insert in the Record Correspondence of the Speaker and the Minority Leader to the President Regarding Veto of House Joint Resolution 390, Authorizing Hand Enrollment of H.R. 1278, Financial Institutions Reform, Recovery and Enforcement Act of 1989, Along With Response From the Attorney General (House of Representatives - January 23, 1990),” Congressional Record, 101st Cong., 2nd sess., (January 23, 1990): H3. See, “Morris K. Udall Scholarship and Excellence in National Environmental and Native American Public Policy Act of 1992 (House of Representatives - March 03, 1992),” Congressional Record, 102nd Cong., 2nd sess., (March 3, 1992): H885-H889. The President withheld his signature from another measure during an intrasession recess period (H.R. 2699, 102nd Congress, 1st sess.) and from a measure during an intersession recess period (H.R. 2712, 101st Congress, 1st sess.) but returned both measures to the House, which proceeded to reconsider them. The measures are not included as pocket vetoes in this table. 2President William J. Clinton withheld his signature from two measures during intrasession recess periods (H.R. 4810, 106th Congress, 2nd sess., and H.R. 8, 106th Congress, 2nd sess.) but returned the bills to the House, which proceeded to reconsider them. See, “Pocket-Veto Power -- Hon. J. Dennis Hastert – (Extensions of Remarks - September 19, 2000),” Congressional Record, 106th Cong., 2nd sess., (September 19, 2000): E1523. The bills are not included as pocket vetoes in this table. 3President George W. Bush withheld his signature from a measure during an intersession recess period (H.R. 1585, 110th Congress, 1st Sess.) but returned the bill to the House, which proceeded to reconsider it. See, “Pocket-Veto Power – (Extensions of Remarks – October 2, 2008),” Congressional Record, 110th Cong., 1st Sess., (October 2, 2008): E2197. The bill is not included as a pocket veto in this table. 4President Barack H. Obama withheld his signature from a measure during an intersession recess period (H.J. Res 64, 111th Congress, 1st sess.) and from a measure during an intrasession recess period (H.R. 3808, 111th Congress, 2nd sess.) but returned both measures to the House, which proceeded to reconsider them. “Pocket-Veto Power – (Extensions of Remarks – May 26, 2010),” Congressional Record, 111th Cong., 1st sess., (May 26, 2010): E941. The measures are not included as pocket vetoes in this table.

http://history.house.gov/Institution/Presidential-Vetoes/Presidential-Vetoes/

The Elves and the Shoemaker Teacher Resources Find The Elves and the Shoemaker educational ideas and activities Showing 1 - 20 of 22 resources Folk and Fairy Tale Readers: The Elves and the Shoemaker Engage young readers in a unit on fairy tales with their very own copy of the classic story "The Elves and the Shoemaker." Including dialogue, three-syllable words, and up to five lines of text on a page, this printable book is best... Pre-K - 3rd English Language Arts CCSS: Adaptable Guided Reading Lesson: The Elves and the Shoemaker Student participate in a guided reading lesson for the Level H book The Elves and the Shoemaker. In this guided reading lesson, 2nd graders learn vocabulary associated with the book and focus on inferences, connections, and other... 2nd English Language Arts Library Skills and Literature New ReviewThe library is such a valuable resource for kids of all ages. Help elementary readers learn all about parts of the library, text features for both fiction and nonfiction text, and different ways to find books that they want to read. K - 5th English Language Arts CCSS: Adaptable Discovering Shoes, Step by Step Students discuss and write about their interpretations of the art. They compare and contrast the numerous types of shoes and how they are used in a certain time and place. Students memorize the Shel Silverstein's poem, "Ickle Me, Pickle... 4th - 6th English Language Arts

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An Introduction to Vegetative Reproduction is a form of asexual reproduction in plants. It does not involve flowers, pollination and seed production. Instead, a new plant grows from a vegetative part, usually a stem, of the parent plant. However, plants which reproduce asexually almost always reproduce sexually as well, bearing flowers, fruits and seeds. Vegetative reproduction from a stem usually involves the buds. Instead of producing a branch, the bud grows into a complete plant which eventually becomes self-supporting. Since no gametes are involved, the plants produced asexually have identical genomes and the offspring form what is known as a clone. In some cases of vegetative reproduction, the structures involved also become storage organs and swell with stored food, e.g. potatoes. The principal types of vegetative reproduction structures are bulbs, corms, rhizomes and runners. Bulbs consist of very short stems with closely packed leaves arranged in concentric circles round the stem. These leaves are swollen with stored food e.g. onion. A terminal bud will produce next year’s flowering shoot and the lateral (axillary) buds will produce new plants. Corms also have a short stem but in this case it is the stem itself which swells and stores food. The circular leaves form only papery scales. As with bulbs, the terminal bud grows into a flowering shoot and the lateral buds produce new plants. Rhizomes are stems which grow horizontally under the ground. In some cases the underground stems are swollen with food reserves e.g. iris. The terminal bud turns upwards to produce the flowering shoot and the lateral buds may grow out to form new rhizomes. Runners are also horizontal stems growing from the parent plant, but they grow above ground. When their terminal buds touch the ground they take root and produce new plants. Advantages of vegetative reproduction Since food stores are available throughout the year and the parent plant with its root system can absorb water from quite a wide area, two of the hazards of seed germination are reduced. Buds are produced in an environment where the parent is able to flourish, but many seeds dispersed from plants never reach a suitable situation for effective germination. Vegetative reproduction does not usually result in rapid and widespread distribution of offspring in the same way as seed dispersal, but tends to produce a dense clump of plants with little room for competitors between them. Such groups of plants are very persistent and, because of their buds and underground food stores, can still grow after their foliage has been destroyed by insects, fire, or cultivation. Those of them regarded as weeds are difficult to eradicate, since even a small piece of rhizome bearing a bud can give rise to a new colony (clone). Bulbs - Snowdrop In the snowdrop and daffodil, the bulb is formed by the leaf bases which completely encircle the short, conical stem. The part of the leaf above ground makes food by photosynthesis and sends it to the leaf bases which swell as they store the food. In the following year the stored food is used for the early growth of the bulb. Life cycle of Daffodil In the spring, adventitious roots grow out of the stem, and the leaves begins to grow above ground, making use of the stored food in the fleshy leaf bases which consequently shrivel. During late spring some of the food made in the leaves in the daffodil is sent to the leaf bases which swell and form a new bulb inside the old one. Life cycle of Tulip and Onion In these bulbs, the food is not sent to the leaf bases but to the lateral buds. As these buds enlarge they form two or more ‘daughter’ bulbs inside the old bulb. The leaves of the old bulb shrivel and dry out forming the dry scales which surround the daughter bulbs. In both cases, when the daughter bulbs grow, they form a clump, together with the parent bulb. Corms - Crocus Plants with bulbs store food in special leaves or leaf bases. Plants with corms store food in the stem, which is very short and swollen. When the foliage has died off, the leaf bases, where they encircle the short stem, form protective scaly coverings. A familiar corms is that of the crocus, and the wild arum corm is illustrated on p.1. Since the corm is a stem, it has lateral buds which can grow into new plants. The stem remains below ground all its life, only the leaves and flower stalk coming above ground. Life Cycle of Corm In Spring, the food stored in the corm enables the terminal bud to grow rapidly and produce leaves and flowers above ground. Later in the year, food made by the leaves is sent back, not to the old corm, but to the base of the stem immediately above it. This region swells and forms a new corm on top of the old, now shrivelled, corm. Some of the lateral buds on the old corm have also grown and produced new plants with corms. The formation of one corm on top of another tends to bring the successive corms nearer and nearer to the soil surface. Adventitious roots develop from the base of the new corm. Once these have grown firmly into the soil, a region near their junction with the stem contracts and pulls the new corm down, keeping it at a constant level in the soil. Wrinkles can be seen on these contractile roots where shrinkage has taken place. Bulbs also have contractile roots which counteract the tendency in successive generations to grow out of the soil. In plants with rhizomes, the stem remains below ground but continues to grow horizontally. The old part of the stem does not die away as in bulbs and corms, but lasts for several years. In the iris, the terminal bud turns up and produces leaves and flowers above ground. The old leaf bases form circular scales round the rhizome, which is swollen with food reserves. Lateral buds grow into new rhizomes. Life Cycle of Rhizome The annual cycle of a rhizome is similar to that of a corm. In late spring/early summer, food from the leaves passes back to the rhizome, and a lateral bud uses it, grows horizontally underground, and so continues the rhizome. Other lateral buds produce new rhizomes which branch from the parent stem. The terminal buds of these branches curve upwards and produce new leafy shoots and flowers. Contractile, adventitious roots grow from the nodes of the underground stem and keep it at a constant depth. Plants such as the strawberry have a very short stem, called a rootstock, with thin scale leaves,. Foliage leaves and flowers grow from the buds in the axils of the scale leaves. Some of the lower buds produce shoots which grow horizontally over the surface of the ground and bear scale leaves and buds. The terminal buds of these runners turn up and produce daughter plants some distance away from the parent, the new plants developing adventitious roots. Later, the runner shrivels away. The runner does not store food but conducts it from the parent plant to the daughters, until they are well developed. Stem tubers - potato In the potato plant, lateral buds at the base of the stem produce shoots which grow laterally at first and then down into the ground. These are comparable to rhizomes, as they are underground stems with tiny scale leaves and lateral buds. They do not swell evenly along their length with stored food. Annual cycle - potato Food made in the leaves passes to the ends of these rhizomes, which swell and form the tubers we call potatoes. Since the potato tuber is a stem, it has leaves and lateral buds; these are the familiar ‘eyes’. Each one of these can produce a new shoot in the following year, using the food stored in the tuber. The old tubers shrivel and rot away at the end of the season Blackberry stems form a rather different type of runner in which the main shoot forms the new individual. When the growing end of a shoot arches over and touches the ground, the terminal bud curves up, producing a new shoot which soon develops adventitious roots. A bud or shoot from one plant is inserted into a cleft or under the bark on the stem of a closely related variety. The rooted portion is called the stock; the bud or shoot being grafted is the scion. The stock is obtained by growing a plant from seed then cutting away the shoot. The scion is a branch or a bud cut from a cultivated variety with the required characteristics of flower colour, fruit quality, etc. Rose plants grown from seed would produce a wide variety of plants, only a few of which would retain all the desirable features of the parent plant. Most of them would be like wild roses. Similarly, most of the apple trees grown from seed would bear only small, sour ‘crab-apples'. By taking cuttings and making grafts, the inbred characteristics of the plant are preserved and you can guarantee that all the new individuals produced by this kind of artificial propagation will be the same. It is possible to produce new individuals from certain plants by putting the cut end of a shoot into water or moist earth. Roots grow from the base of the stem into the soil while the shoot continues to grow and produce leaves. In some cases the cut end of the stem may be treated with a rooting 'hormone' to promote root growth. Evaporation from the shoot is reduced by covering it with polythene or a glass jar. Carnations, geraniums and chrysanthemums are commonly propagated from cuttings. Once a cell has become part of a tissue it usually loses the ability to reproduce. However, the nucleus of any plant cell still holds all the 'instructions' (genes) for making a complete plant and in certain circumstances they can be brought back into action. In laboratory conditions single plant cells can be induced to divide and grow into complete plants. One technique is to take a small piece of plant tissue from a root or stem and treat it with enzymes to separate it into individual cells The cells are then provided with particular plant 'hormones’ which induce cell division and, eventually the formation of roots, stems and leaves. An alternative method is to start with a small piece of tissue and place it on a nutrient jelly (agar). Cells in the tissue start to divide and produce many cells forming a shapeless mass called a callus. If the callus is then provided with the appropriate ‘hormones’ it develops into a complete plant. |Search this site| |Search the web| © Copyright 2004 - 2016 D G Mackean & Ian Mackean. All rights reserved.

http://www.biology-resources.com/plants-vegetative-reproduction-01.html

A Recycling for Kids, theme need not be complicated or abstract. Instead take advantage of daily classroom routines when planning recycling type activities. Each day at snack or lunch time, set up a classification area on one table in your classroom. Make it one of your special helper’s jobs. Students recycle their garbage into three categories when they finish eating, recycling, compost, and landfill garbage. Kids are very enthusiastic to do this and quickly learn what can be recycled, what goes in the compost and what ends up in the landfill. Teach young children about recycling - a plastic ice cream bucket with a lid for the compost - a photocopy paper box lid for items that go in the garbage - a photocopy paper box lid for items that can be recycled. The size of the photocopy paper box lids allow the items to be spread out for easy counting. Keep it Simple To keep recycling for kids simple, have the special helper weigh the compost (with help) and then point to each piece of recycling and garbage in the lids as the whole class counts along. The same student records the observations on a chart similar to the sample on the above. The teacher records the date as the students suggest letters and assists the students when necessary. The children benefited from the extra practice counting the recycling and garbage items, and recording, or watching the numbers be recorded under the correct pictures on the chart. Each day use the chart as a teaching tool, reinforcing the children’s knowledge of numbers, letters and letter sounds. In many schools, older students pick up the compost after lunch for the school garden. The recycling items went into the class blue box and the landfill garbage went into the wastebasket. Next, have the students record their classroom activities.

http://www.kindergarten-lessons.com/recycling_for_kids/

Science Fair1. Present Scientific Method slides containing information on a display board Exhibit the experiment with the display board2. Parents are encouraged to assist children but are asked to maintain the role of assistant only.3. If parents type the display board information, they should have children dictate the information as it is being typed.4. Construction paper or colored paper should be used behind all display board items to give a neat, organized, and creative effect.5. Students are expected to use these Scientific Method slides in the given order for continuity on the display boards as it is helpful to judges. Students may change the theme of the slides to make their presentations more unique or creative.Remember that the Science Fair should be a fun method of learning and is intended to help in preparing students for the STAAR test. Science Fair Project Type/write your project title here. Put your assigned number without your name for the display board QUESTIONType/write your question here. (This is the question that your experiment answers.) Materials• Type/write a detailed list of the items you need to complete your experiments.• Be specific about the amounts used. Procedure• List all of the steps used in completing your experiment.• Remember to number your steps.• If possible, add photos of your experiments. Research Summarize your research here in five or more bullet points: Do not give more than ten points, however.• 1st bullet point• 2nd bullet point• 3rd bullet point• 4th bullet point• 5th bullet point HypothesisWhat do you think the result of your experiment willshow?It is okay if you are incorrect. Variables• Controlled variables: These are the things that are kept the same throughout your experiments.• Independent variable: The ONE variable that you purposely change and test.• For example, a plant experiment would use the type of plant as the control variable while the type of soil may be the independent variable that changes. Data/Observations• It is easier to understand the data (information) if it is put into a table or graph.• Draw a graph or make one in a software package such as PowerPoint or Microsoft Excel.• Make sure all data (information) is clearly labeled. ConclusionType/write a brief summary here of your discoverybased on the results of the experiment. Indicatewhether or not the data (information) supports thehypothesis and explain why or why not. Works Cited• Be sure to list books, magazines, encyclopedias, and Internet sources that you used.• Put them in alphabetical order. Presentation• Students will be expected to stand beside their projects on Parent Night and during Science class period as visitors tour the exhibits.• The students’ behavior while presenting their information will be included in the grade for the Science Fair.• Each topic will be judged separately.• Topics Include: Earth Science, Life Science, and Physical Science• The students’ grades will be determined by Mrs. Hagood and will not be influenced by judges decisions.• Scoring will be based on how well students meet the criteria. Unacceptable Projects• Remember that the experiment must prove a hypothesis and cannot be a model of how something works or a report on another person’s research . For example, students cannot actually go into outer space to research how plants grow; therefore, a project of this type would be research. Also a volcano could be created to show chemical reactions with research on that area but should not be presented as an experiment on how a volcano works as vinegar and baking soda are not materials found in lava and cannot be utilized as an experiment on Earth processes where a variable is used to prove a hypothesis. Tips on Placing in the Fair• Select an experiment that few others have chosen.• The more technical the experiment, the better.• Colorful and creative displays are a plus.• Multiple charts and graphs are beneficial.• Research that goes beyond the basic requirement of five points but does not exceed ten points . Remember to Follow Safety Procedures• Students should wear safety glasses,• Do not taste anything unless it is approved by a responsible adult,• Wear protective clothing if needed,• Avoid using dangerous materials or fire/heat unless supervised by a responsible adult,• If animals are used, take precautions that prevent students from being bitten or harmed in any way,• Use good judgment and common sense at all times. Judging SheetJudges should score projects within a range of 1 to 5. Consider a 1 to be very low, a 2 is close to meetingthe expectation, 3 should reflect that the student did what was expected. A 4 will be above average with a 5 being a superior rating. If possible, it will be appreciated if scores are totaled. 1. Title Page includes Number & Tchr.___________________ 2. Question_________________ 3. Materials _________________________________________ 4. Procedure- detailed steps of experiment__________________ 5. Research (Must have 5 facts)__________________________ Project Number 6. Variables _________________________________________ ___________________ 1. Controlled variable- is the step/material that stays the same 2. Independent variable- is the step/material that changes. 7. Hypothesis –What student thinks will happen______________ Project Title 8. Data/Observations – In graph or table form________________ ______________________ 9. Conclusion- What they discovered__ _____________________ 10. Works Cited- Where information was found ________________ 11. Neat -Board has slides displayed with construction paper backing Judge Number and is neatly arranged__________________________________ ____________________ 12. Technical (Bonus Points)_______________________________ 13. Creative-(Bonus Points)_________________________________ Total ____________

http://www.slideshare.net/test3student/science-fair-training-slides

|Part of a series on the| |Spanish conquest of the Maya| The Spanish conquest of Yucatán was the campaign undertaken by the Spanish conquistadores against the Late Postclassic Maya states and polities in the Yucatán Peninsula, a vast limestone plain covering south-eastern Mexico, northern Guatemala, and all of Belize. The Spanish conquest of the Yucatán Peninsula was hindered by its politically fragmented state. The Spanish engaged in a strategy of concentrating native populations in newly founded colonial towns. Native resistance to the new nucleated settlements took the form of the flight into inaccessible regions such as the forest or joining neighbouring Maya groups that had not yet submitted to the Spanish. Among the Maya, ambush was a favoured tactic. Spanish weaponry included broadswords, rapiers, lances, pikes, halberds, crossbows, matchlocks and light artillery. Maya warriors fought with flint-tipped spears, bows and arrows and stones, and wore padded cotton armour to protect themselves. The Spanish introduced a number of Old World diseases previously unknown in the Americas, initiating devastating plagues that swept through the native populations. The first encounter with the Yucatán Maya occurred in 1502, when the fourth voyage of Christopher Columbus came across a large Maya trading canoe off Honduras. In 1517, Francisco Hernández de Córdoba made landfall on the tip of the peninsula. His expedition continued along the coast and suffered heavy losses in a pitched battle at Champotón, forcing a retreat to Cuba. Juan de Grijalva explored the coast in 1518, and heard tales of the wealthy Aztec Empire further west. As a result of these rumours, Hernán Cortés set sail with another fleet. From Cozumel he continued around the peninsula to Tabasco where he fought a battle at Potonchán; from there Cortés continued onward to conquer the Aztec Empire. In 1524, Cortés led a sizeable expedition to Honduras, cutting across southern Campeche, and through Petén in what is now northern Guatemala. In 1527 Francisco de Montejo set sail from Spain with a small fleet. He left garrisons on the east coast, and subjugated the northeast of the peninsula. Montejo then returned to the east to find his garrisons had almost been eliminated; he used a supply ship to explore southwards before looping back around the entire peninsula to central Mexico. Montejo pacified Tabasco with the aid of his son, also named Francisco de Montejo. In 1531 the Spanish moved their base of operations to Campeche, where they repulsed a significant Maya attack. After this battle, the Spanish founded a town at Chichen Itza in the north. Montejo carved up the province amongst his soldiers. In mid-1533 the local Maya rebelled and laid siege to the small Spanish garrison, which was forced to flee. Towards the end of 1534, or the beginning of 1535, the Spanish retreated from Campeche to Veracruz. In 1535, peaceful attempts by the Franciscan Order to incorporate Yucatán into the Spanish Empire failed after a renewed Spanish military presence at Champotón forced the friars out. Champotón was by now the last Spanish outpost in Yucatán, isolated among a hostile population. In 1541–42 the first permanent Spanish town councils in the entire peninsula were founded at Campeche and Mérida. When the powerful lord of Mani converted to the Roman Catholic religion, his submission to Spain and conversion to Christianity encouraged the lords of the western provinces to accept Spanish rule. In late 1546 an alliance of eastern provinces launched an unsuccessful uprising against the Spanish. The eastern Maya were defeated in a single battle, which marked the final conquest of the northern portion of the Yucatán Peninsula. The polities of Petén in the south remained independent and received many refugees fleeing from Spanish jurisdiction. In 1618 and in 1619 two unsuccessful Franciscan missions attempted the peaceful conversion of the still pagan Itza. In 1622 the Itza slaughtered two Spanish parties trying to reach their capital Nojpetén. These events ended all Spanish attempts to contact the Itza until 1695. Over the course of 1695 and 1696 a number of Spanish expeditions attempted to reach Nojpetén from the mutually independent Spanish colonies in Yucatán and Guatemala. In early 1695 the Spanish began to build a road from Campeche south towards Petén and activity intensified, sometimes with significant losses on the part of the Spanish. Martín de Urzúa y Arizmendi, governor of Yucatán, launched an assault upon Nojpetén in March 1697; the city fell after a brief battle. With the defeat of the Itza, the last independent and unconquered native kingdom in the Americas fell to the Spanish. - Yucatán before the conquest - Impact of Old World diseases - Weaponry, strategies and tactics - First encounters: 1502 and 1511 - Francisco Hernández de Córdoba, 1517 - Juan de Grijalva, 1518 - Hernán Cortés, 1519 - Hernán Cortés in the Maya lowlands, 1524–25 - Francisco de Montejo, 1527–28 - Francisco de Montejo and Alonso d'Avila, 1531–35 - Conquest and settlement in northern Yucatán, 1540–46 - Petén Basin, 1618–97 - Further reading The Yucatán Peninsula is bordered by the Caribbean Sea to the east and by the Gulf of Mexico to the north and west. It can be delimited by a line running from the Laguna de Términos on the Gulf coast through to the Gulf of Honduras on the Caribbean coast. It incorporates the modern Mexican states of Yucatán, Quintana Roo and Campeche, the eastern portion of the state of Tabasco, most of the Guatemalan department of Petén, and all of Belize. Most of the peninsula is formed by a vast plain with few hills or mountains and a generally low coastline. A 15-kilometre (9.3 mi) stretch of high, rocky coast runs south from the city of Campeche on the Gulf Coast. A number of bays are situated along the east coast of the peninsula, from north to south they are Ascensión Bay, Espíritu Santo Bay, Chetumal Bay and Amatique Bay. The north coast features a wide, sandy littoral zone. The extreme north of the peninsula, roughly corresponding to Yucatán State, has underlying bedrock consisting of flat Cenozoic limestone. To the south of this the limestone rises to form the low chain of Puuc Hills, with a steep initial scarp running 160 kilometres (99 mi) east from the Gulf coast near Champotón, terminating some 50 kilometres (31 mi) from the Caribbean coast near the border of Quintana Roo. The hills reach a maximum altitude of 170 metres (560 ft). The northwestern and northern portions of the Yucatán Peninsula experience lower rainfall than the rest of the peninsula; these regions feature highly porous limestone bedrock resulting in less surface water. This limestone geology results in most rainwater filtering directly through the bedrock to the phreatic zone, from whence it slowly flows to the coasts to form large submarine springs. Various freshwater springs rise along the coast to form watering holes. The filtering of rainwater through the limestone has caused the formation of extensive cave systems. These cave rooves are subject to collapse forming deep sinkholes; if the bottom of the cave is deeper than the groundwater level then a cenote is formed. In contrast, the northeastern portion of the peninsula is characterised by forested swamplands. The northern portion of the peninsula lacks rivers, except for the Champotón River – all other rivers are located in the south. The Sibun River flows from west to east from south central Quintana Roo to Lake Bacalar on the Caribbean Coast; the Río Hondo flows northwards from Belize to empty into the same lake. Bacalar Lake empties into Chetumal Bay. The Río Nuevo flows from Lamanai Lake in Belize northwards to Chetumal Bay. The Mopan River and the Macal River flow through Belize and join to form the Belize River, which empties into the Caribbean Sea. In the southwest of the peninsula, the San Pedro River, the Candelaría River and the Mamantel River, which all form a part of the Gulf of Mexico drainage. The Petén region consists of densely forested low-lying limestone plain featuring karstic topography. The area is crossed by low east–west oriented ridges of Cenozoic limestone and is characterised by a variety of forest and soil types; water sources include generally small rivers and low-lying seasonal swamps known as bajos. A chain of fourteen lakes runs across the central drainage basin of Petén; during the rainy season some of these lakes become interconnected. This drainage area measures approximately 100 kilometres (62 mi) east–west by 30 kilometres (19 mi) north–south. The largest lake is Lake Petén Itza, near the centre of the drainage basin; it measures 32 by 5 kilometres (19.9 by 3.1 mi). A broad savannah extends south of the central lakes. To the north of the lakes region bajos become more frequent, interspersed with forest. In the far north of Petén the Mirador Basin forms another interior drainage region. To the south the plain gradually rises towards the Guatemalan Highlands. The canopy height of the forest gradually decreases from Petén northwards, averaging from 25 to 35 metres (82 to 115 ft). This dense forest covers northern Petén and Belize, most of Quinatana Roo, southern Campeche and a portion of the south of Yucatán State. Further north, the vegetation turns to lower forest consisting of dense scrub. The climate becomes progressively drier towards the north of the peninsula. In the north, the annual mean temperature is 27 °C (81 °F) in Mérida. Average temperature in the peninsula varies from 24 °C (75 °F) in January to 29 °C (84 °F) in July. The lowest temperature on record is 6 °C (43 °F). For the peninsula as a whole, the mean annual precipitation is 1,100 millimetres (43 in). The rainy season lasts from June to September, while the dry season runs from October to May. During the dry season, rainfall averages 300 millimetres (12 in); in the wet season this increases to an average 800 to 900 millimetres (31 to 35 in). The prevailing winds are easterly and have created an east-west precipitation gradient with average rainfall in the east exceeding 1,400 millimetres (55 in) and the north and northwestern portions of the peninsula receiving a maximum of 800 millimetres (31 in). The southeastern portion of the peninsula has a tropical rainy climate with a short dry season in winter. Petén has a hot climate and receives the highest rainfall in all Mesoamerica. The climate is divided into wet and dry seasons, with the rainy season lasting from June to December, although these seasons are not clearly defined in the south; with rain occurring through most of the year. The climate of Petén varies from tropical in the south to semitropical in the north; temperature varies between 12 and 40 °C (54 and 104 °F), although it does not usually drop beneath 18 °C (64 °F). Mean temperature varies from 24.3 °C (75.7 °F) in the southeast to 26.9 °C (80.4 °F) in the northeast. Highest temperatures are reached from April to June, while January is the coldest month; all Petén experiences a hot dry period in late August. Annual precipitation is high, varying from a mean of 1,198 millimetres (47.2 in) in the northeast to 2,007 millimetres (79.0 in) in central Petén. Yucatán before the conquest The first large Maya cities developed in the Petén Basin in the far south of the Yucatán Peninsula as far back as the Middle Preclassic (c. 600–350 BC), and Petén formed the heartland of the ancient Maya civilization during the Classic period (c. AD 250–900). The 16th century Maya provinces of northern Yucatán are likely to have evolved out of polities of the Maya Classic period. From the mid-13th century AD through to the mid-15th century, the League of Mayapán united several of the northern provinces; for a time they shared a joint form of government. The great cities that dominated Petén had fallen into ruin by the beginning of the 10th century AD with the onset of the Classic Maya collapse. A significant Maya presence remained in Petén into the Postclassic period after the abandonment of the major Classic period cities; the population was particularly concentrated near permanent water sources. In the early 16th century, when the Spanish discovered the Yucatán Peninsula, the region was still dominated by the Maya civilization. It was divided into a number of independent provinces referred to as kuchkabal (plural kuchkabaloob) in the Yucatec Maya language. The various provinces shared a common culture but the internal sociopolitical organisation varied from one province to the next, as did access to important resources. These differences in political and economic makeup often led to hostilities between the provinces. The politically fragmented state of the Yucatán Peninsula at the time of conquest hindered the Spanish invasion, since there was no central political authority to be overthrown. However, the Spanish were also able to exploit this fragmentation by taking advantage of pre-existing rivalries between polities. Estimates of the number of kuchkabal in the northern Yucatán vary from sixteen to twenty-four. The boundaries between polities were not stable, being subject to the effects of alliances and wars; those kuchkabaloob with more centralised forms of government were likely to have had more stable boundaries than those of loose confederations of provinces. When the Spanish discovered Yucatán, the provinces of Mani and Sotuta were two of the most important polities in the region. They were mutually hostile; the Xiu Maya of Mani allied themselves with the Spanish, while the Cocom Maya of Sotuta became the implacable enemies of the European colonisers. At the time of conquest, polities in the north included Mani, Cehpech and Chakan. Chakan was largely landlocked with a small stretch of coast on the north of the peninsula. Cehpech was a coastal province to its east; further east along the north coast were Ah Kin Chel, Cupul, and Chikinchel. The modern city of Valladolid is situated upon the site of the former capital of Cupul. Cupul and Chinkinchel are known to have been mutually hostile, and to have engaged in wars to control the salt beds of the north coast. Tases was a small landlocked province south of Chikinchel. Ecab was a large province in the east. Uaymil was in the southeast, and Chetumal was to the south of it; all three bordered on the Caribbean Sea. Cochuah was also in the eastern half of the peninsula; it was southwest of Ecab and northwest of Uaymil. Its borders are poorly understood and it may have been landlocked, or have extended to occupy a portion of the Caribbean coast between the latter two kuchkabaloob. The capital of Cochuah was Tihosuco. Hocaba and Sotuta were landlocked provinces north of Mani and southwest of Ah Kin Chel and Cupul. Ah Canul was the northernmost province on the Gulf coast of the peninsula. Canpech (modern Campeche) was to the south of it, followed by Chanputun (modern Champotón). South of Chanputun, and extending west along the Gulf coast was Acalan. This Chontal Maya-speaking province extended east of the Usumacinta River in Tabasco, as far as what is now the southern portion of Campeche state, where their capital was located. In the southern portion of the peninsula, a number of polities occupied the Petén Basin. The Kejache occupied a territory to the north of the Itza and east of Acalan, between the Petén lakes and what is now Campeche, and to the west of Chetumal. The Cholan Maya-speaking Lakandon (not to be confused with the modern inhabitants of Chiapas by that name) controlled territory along the tributaries of the Usumacinta River spanning southwestern Petén in Guatemala and eastern Chiapas. The Lakandon had a fierce reputation amongst the Spanish. Although there is insufficient data to accurately estimate population sizes at the time of contact with the Spanish, early Spanish reports suggest that sizeable Maya populations existed in Petén, particularly around the central lakes and along the rivers. Before their defeat in 1697 the Itza controlled or influenced much of Petén and parts of Belize. The Itza were warlike, and their martial prowess impressed both neighbouring Maya kingdoms and their Spanish enemies. Their capital was Nojpetén, an island city upon Lake Petén Itzá; it has developed into the modern town of Flores, which is the capital of the Petén department of Guatemala. The Itza spoke a variety of Yucatecan Maya. The Kowoj were the second in importance; they were hostile towards their Itza neighbours. The Kowoj were located to the east of the Itza, around the eastern Petén lakes: Lake Salpetén, Lake Macanché, Lake Yaxhá and Lake Sacnab. The Yalain appear to have been one of the three dominant polities in Postclassic central Petén, alongside the Itza and the Kowoj. The Yalain territory had its maximum extension from the east shore of Lake Petén Itzá eastwards to Tipuj in Belize. In the 17th century the Yalain capital was located at the site of that name on the north shore of Lake Macanché. At the time of Spanish contact the Yalain were allied with the Itza, an alliance cemented by intermarriage between the elites of both groups. In the late 17th century Spanish colonial records document hostilities between Maya groups in the lakes region, with the incursion of the Kowoj into former Yalain sites including Zacpeten on Lake Macanché and Ixlu on Lake Salpetén. Other groups in Petén are less well known, and their precise territorial extent and political makeup remains obscure; among them were the Chinamita, the Icaiche, the Kejache, the Lakandon Ch'ol, the Manche Ch'ol, and the Mopan. Impact of Old World diseases A single soldier arriving in Mexico in 1520 was carrying smallpox and thus initiated the devastating plagues that swept through the native populations of the Americas. The European diseases that ravaged the indigenous inhabitants of the Americas also severely affected the various Maya groups of the entire Yucatán Peninsula. Modern estimates of native population decline vary from 75% to 90% mortality. The terrible plagues that swept the peninsula were recorded in Yucatec Maya written histories, which combined with those of neighbouring Maya peoples in the Guatemalan Highlands, suggest that smallpox was rapidly transmitted throughout the Maya area the same year that it arrived in central Mexico with the forces under the command of Pánfilo Narváez. Old World diseases are often mentioned only briefly in indigenous accounts, making it difficult to identify the exact culprit. Among the most deadly were the aforementioned smallpox, influenza, measles and a number of pulmonary diseases, including tuberculosis; the latter disease was attributed to the arrival of the Spanish by the Maya inhabitants of Yucatán. These diseases swept through Yucatán in the 1520s and 1530s, with periodic recurrences throughout the 16th century. By the late 16th century, the reports of high fevers suggest the arrival of malaria in the region, and yellow fever was first reported in the mid-17th century, with a terse mention in the Chilam Balam of Chumayel for 1648. That particular outbreak was traced back to the island of Guadaloupe in the Caribbean, from whence it was introduced to the port city of Campeche, and from there was transmitted to Mérida. Mortality was high, with approximately 50% of the population of some Yucatec Maya settlements being wiped out. Sixteen Franciscan friars are reported to have died in Mérida, probably the majority of the Franciscans based there at the time, and who had probably numbered not much more than twenty before the outbreak. Those areas of the peninsula that experience damper conditions, particularly those possessing swamplands, became rapidly depopulated after the conquest with the introduction of malaria and other waterborne parasites. An example was the one-time well-populated province of Ecab occupying the northeastern portion of the peninsula. In 1528, when Francisco de Montejo occupied the town of Conil for two months, the Spanish recorded approximately 5,000 houses in the town; the adult male population at the time has been conservatively estimated as 3,000. By 1549, Spanish records show that only 80 tributaries were registered to be taxed, indicating a population drop in Conil of more than 90% in 21 years. The native population of the northeastern portion of the peninsula was almost completely eliminated within fifty years of the conquest. In the south, conditions conducive to the spread of malaria existed throughout Petén and Belize. At the time of the fall of Nojpetén in 1697, there are estimated to have been 60,000 Maya living around Lake Petén Itzá, including a large number of refugees from other areas. It is estimated that 88% of them died during the first ten years of colonial rule owing to a combination of disease and war. Likewise, in Tabasco the population of approximately 30,000 was reduced by an estimated 90%, with measles, smallpox, catarrhs, dysentery and fevers being the main culprits. Weaponry, strategies and tactics The Spanish engaged in a strategy of concentrating native populations in newly founded colonial towns, or reducciones (also known as congregaciones). Native resistance to the new nucleated settlements took the form of the flight of the indigenous inhabitants into inaccessible regions such as the forest or joining neighbouring Maya groups that had not yet submitted to the Spanish. Those that remained behind in the reducciones often fell victim to contagious diseases. An example of the effect on populations of this strategy is the province of Acalan, which occupied an area spanning southern Campeche and eastern Tabasco. When Hernán Cortés passed through Acalan in 1525 he estimated the population size as at least 10,000. In 1553 the population was recorded at around 4,000. In 1557 the population was forcibly moved to Tixchel on the Gulf coast, so as to be more easily accessible to the Spanish authorities. In 1561 the Spanish recorded only 250 tribute-paying inhabitants of Tixchel, which probably had a total population of about 1,100. This indicates a 90% drop in population over a 36-year span. Some of the inhabitants had fled Tixchel for the forest, while others had succumbed to disease, malnutrition and inadequate housing in the Spanish reducción. Coastal reducciones, while convenient for Spanish administration, were vulnerable to pirate attacks; in the case of Tixchel, pirate attacks and contagious European diseases led to the eradication of the reducción town and the extinction of the Chontal Maya of Campeche. Among the Maya, ambush was a favoured tactic. Spanish weaponry and armour The 16th-century Spanish conquistadors were armed with broadswords, rapiers, crossbows, matchlocks and light artillery. Mounted conquistadors were armed with a 3.7-metre (12 ft) lance, that also served as a pike for infantrymen. A variety of halberds and bills were also employed. As well as the one-handed broadsword, a 1.7-metre (5.5 ft) long two-handed version was also used. Crossbows had 0.61-metre (2 ft) arms stiffened with hardwoods, horn, bone and cane, and supplied with a stirrup to facilitate drawing the string with a crank and pulley. Crossbows were easier to maintain than matchlocks, especially in the humid tropical climate of the Caribbean region that included much of the Yucatán Peninsula. Native weaponry and armour Maya warriors entered battle against the Spanish with flint-tipped spears, bows and arrows and stones. They wore padded cotton armour to protect themselves. Members of the Maya aristocracy wore quilted cotton armour, and some warriors of lesser rank wore twisted rolls of cotton wrapped around their bodies. Warriors bore wooden or animal hide shields decorated with feathers and animal skins. First encounters: 1502 and 1511 On 30 July 1502, during his fourth voyage, Christopher Columbus arrived at Guanaja, one of the Bay Islands off the coast of Honduras. He sent his brother Bartholomew to scout the island. As Bartholomew explored the island with two boats, a large canoe approached from the west, apparently en route to the island. The canoe was carved from one large tree trunk and was powered by twenty-five naked rowers. Curious as to the visitors, Bartholomew Columbus seized and boarded it. He found it was a Maya trading canoe from Yucatán, carrying well-dressed Maya and a rich cargo that included ceramics, cotton textiles, yellow stone axes, flint-studded war clubs, copper axes and bells, and cacao. Also among the cargo were a small number of women and children, probably destined to be sold as slaves, as were a number of the rowers. The Europeans looted whatever took their interest from amongst the cargo and seized the elderly Maya captain to serve as an interpreter; the canoe was then allowed to continue on its way. This was the first recorded contact between Europeans and the Maya. It is likely that news of the piratical strangers in the Caribbean passed along the Maya trade routes – the first prophecies of bearded invaders sent by Kukulkan, the northern Maya feathered serpent god, were probably recorded around this time, and in due course passed into the books of Chilam Balam. In 1511 the Spanish caravel Santa María de la Barca set sail along the Central American coast under the command of Pedro de Valdivia. The ship was sailing to Santo Domingo from Darién to inform the colonial authorities there of ongoing conflict between conquistadors Diego de Nicuesa and Vasco Nuñez de Balboa in Darién. The ship foundered upon a reef known as Las Víboras ("The Vipers") or, alternatively, Los Alacranes ("The Scorpions"), somewhere off Jamaica. There were just twenty survivors from the wreck, including Captain Valdivia, Gerónimo de Aguilar and Gonzalo Guerrero. They set themselves adrift in one of the ship's boats, with bad oars and no sail; after thirteen days during which half of the survivors died, they made landfall upon the coast of Yucatán. There they were seized by Halach Uinik, a Maya lord. Captain Vildivia was sacrificed with four of his companions, and their flesh was served at a feast. Aguilar and Guerrero were held prisoner and fattened for killing, together with five or six of their shipmates. Aguilar and Guerrero managed to escape their captors and fled to a neighbouring lord who was an enemy of Halach Uinik; he took them prisoner and kept them as slaves. After a time, Gonzalo Guerrero was passed as a slave to the lord Nachan Can of Chetumal. Guerrero became completely Mayanised and served his new lord with such loyalty that he was married to one of Nachan Chan's daughters, Zazil Ha, by whom he had three children. By 1514, Guerrero had achieved the rank of nacom, a war leader who served against Nachan Chan's enemies. Francisco Hernández de Córdoba, 1517 In 1517, Francisco Hernández de Córdoba set sail from Cuba with a small fleet, consisting of two caravels and a brigantine, with the dual intention of exploration and of rounding up slaves. The experienced Antón de Alaminos served as pilot; he had previously served as pilot under Christopher Columbus on his final voyage. Also among the approximately 100-strong expedition members was Bernal Díaz del Castillo. The expedition sailed west from Cuba for three weeks, and weathered a two-day storm a week before sighting the coast of the northeastern tip of the Yucatán Peninsula. The ships could not put in close to the shore due to the shallowness of the coastal waters. However, they could see a Maya city some two leagues inland, upon a low hill. The Spanish called it Gran Cairo (literally "Great Cairo") due to its size and its pyramids. Although the location is not now known with certainty, it is believed that this first sighting of Yucatán was at Isla Mujeres. The following morning, the Spanish sent the two ships with a shallower draught to find a safe approach through the shallows. The caravels anchored about one league from the shore. Ten large canoes powered by both sails and oars rowed out to meet the Spanish ships. Over thirty Maya boarded the vessels and mixed freely with the Spaniards. The Maya visitors accepted gifts of beads, and the leader indicated with signs that they would return to take the Spanish ashore the following day. The Maya leader returned the following day with twelve canoes, as promised. The Spanish could see from afar that the shore was packed with natives. The conquistadors put ashore in the brigantine and the ships' boats; a few of the more daring Spaniards boarded the native canoes. The Spanish named the headland Cape Catoche, after some words spoken by the Maya leader, which sounded to the Spanish like cones catoche. Once ashore, the Spaniards clustered loosely together and advanced towards the city along a path among low, scrub-covered hillocks. At this point the Maya leader gave a shout and the Spanish party was ambushed by Maya warriors armed with spears, bows and arrows, and stones. Thirteen Spaniards were injured by arrows in the first assault, but the conquistadors regrouped and repulsed the Maya attack. They advanced to a small plaza bordered by temples upon the outskirts of the city. When the Spaniards ransacked the temples they found a number of low-grade gold items, which filled them with enthusiasm. The expedition captured two Mayas to be used as interpreters and retreated to the ships. Over the following days the Spanish discovered that although the Maya arrows had struck with little force, the flint arrowheads tended to shatter on impact, causing infected wounds and a slow death; two of the wounded Spaniards died from the arrow-wounds inflicted in the ambush. Over the next fifteen days the fleet slowly followed the coastline west, and then south. The casks brought from Cuba were leaking and the expedition was now running dangerously low on fresh water; the hunt for more became an overriding priority as the expedition advanced, and shore parties searching for water were left dangerously exposed because the ships could not pull close to the shore due to the shallows. On 23 February 1517, the day of Saint Lazarus, another city was spotted and named San Lázaro by the Spanish – it is now known by its original Maya name, Campeche. A large contingent put ashore in the brigantine and the ships' boats to fill their water casks in a freshwater pool. They were approached by about fifty finely dressed and unarmed Indians while the water was being loaded into the boats; they questioned the Spaniards as to their purpose by means of signs. The Spanish party then accepted an invitation to enter the city. They were led amongst large buildings until they stood before a blood-caked altar, where many of the city's inhabitants crowded around. The Indians piled reeds before the visitors; this act was followed by a procession of armed Maya warriors in full war paint, followed by ten Maya priests. The Maya set fire to the reeds and indicated that the Spanish would be killed if they were not gone by the time the reeds had been consumed. The Spanish party withdrew in defensive formation to the shore and rapidly boarded their boats to retreat to the safety of the ships. The small fleet continued for six more days in fine weather, followed by four stormy days. By this time water was once again dangerously short. The ships spotted an inlet close to another city, Champotón, and a landing party discovered fresh water. Armed Maya warriors approached from the city while the water casks were being filled. Communication was once again attempted with signs. Night fell by the time the water casks had been filled and the attempts at communication concluded. In the darkness the Spaniards could hear the movements of large numbers of Maya warriors. They decided that a night-time retreat would be too risky; instead, they posted guards and waited for dawn. At sunrise, the Spanish saw that they had been surrounded by a sizeable army. The massed Maya warriors launched an assault with missiles, including arrows, darts and stones; they then charged into hand-to-hand combat with spears and clubs. Eighty of the defenders were wounded in the initial barrage of missiles, and two Spaniards were captured in the frantic melee that followed. All of the Spanish party received wounds, including Hernández de Córdoba. The Spanish regrouped in a defensive formation and forced passage to the shore, where their discipline collapsed and a frantic scramble for the boats ensued, leaving the Spanish vulnerable to the pursuing Maya warriors who waded into the sea behind them. Most of the precious water casks were abandoned on the beach. When the surviving Spanish reached the safety of the ships, they realised that they had lost over fifty men, more than half their number. Five men died from their wounds in the following days. The battle had lasted only an hour, and the Spanish named the locale as the Coast of the Disastrous Battle. They were now far from help and low on supplies; too many men had been lost and injured to sail all three ships back to Cuba. They decided to abandon their smallest ship, the brigantine, although it was purchased on credit from Governor Velásquez of Cuba. The few men who had not been wounded because they were manning the ships during the battle were reinforced with three men who had suffered relatively minor wounds; they put ashore at a remote beach to dig for water. They found some and brought it back to the ships, although it sickened those who drank it. The two ships sailed through a storm for two days and nights; Alaminos, the pilot, then steered a course for Florida, where they found good drinking water, although they lost one man to the local Indians and another drank so much water that he died. The ships finally made port in Cuba, where Hernández de Cordóba wrote a report to Governor Velázquez describing the voyage, the cities, the plantations, and, most importantly, the discovery of gold. Hernández died soon after from his wounds. The two captured Maya survived the voyage to Cuba and were interrogated; they swore that there was abundant gold in Yucatán. Juan de Grijalva, 1518 Diego Velázquez, the governor of Cuba, was enthused by Hernández de Córdoba's report of gold in Yucatán. He organised a new expedition consisting of four ships and 260 men. He placed his nephew Juan de Grijalva in command. Francisco de Montejo, who would eventually conquer much of the peninsula, was captain of one of the ships,; Pedro de Alvarado and Alonso d'Avila captained the other ships. Bernal Díaz del Castillo served on the crew; he was able to secure a place on the expedition as a favour from the governor, who was his kinsman. Antón de Alaminos once again served as pilot. Governor Velázquez provided all four ships, in an attempt to protect his claim over the peninsula. The small fleet was stocked with crossbows, muskets, barter goods, salted pork and cassava bread. Grijalva also took one of the captured Indians from the Hernández expedition. The fleet left Cuba in April 1518, and made its first landfall upon the island of Cozumel, off the east coast of Yucatán. The Maya inhabitants of Cozumel fled the Spanish and would not respond to Grijalva's friendly overtures. The fleet sailed south from Cozumel, along the east coast of the peninsula. The Spanish spotted three large Maya cities along the coast, one of which was probably Tulum. On Ascension Thursday the fleet discovered a large bay, which the Spanish named Bahía de la Ascensión. Grijalva did not land at any of these cities and turned back north from Ascensión Bay. He looped around the north of the Yucatán Peninsula to sail down the west coast. At Campeche the Spanish tried to barter for water but the Maya refused, so Grijalva opened fire against the city with small cannon; the inhabitants fled, allowing the Spanish to take the abandoned city. Messages were sent with a few Maya who had been too slow to escape but the Maya remained hidden in the forest. The Spanish boarded their ships and continued along the coast. At Champotón, where the inhabitants had routed Hernández and his men, the fleet was approached by a small number of large war canoes, but the ships' cannon soon put them to flight. At the mouth of the Tabasco River the Spanish sighted massed warriors and canoes but the natives did not approach. By means of interpreters, Grijalva indicated that he wished to trade and bartered wine and beads in exchange for food and other supplies. From the natives they received a few gold trinkets and news of the riches of the Aztec Empire to the west. The expedition continued far enough to confirm the reality of the gold-rich empire, sailing as far north as Pánuco River. As the fleet returned to Cuba, the Spanish attacked Champotón to avenge the previous year's defeat of the Spanish expedition led by Hernández. One Spaniard was killed and fifty were wounded in the ensuing battle, including Grijalva. Grijalva put into the port of Havana five months after he had left. Hernán Cortés, 1519 Grijalva's return aroused great interest in Cuba, and Yucatán was believed to be a land of riches waiting to be plundered. A new expedition was organised, with a fleet of eleven ships carrying 500 men and some horses. Hernán Cortés was placed in command, and his crew included officers that would become famous conquistadors, including Pedro de Alvarado, Cristóbal de Olid, Gonzalo de Sandoval and Diego de Ordaz. Also aboard were the Francisco de Montejo and Bernal Díaz del Castillo, veterans of the Grijalva expedition. The fleet made its first landfall at Cozumel, and Cortés remained there for several days. Maya temples were cast down and a Christian cross was put up on one of them. At Cozumel Cortés heard rumours of bearded men on the Yucatán mainland, who he presumed were Europeans. Cortés sent out messengers to them and was able to rescue the shipwrecked Gerónimo de Aguilar, who had been enslaved by a Maya lord. Aguilar had learnt the Yucatec Maya language and became Cortés' interpreter. From Cozumel, the fleet looped around the north of the Yucatán Peninsula and followed the coast to the Tabasco River, which Cortés renamed as the Grijalva River in honour of the Spanish captain who had discovered it. In Tabasco, Cortés anchored his ships at Potonchán, a Chontal Maya town. The Maya prepared for battle but the Spanish horses and firearms quickly decided the outcome. The defeated Chontal Maya lords offered gold, food, clothing and a group of young women in tribute to the victors. Among these women was a young Maya noblewoman called Malintzin, who was given the Spanish name Marina. She spoke Maya and Nahuatl and became the means be which Cortés was able to communicate with the Aztecs. Marina became Cortés' consort and eventually bore him a son. From Tabasco, Cortés continued to Cempoala in Veracruz, a subject city of the Aztec Empire, and from there went on to conquer the Aztecs. In 1519 Cortés sent the veteran Francisco de Montejo back to Spain with treasure for the king. While he was in Spain he pleaded Cortés' cause against the supporters of Diego de Velasquez. Montejo remained in Spain for seven years, and eventually succeeded in acquiring the hereditary military title of adelantado. Hernán Cortés in the Maya lowlands, 1524–25 In 1524, after the Spanish conquest of the Aztec Empire, Hernán Cortés led an expedition to Honduras over land, cutting across Acalan in southern Campeche and the Itza kingdom in what is now the northern Petén Department of Guatemala. His aim was to subdue the rebellious Cristóbal de Olid, whom he had sent to conquer Honduras; Olid had, however, set himself up independently on his arrival in that territory. Cortés left Tenochtitlan on 12 October 1524 with 140 Spanish soldiers, 93 of them mounted, 3,000 Mexican warriors, 150 horses, a herd of pigs, artillery, munitions and other supplies. He also had with him the captured Aztec emperor Cuauhtemoc, and Cohuanacox and Tetlepanquetzal, the captive Aztec lords of Texcoco and Tlacopan. Cortés marched into Maya territory in Tabasco; the army crossed the Usumacinta River near Tenosique and crossed into the Chontal Maya province of Acalan, where he recruited 600 Chontal Maya carriers. In Acalan, Cortés believed that the captive Aztec lords were plotting against him and he ordered Cuauhtemoc and Tetlepanquetzal to be hanged. Cortés and his army left Acalan on 5 March 1525. The expedition passed onwards through Kejache territory and reported that the Kejache towns were situated in easily defensible locations and were often fortified. One of these was built on a rocky outcrop near a lake and a river that fed into it. The town was fortified with a wooden palisade and was surrounded by a moat. Cortés reported that the town of Tiac was even larger and was fortified with walls, watchtowers and earthworks; the town itself was divided into three individually fortified districts. Tiac was said to have been at war with the unnamed smaller town. The Kejache claimed that their towns were fortified against the attacks of their aggressive Itza neighbours. They arrived at the north shore of Lake Petén Itzá on 13 March 1525. The Roman Catholic priests accompanying the expedition celebrated mass in the presence of Aj Kan Ek', the king of the Itza, who was said to be so impressed that he pledged to worship the cross and to destroy his idols. Cortés accepted an invitation from Kan Ek' to visit Nojpetén (also known as Tayasal), and crossed to the Maya city with 20 Spanish soldiers while the rest of his army continued around the lake to meet him on the south shore. On his departure from Nojpetén, Cortés left behind a cross and a lame horse that the Itza treated as a deity, attempting to feed it poultry, meat and flowers, but the animal soon died. The Spanish did not officially contact the Itza again until the arrival of Franciscan priests in 1618, when Cortés' cross was said to still be standing at Nojpetén. From the lake, Cortés continued south along the western slopes of the Maya Mountains, a particularly arduous journey that took 12 days to cover 32 kilometres (20 mi), during which he lost more than two-thirds of his horses. When he came to a river swollen with the constant torrential rains that had been falling during the expedition, Cortés turned upstream to the Gracias a Dios rapids, which took two days to cross and cost him more horses. On 15 April 1525 the expedition arrived at the Maya village of Tenciz. With local guides they headed into the hills north of Lake Izabal, where their guides abandoned them to their fate. The expedition became lost in the hills and came close to starvation before they captured a Maya boy who led them to safety. Cortés found a village on the shore of Lake Izabal, perhaps Xocolo. He crossed the Dulce River to the settlement of Nito, somewhere on the Amatique Bay, with about a dozen companions, and waited there for the rest of his army to regroup over the next week. By this time the remnants of the expedition had been reduced to a few hundred; Cortés succeeded in contacting the Spaniards he was searching for, only to find that Cristóbal de Olid's own officers had already put down his rebellion. Cortés then returned to Mexico by sea. Francisco de Montejo, 1527–28 The richer lands of Mexico engaged the main attention of the Conquistadors for some years, then in 1526 Francisco de Montejo (a veteran of the Grijalva and Cortés expeditions) successfully petitioned the King of Spain for the right to conquer Yucatán. On 8 December of that year he was issued with the hereditary military title of adelantado and permission to colonise the Yucatán Peninsula. In 1527 he left Spain with 400 men in four ships, with horses, small arms, cannon and provisions. He set sail for Santo Domingo, where more supplies and horses were collected, allowing Montejo to increase his cavalry to fifty. One of the ships was left at Santo Domingo as a supply ship to provide later support; the other ships set sail and reached Cozumel, an island off the east coast of Yucatán, in the second half of September 1527. Montejo was received in peace by the lord of Cozumel, Aj Naum Pat, but the ships only stopped briefly before making for the Yucatán coast. The expedition made landfall somewhere near Xelha in the Maya province of Ekab, in what is now Mexico's Quintana Roo state. Montejo garrisoned Xelha with 40 soldiers under his second-in-command, Alonso d'Avila, and posted 20 more at nearby Pole. Xelha was renamed Salamanca de Xelha and became the first Spanish settlement in the peninsula. The provisions were soon exhausted and additional food was seized from the local Maya villagers; this too was soon consumed. Many local Maya fled into the forest and Spanish raiding parties scoured the surrounding area for food, finding little. With discontent growing among his men, Montejo took the drastic step of burning his ships; this strengthened the resolve of his troops, who gradually acclimatised to the harsh conditions of Yucatán. Montejo was able to get more food from the still-friendly Aj Nuam Pat, when the latter made a visit to the mainland. Montejo took 125 men and set out on an expedition to explore the north-eastern portion of the Yucatán peninsula. His expedition passed through the towns of Xamanha, Mochis and Belma, none of which survives today.[nb 1] At Belma, Montejo gathered the leaders of the nearby Maya towns and ordered them to swear loyalty to the Spanish Crown. After this, Montejo led his men to Conil, a town in Ekab that was described as having 5,000 houses, where the Spanish party halted for two months. In the spring of 1528, Montejo left Conil for the city of Chauaca, which was abandoned by its Maya inhabitants under cover of darkness. The following morning the inhabitants attacked the Spanish party but were defeated. The Spanish then continued to Ake, some 16 kilometres (9.9 mi) north of Tizimín, where they engaged in a major battle against the Maya, which killed more than 1,200 Maya. After this Spanish victory, the neighbouring Maya leaders all surrendered. Montejo's party then continued to Sisia and Loche before heading back to Xelha. Montejo arrived at Xelha with only 60 of his party, and found that only 12 of his 40-man garrison survived, while the garrison at Pole had been entirely wiped out. The support ship eventually arrived from Santo Domingo, and Montejo used it to sail south along the coast, while he sent D'Avila via land. Montejo discovered the thriving port city of Chaktumal (modern Chetumal). At Chaktumal, Montejo learnt that shipwrecked Spanish sailor Gonzalo de Guerrero was in the region, and Montejo sent messages to him, inviting him to return to join his compatriots, but Guerrero declined. The Maya at Chaktumal fed false information to the Spanish, and Montejo was unable to find d'Avila and link up with him. D'Avila returned overland to Xelha, and transferred the fledgling Spanish colony to nearby Xamanha, modern Playa del Carmen, which Montejo considered to be a better port. After waiting for d'Avila without result, Montejo sailed south as far as the Ulúa River in Honduras before turning around and heading back up the coast to finally meet up with his lieutenant at Xamanha. Late in 1528, Montejo left d'Avila to oversee Xamanha and sailed north to loop around the Yucatán Peninsula and head for the Spanish colony of New Spain in central Mexico. Francisco de Montejo and Alonso d'Avila, 1531–35 Montejo was appointed alcalde mayor (a local colonial governor) of Tabasco in 1529, and pacified that province with the aid of his son, also named Francisco de Montejo. D'Avila was sent from eastern Yucatán to conquer Acalan, which extended southeast of the Laguna de Terminos. Montejo the Younger founded Salamanca de Xicalango as a base of operations. In 1530 D'Avila established Salamanca de Acalán as a base from which to launch new attempts to conquer Yucatán. Salamanca de Acalán proved a disappointment, with no gold for the taking and with lower levels of population than had been hoped. D'Avila soon abandoned the new settlement and set off across the lands of the Kejache to Champotón, arriving there towards the end of 1530. During a colonial power struggle in Tabasco, the elder Montejo was imprisoned for a time. Upon his release he met up with his son in Xicalango, Tabasco, and they then both rejoined d'Avila at Champotón. In 1531 Montejo moved his base of operations to Campeche. Alonso d'Avila was sent overland to Chauaca in the east of the peninsula, passing through Maní where he was well received by the Xiu Maya. D'Avila continued southeast to Chetumal where he founded the Spanish town of Villa Real ("Royal Town"). The local Maya fiercely resisted the placement of the new Spanish colony and d'Avila and his men were forced to abandon Villa Real and make for Honduras in canoes. At Campeche, the Maya amassed a strong force and attacked the city; the Spanish were able to fight them off, although the elder Montejo was almost killed. Aj Canul, the lord of the attacking Maya, surrendered to the Spanish. After this battle, the younger Francisco de Montejo was despatched to the northern Cupul province, where the lord Naabon Cupul reluctantly allowed him to found the Spanish town of Ciudad Real at Chichen Itza. Montejo carved up the province amongst his soldiers and gave each of his men two to three thousand Maya in encomienda. After six months of Spanish rule, Cupul dissatisfaction could no longer be contained and Naabon Cupul was killed during a failed attempt to kill Montejo the Younger. The death of their lord only served to inflame Cupul anger and, in mid 1533, they laid siege to the small Spanish garrison at Chichen Itza. Montejo the Younger abandoned Ciudad Real by night after arranging a distraction for their attackers, and he and his men fled west, where the Chel, Pech and Xiu provinces remained obedient to Spanish rule. Montejo the Younger was received in friendship by Namux Chel, the lord of the Chel province, at Dzilam. In the spring of 1534 he rejoined his father in the Chakan province at Dzikabal, near T'ho (the modern city of Mérida). While his son had been attempting to consolidate the Spanish control of Cupul, Francisco de Montejo the Elder had met the Xiu ruler at Maní. The Xiu Maya maintained their friendship with the Spanish throughout the conquest and Spanish authority was eventually established over Yucatán in large part due to Xiu support. The Montejos, after reuniting at Dzikabal, founded a new Spanish town at Dzilam, although the Spanish suffered hardships there. Montejo the Elder returned to Campeche, where he was received with friendship by the local Maya. He was accompanied by the friendly Chel lord Namux Chel, who travelled on horseback, and two of the lord's cousins, who were taken in chains. Montejo the Younger remained behind in Dzilam to continue his attempts at conquest of the region but, finding the situation too difficult, he soon retreated to Campeche to rejoin his father and Alonso d'Avila, who had returned to Campeche shortly before Montejo the Younger. Around this time the news began to arrive of Francisco Pizarro's conquests in Peru and the rich plunder that his soldiers were taking there, undermining the morale of Montejo's already disenchanted band of followers. Montejo's soldiers began to abandon him to seek their fortune elsewhere; in seven years of attempted conquest in the northern provinces of the Yucatán Peninsula, very little gold had been found. Towards the end of 1534 or the beginning of the next year, Montejo the Elder and his son retreated from Campeche to Veracruz, taking their remaining soldiers with them. Montejo the Elder became embroiled in colonial infighting over the right to rule Honduras, a claim that put him in conflict with Pedro de Alvarado, captain general of Guatemala, who also claimed Honduras as part of his jurisdiction. Alvarado was ultimately to prove successful. In Montejo the Elder's absence, first in central Mexico, and then in Honduras, Montejo the Younger acted as lieutenant governor and captain general in Tabasco. Conflict at Champoton The Franciscan friar Jacobo de Testera arrived in Champoton in 1535 to attempt the peaceful incorporation of Yucatán into the Spanish Empire. Testera had been assured by the Spanish authorities that no military activity would be undertaken in Yucatán while he was attempting its conversion to the Roman Catholic faith, and that no soldiers would be permitted to enter the peninsula. His initial efforts were proving successful when Captain Lorenzo de Godoy arrived in Champoton at the command of soldiers despatched there by Montejo the Younger. Godoy and Testera were soon in conflict and the friar was forced to abandon Champoton and return to central Mexico. Godoy's attempt to subdue the Maya around Champoton was unsuccessful and the local Kowoj Maya resisted his attempts to assert Spanish dominance of the region. This resistance was sufficiently tenacious that Montejo the Younger sent his cousin from Tabasco to Champoton to take command. His diplomatic overtures to the Champoton Kowoj were successful and they submitted to Spanish rule. Champoton was the last Spanish outpost in the Yucatán Peninsula; it was increasingly isolated and the situation there became difficult. Conquest and settlement in northern Yucatán, 1540–46 In 1540 Montejo the Elder, who was now in his late 60s, turned his royal rights to colonise Yucatán over to his son, Francisco Montejo the Younger. In early 1541 Montejo the Younger joined his cousin in Champton; he did not remain there long, and quickly moved his forces to Campeche. Once there Montejo the Younger, commanding between three and four hundred Spanish soldiers, established the first permanent Spanish town council in the Yucatán Peninsula. Shortly after establishing the Spanish presence in Campeche, Montejo the Younger summoned the local Maya lords and commanded them to submit to the Spanish Crown. A number of lords submitted peacefully, including the ruler of the Xiu Maya. The lord of the Canul Maya refused to submit and Montejo the Younger sent his cousin against them; Montejo himself remained in Campeche awaiting reinforcements. Montejo the Younger's cousin met the Canul Maya at Chakan, not far from T'ho. On 6 January 1542 he founded the second permanent town council, calling the new colonial town Mérida. On 23 January, Tutul Xiu, the lord of Mani, approached the Spanish encampment at Mérida in peace, bearing sorely needed food supplies. He expressed interest in the Spanish religion and witnessed a Roman Catholic mass celebrated for his benefit. Tutul Xiu was greatly impressed and converted to the new religion; he was baptised as Melchor and stayed with the Spanish at Mérida for two months, receiving instruction in the Catholic faith. Tutul Xiu was the ruler of the most powerful province of northern Yucatán and his submission to Spain and conversion to Christianity had repercussions throughout the peninsula, and encouraged the lords of the western provinces of the peninsula to accept Spanish rule. The eastern provinces continued to resist Spanish overtures. Montejo the Younger then sent his cousin to Chauaca where most of the eastern lords greeted him in peace. The Cochua Maya resisted fiercely but were soon defeated by the Spanish. The Cupul Maya also rose up against the newly imposed Spanish domination, but their opposition was quickly put down. Montejo continued to the eastern Ekab province, reaching the east coast at Pole. Stormy weather prevented the Spanish from crossing to Cozumel, and nine Spaniards were drowned in the attempted crossing. A further Spaniard was killed by hostile Maya. Rumours of this setback grew in the telling and both the Cupul and Cochua provinces once again rose up against their would-be European overlords. The Spanish hold on the eastern portion of the peninsula remained tenuous and a number of Maya polities remained independent, including Chetumal, Cochua, Cupul, Sotuta and the Tazes. On 8 November 1546 and alliance of eastern provinces launched a coordinated uprising against the Spanish. The provinces of Cupul, Cochua, Sotuta, Tazes, Uaymil, Chetumal and Chikinchel united in a concerted effort to drive the invaders from the peninsula; the uprising lasted four months. Eighteen Spaniards were surprised in the eastern towns, and were sacrificed. A contemporary account described the slaughter of over 400 allied Maya, as well as livestock. Mérida and Campeche were forewarned of the impending attack; Montejo the Younger and his cousin were in Campeche. Montejo the Elder arrived in Mérida from Chiapas in December 1546, with reinforcements gathered from Champoton and Campeche. The rebellious eastern Maya were finally defeated in a single battle, in which twenty Spaniards and several hundred allied Maya were killed. This battle marked the final conquest of the northern portion of the Yucatán Peninsula. As a result of the uprising and the Spanish response, many of the Maya inhabitants of the eastern and southern territories fled to the still unconquered Petén Basin, in the extreme south. The Spanish only achieved dominance in the north and the polities of Petén remained independent and continued to receive many refugees from the north. Petén Basin, 1618–97 The Petén Basin covers an area that is now part of Guatemala; in colonial times it originally fell under the jurisdiction of the Governor of Yucatán, before being transferred to the jurisdiction of the Audiencia Real of Guatemala in 1703. The Itza kingdom centred upon Lake Petén Itzá had been visited by Hernán Cortés on his march to Honduras in 1525. Early 17th century Following Cortés' visit, no Spanish attempted to visit the warlike Itza inhabitants of Nojpetén for almost a hundred years. In 1618 two Franciscan friars set out from Mérida on a mission to attempt the peaceful conversion of the still-pagan Itza in central Petén. Bartolomé de Fuensalida and Juan de Orbita were accompanied by some Christianised Maya. After an arduous six-month journey the travellers were well received at Nojpetén by the current Kan Ek'. They stayed for some days in an attempt to evangelise the Itza, but the Aj Kan Ek' refused to renounce his Maya religion, although he showed interest in the masses held by the Catholic missionaries. Attempts to convert the Itza failed, and the friars left Nojpetén on friendly terms with Kan Ek'. The friars returned in October 1619, and again Kan Ek' welcomed them in a friendly manner, but this time the Maya priesthood were hostile and the missionaries were expelled without food or water, but survived the journey back to Mérida. In March 1622, the governor of Yucatán, Diego de Cardenas, ordered Captain Francisco de Mirones Lezcano to launch an assault upon the Itza; he set out from Yucatán with 20 Spanish soldiers and 80 Mayas from Yucatán. His expedition was later joined by Franciscan friar Diego Delgado. In May the expedition advanced to Sakalum, southwest of Bacalar, where there was a lengthy delay while they waited for reinforcements. En route to Nojpetén, Delgado believed that the soldiers' treatment of the Maya was excessively cruel, and he left the expedition to make his own way to Nojpetén with eighty Christianised Maya from Tipuj in Belize. In the meantime the Itza had learnt of the approaching military expedition and had become hardened against further Spanish missionary attempts. When Mirones learnt of Delgado's departure, he sent 13 soldiers to persuade him to return or continue as his escort should he refuse. The soldiers caught up with him just before Tipuj, but he was determined to reach Nojpetén. From Tipuj, Delgado sent a messenger to Kan Ek', asking permission to travel to Nojpetén; the Itza king replied with a promise of safe passage for the missionary and his companions. The party was initially received in peace at the Itza capital, but as soon as the Spanish soldiers let their guard down, the Itza seized and bound the new arrivals. The soldiers were sacrificed to the Maya gods. After their sacrifice, the Itza took Delgado, cut his heart out and dismembered him; they displayed his head on a stake with the others. The fortune of the leader of Delgado's Maya companions was no better. With no word from Delgado's escort, Mirones sent two Spanish soldiers with a Maya scout to learn their fate. When they arrived upon the shore of Lake Petén Itzá, the Itza took them across to their island capital and imprisoned them. Bernardino Ek, the scout, escaped and returned to Mirones with the news. Soon afterwards, on 27 January 1624, an Itza war party led by AjK'in P'ol caught Mirones and his soldiers off guard and unarmed in the church at Sakalum, and killed them all. Spanish reinforcements arrived too late. A number of local Maya men and women were killed by Spanish attackers, who also burned the town. Following these killings, Spanish garrisons were stationed in several towns in southern Yucatán, and rewards were offered for the whereabouts of AjK'in P'ol. The Maya governor of Oxkutzcab, Fernando Kamal, set out with 150 Maya archers to track the warleader down; they succeeded in capturing the Itza captain and his followers, together with silverware from the looted Sakalum church and items belonging to Mirones. The prisoners were taken back to the Spanish Captain Antonio Méndez de Canzo, interrogated under torture, tried, and condemned to be hanged, drawn and quartered. They were decapitated, and the heads were displayed in the plazas of towns throughout the colonial Partido de la Sierra in what is now Mexico's Yucatán state. These events ended all Spanish attempts to contact the Itza until 1695. In the 1640s internal strife in Spain distracted the government from attempts to conquer unknown lands; the Spanish Crown lacked the time, money or interest in such colonial adventures for the next four decades. Late 17th century Part of a series on the |History of New Spain| |New Spain portal| In 1692 Basque nobleman Martín de Ursúa y Arizmendi proposed to the Spanish king the construction of a road from Mérida southwards to link with the Guatemalan colony, in the process "reducing" any independent native populations into colonial congregaciones; this was part of a greater plan to subjugate the Lakandon and Manche Ch'ol of southern Petén and the upper reaches of the Usumacinta River. The original plan was for the province of Yucatán to build the northern section and for Guatemala to build the southern portion, with both meeting somewhere in Ch'ol territory; the plan was later modified to pass further east, through the kingdom of the Itza. The governor of Yucatán, Martín de Ursúa y Arizmendi, began to build the road from Campeche south towards Petén. At the beginning of March 1695, Captain Alonso García de Paredes led a group of 50 Spanish soldiers, accompanied by native guides, muleteers and labourers. The expedition advanced south into Kejache territory, which began at Chunpich, about 5 kilometres (3.1 mi) north of the modern border between Mexico and Guatemala. He rounded up some natives to be moved into colonial settlements, but met with armed Kejache resistance. García decided to retreat around the middle of April. In March 1695, Captain Juan Díaz de Velasco set out from Cahabón in Alta Verapaz, Guatemala, with 70 Spanish soldiers, accompanied by a large number of Maya archers from Verapaz, native muleteers, and four Dominican friars. The Spanish pressed ahead to Lake Petén Itzá and engaged in a series of fierce skirmishes with Itza hunting parties. At the lakeshore, within sight of Nojpetén, the Spanish encountered such a large force of Itzas that they retreated south, back to their main camp. Interrogation of an Itza prisoner revealed that the Itza kingdom was in a state of high alert to repel the Spanish; the expedition almost immediately withdrew back to Cahabón. In mid-May 1695 García again marched southwards from Campeche, with 115 Spanish soldiers and 150 Maya musketeers, plus Maya labourers and muleteers; the final tally was more than 400 people, which was regarded as a considerable army in the impoverished Yucatán province. Ursúa also ordered two companies of Maya musketeers from Tek'ax and Oxk'utzkab' to join the expedition at B'olonch'en Kawich, some 60 kilometres (37 mi) southeast of the city of Campeche. At the end of May three friars were assigned to join the Spanish force, accompanied by a lay brother. A second group of Franciscans would continue onwards independently to Nojpetén to make contact with the Itzas; it was led by friar Andrés de Avendaño, who was accompanied by another friar and a lay brother. García ordered the construction of a fort at Chuntuki, some 25 leagues (approximately 65 miles or 105 km) north of Lake Petén Itzá, which would serve as the main military base for the Camino Real ("Royal Road") project. The Sajkab'chen company of native musketeers pushed ahead with the road builders from Tzuktzok' to the first Kejache town at Chunpich, which the Kejache had fled. The company's officers sent for reinforcements from García at Tzuktok' but before any could arrive some 25 Kejache returned to Chunpich with baskets to collect their abandoned food. The nervous Sajkab'chen sentries feared that the residents were returning en masse and discharged their muskets at them, with both groups then retreating. The musketeer company then arrived to reinforce their sentries and charged into battle against approaching Kejache archers. Several musketeers were injured in the ensuing skirmish and, the Kejache retreated along a forest path without injury. The Sajkab'chen company followed the path and found two more deserted settlements with large amounts of abandoned food. They seized the food and retreated back along the path. Around 3 August García moved his entire army forward to Chunpich, and by October Spanish soldiers had established themselves near the source of the San Pedro River. By November Tzuktok' was garrisoned with 86 soldiers and more at Chuntuki. In December 1695 the main force was reinforced with 250 soldiers, of which 150 were Spanish and pardo and 100 were Maya, together with labourers and muleteers. Avendaño's expedition, June 1695 In May 1695 Antonio de Silva had appointed two groups of Franciscans to head for Petén; the first group was to join up with García's military expedition. The second group was to head for Lake Petén Itza independently. This second group was headed by friar Andrés de Avendaño. Avendaño was accompanied by another friar, a lay brother, and six Christian Maya. This latter group left Mérida on 2 June 1695. Avendaño continued south along the course of the new road, finding increasing evidence of Spanish military activity. The Franciscans overtook García at B'uk'te, about 12 kilometres (7.5 mi) before Tzuktok'. On 3 August García advanced to Chunpich but tried to persuade Avendaño to stay behind to minister to the prisoners from B'uk'te. Avendaño instead split his group and left in secret with just four Christian Maya companions, seeking the Chunpich Kejache that had attacked one of García's advance companies and had now retreated into the forest. He was unable to find the Kejache but did manage to get information regarding a path that led southwards to the Itza kingdom. Avendaño returned to Tzuktok' and reconsidered his plans; the Franciscans were short of supplies, and the forcefully congregated Maya that they were charged with converting were disappearing back into the forest daily. Antonio de Silva ordered Avendaño to return to Mérida, and he arrived there on 17 September 1695. Meanwhile, the other group of Franciscans, led by Juan de San Buenaventura Chávez, continued following the roadbuilders into Kejache territory, through IxB'am, B'atkab' and Chuntuki (modern Chuntunqui near Carmelita, Petén). Juan de San Buenaventura's small group of Franciscans arrived in Chuntuki on 30 August 1695, and found that the army had opened the road southwards for another seventeen leagues (approximately 44.2 miles or 71.1 km), almost half way to Lake Petén Itzá, but returned to Chuntuki due to the seasonal rains. San Buenaventura was accompanied by two friars and a lay brother. With Avendaño's return to Mérida, provincial superior Antonio de Silva despatched two additional friars to join San Buenaventura's group. One of these was to convert the Kejache in Tzuktok', and the other was to do the same at Chuntuki. On 24 October San Buenaventura wrote to the provincial superior reporting that the warlike Kejache were now pacified and that they had told him that the Itza were ready to receive the Spanish in friendship. On that day 62 Kejache men had voluntarily come to Chuntuki from Pak'ek'em, where another 300 Kejache resided. In early November 1695, friar Tomás de Alcoser and brother Lucas de San Francisco were sent to establish a mission at Pak'ek'em, where they were well received by the cacique (native chief) and his pagan priest. Pak'ek'em was sufficiently far from the new Spanish road that it was free from military interference, and the friars oversaw the building of a church in what was the largest mission town in Kejache territory. A second church was built at B'atkab' to attend to over 100 K'ejache refugees who had been gathered there under the stewardship of a Spanish friar; a further church was established at Tzuktok', overseen by another friar. Avendaño's expedition, December 1695 – January 1696 Franciscan Andrés de Avendaño left Mérida on 13 December 1695, and arrived in Nojpetén around 14 January 1696, accompanied by four companions. From Chuntuki they followed an Indian trail that led them past the source of the San Pedro River and across steep karst hills to a watering hole by some ruins. From there they followed the small Acté River to a Chak'an Itza town called Saklemakal. They arrived at the western end of Lake Petén Itzá to an enthusiastic welcome by the local Itza. The following day, the current Aj Kan Ek' travelled across the lake with eighty canoes to greet the visitors at the Chak'an Itza port town of Nich, on the west shore of Lake Petén Itza. The Franciscans returned to Nojpetén with Kan Ek' and baptised over 300 Itza children over the following four days. Avendaño tried to convince Kan Ek' to convert to Christianity and surrender to the Spanish Crown, without success. The king of the Itza, cited Itza prophecy and said the time was not yet right. On 19 January AjKowoj, the king of the Kowoj, arrived at Nojpetén and spoke with Avendaño, arguing against the acceptance of Christianity and Spanish rule. The discussions between Avendaño, Kan Ek' and AjKowoj exposed deep divisions among the Itza. Kan Ek' learnt of a plot by the Kowoj and their allies to ambush and kill the Franciscans, and the Itza king advised them to return to Mérida via Tipuj. The Spanish friars became lost and suffered great hardships, including the death of one of Avendaño's companions, but after a month wandering in the forest found their way back to Chuntuki, and from there returned to Mérida. Battle at Ch'ich', 2 February 1696 By mid-January Captain García de Paredes had arrived at the advance portion of the Camino Real at Chuntuki. By now he only had 90 soldiers plus labourers and porters. Captain Pedro de Zubiaur, García’s senior officer, arrived at Lake Petén Itza with 60 musketeers, two Franciscans, and allied Yucatec Maya warriors. They were also accompanied by about 40 Maya porters. They were approached by about 300 canoes carrying approximately 2,000 Itza warriors. The warriors began to mingle freely with the Spanish party and a scuffle then broke out; a dozen of the Spanish party were forced into canoes, and three of them were killed. At this point the Spanish soldiers opened fire with their muskets, and the Itza retreated across the lake with their prisoners, who included the two Franciscans. The Spanish party retreated from the lake shore and regrouped on open ground where they were surrounded by thousands of Itza warriors. Zubiaur ordered his men to fire a volley that killed between 30 and 40 Itzas. Realising that they were hopelessly outnumbered, the Spanish retreated towards Chuntuki, abandoning their captured companions to their fate. Martín de Ursúa was now convinced that Kan Ek' would not surrender peacefully, and he began to organise an all-out assault on Nojpetén. Work on the road was redoubled and about a month after the battle at Ch'ich' the Spanish arrived at the lakeshore, now supported by artillery. Again a large number of canoes gathered, and the nervous Spanish soldiers opened fire with cannons and muskets; no casualties were reported among the Itza, who retreated and raised a white flag from a safe distance. Expedition from Verapaz, February – March 1696 Oidor Bartolomé de Amésqueta led the next Guatemalan expedition against the Itza. He marched his men from Cahabón to Mopán, arriving on 25 February 1696. On 7 March, Captain Díaz de Velasco led a party ahead to the lake; he was accompanied by two Dominican friars and by AjK'ixaw, an Itza nobleman who had been taken prisoner on Díaz's previous expedition. When they drew close to the shore of Lake Petén Itzá, AjK'ixaw was sent ahead as an emissary to Nojpetén. Díaz's party was lured into an Itza trap and the expedition members were killed to a man. The two friars were captured and sacrificed. The Itza killed a total of 87 expedition members, including 50 soldiers, two Dominicans and about 35 Maya helpers. Amésqueta left Mopán three days after Díaz and followed Díaz’s trail to the lakeshore. He arrived at the lake over a week later with 36 men. As they scouted along the south shore near Nojpetén they were shadowed by about 30 Itza canoes and more Itzas approached by land but kept a safe distance. Amésqueta was extremely suspicious of the small canoes being offered by the Itza to transport his party across to Nojpetén; as nightfall approached Amésqueta retreated from the lakeshore and his men took up positions on a small hill nearby. In the early hours of the morning he ordered a retreat by moonlight. At San Pedro Martír he received news of an Itza embassy to Mérida in December 1695, and an apparent formal surrender of the Itza to Spanish authority. Unable to reconcile the news with the loss of his men, and with appalling conditions in San Pedro Mártir, Amésqueta abandoned his unfinished fort and retreated to Guatemala. Assault on Nojpetén The Itzas' continued resistance had become a major embarrassment for the Spanish colonial authorities, and soldiers were despatched from Campeche to take Nojpetén once and for all. Martín de Urzúa y Arizmendi arrived on the western shore of Lake Petén Itzá with his soldiers on 26 February 1697, and once there built the heavily armed galeota attack boat. The galeota carried 114 men and at least five artillery pieces. The piragua longboat used to cross the San Pedro River was also transported to the lake to be used in the attack on the Itza capital. On 10 March a number of Itza and Yalain emissaries arrived at Ch'ich' to negotiate with Ursúa. Kan Ek' then sent a canoe with a white flag raised bearing emissaries, who offered peaceful surrender. Ursúa received the embassy in peace and invited Kan Ek' to visit his encampment three days later. On the appointed day Kan Ek' failed to arrive; instead Maya warriors amassed both along the shore and in canoes upon the lake. A waterbourne assault was launched upon Kan Ek's capital on the morning of 13 March. Ursúa boarded the galeota with 108 soldiers, two secular priests, five personal servants, the baptised Itza emissary AjChan and his brother-in-law and an Itza prisoner from Nojpetén. The attack boat was rowed east towards the Itza capital; half way across the lake it encountered a large fleet of canoes spread in an arc across the approach to Nojpetén – Ursúa simply gave the order to row through them. A large quantity of defenders had gathered along the shore of Nojpetén and on the roofs of the city. Itza archers began to shoot at the invaders from the canoes. Ursúa ordered his men not to return fire but arrows wounded a number of his soldiers; one of the wounded soldiers discharged his musket and at that point the officers lost control of their men. The defending Itza soon fled from the withering Spanish gunfire. The city fell after a brief but bloody battle in which many Itza warriors died; the Spanish suffered only minor casualties. The Spanish bombardment caused heavy loss of life on the island; the surviving Itza abandoned their capital and swam across to the mainland with many dying in the water. After the battle the surviving defenders melted away into the forests, leaving the Spanish to occupy an abandoned Maya town. Martín de Ursúa planted his standard upon the highest point of the island and renamed Nojpetén as Nuestra Señora de los Remedios y San Pablo, Laguna del Itza ("Our Lady of Remedy and Saint Paul, Lake of the Itza"). The Itza nobility fled, dispersing to Maya settlements throughout Petén; in response the Spanish scoured the region with search parties. Kan Ek' was soon captured with help from the Yalain Maya ruler Chamach Xulu; The Kowoj king (Aj Kowoj) was also soon captured, together with other Maya nobles and their families. With the defeat of the Itza, the last independent and unconquered native kingdom in the Americas fell to the European colonisers. - Belma has been tentatively identified with the modern settlement and Maya archaeological site of El Meco. - Quezada 2011, p. 13. - Quezada 2011, p. 14. - White and Hood 2004, p. 152. Quezada 2011, p. 14. - Thompson 1966, p. 25. - Quezada 2011, p. 15. - Quezada 2011, pp. 14–15. - Lovell 2005, p. 17. - Sharer and Traxler 2006, pp. 46–47. - Rice and Rice 2009, p. 5. - Quezada 2011, p. 16. - Quezada 2011, p. 17. - White and Hood 2004, p. 152. - Schwartz 1990, p. 17. - Schwartz 1990, p. 18. - Estrada-Belli 2011, p. 52. - Coe 1999, p. 31. Webster 2002, p. 45. - Andrews 1984, p. 589. - Sharer and Traxler 2006, pp. 499–500. - Sharer and Traxler 2006, pp. 613, 616. - Andrews 1984, p. 590. - Caso Barrera 2002, p. 17. - Andrews 1984, p. 591. - Andrews 1984, p. 593. - Andrews 1984, p. 592. - Jones 2000, p. 353. - Houwald 1984, p. 257. - Jones 2000, p. 351. - Jones 2000, p. 352. - Rice and Rice 2009, p. 10. Rice 2009, p. 17. - Cecil et al 1999, p. 788. - Rice and Rice 2005, p. 149. - Rice 2009, p. 17. Feldman 2000, p. xxi. - Smith 2003, p. 279. - Thompson 1966, p. 24. - Thompson 1966, p. 26. - Jones 2000, p. 364. - Rice 2009, p. 83. - Pugh 2009, p. 191. Houwald 1984, p. 256. - Houwald 1984, p. 256. - Clendinnen 2003, p. 7. - Pohl and Hook 2008, p. 26. - Pohl and Hook 2008, pp. 26–27. - Pohl and Hook 2008, p. 27. - Wise and McBride 2008, pp. 33–34. - Wise and McBride 2008, p. 34. - Clendinnen 2003, p. 3. - Perramon 1986, p. 242. Clendinnen 2003, p. 3. - Clendinnen 2003, pp. 3–4. - Sharer and Traxler 2006, p. 758. - Clendinnen 2003, p. 4. - de Díos González 2008, p. 25. Gómez Martín June 2013, p. 56. - Gómez Martín June 2013, p. 56. - de Díos González 2008, pp. 25–26. - de Díos González 2008, p. 26. - Clendinnen 2003, pp. 4–5. - Clendinnen 2003, pp. 6. - Clendinnen 2003, pp. 5. - Clendinnen 2003, p. 8. - Clendinnen 2003, pp. 8–9. - Clendinnen 2003, p. 9. - Clendinnen 2003, pp. 9–10. - Clendinnen 2003, p. 10. - Clendinnen 2003, pp. 10–11. - Clendinnen 2003, p. 11. - Clendinnen 2003, p. 12. - Clendinnen 2003, pp. 11–12. - Clendinnen 2003, pp. 12–13. - Clendinnen 2003, p. 13. - Clendinnen 2003, p. 14. - Sharer and Traxler 2006, p. 759. Clendinnen 2003, p. 14. - Sharer and Traxler 2006, p. 759. Recinos 1986, p. 18. - Recinos 1986, p. 18. - Clendinnen 2003, p. 15. - Clendinnen 2003, pp. 14–15. - Clendinnen 2003, pp. 15–16. - Clendinnen 2003, p. 16. - Sharer and Traxler 2006, pp. 760–761. - Sharer and Traxler 2006, pp. 758–759, 760–761. - Townsend 1995, p. 16. - Hernández et al 2010, p. 26. - Townsend 1995, pp. 16ff. - Jones 2000, p. 358. - Rice and Rice 2009, p. 12. - Rice et al 2009, p. 127. - Rice and Rice 2005, p. 152. - Sharer and Traxler 2006, p. 762. Jones 2000, p. 358. - Sharer and Traxler 2006, p. 773. Jones 2000, p. 358. - Feldman 1998, p. 6. - Webster 2002, p. 83. - Sharer and Traxler 2006, pp. 766–767. - Sharer and Traxler 2006, p. 767. Clendinnen 2003, p. 20. - Clendinnen 2003, p. 20. - ITMB 2000. - Clendinnen 2003, p. 21. - Sharer and Traxler 2006, p. 767. Clendinnen 1989, 2003, p. 21. - Sharer and Traxler 2006, p. 767. - Sharer and Traxler 2006, p. 767-768. - Sharer and Traxler 2006, p. 768. Clendinnen 2003, p. 21. - Quezada 2011, p. 37. - Quezada 2011, pp. 37–38. - Clendinnen 2003, p. 23. Sharer and Traxler 2006, p. 768. - Sharer and Traxler 2006, pp. 768–769. - Sharer and Traxler 2006, pp. 769–770. - Sharer and Traxler 2006, pp. 770–771. - Caso Barrera 2002, pp. 17, 19. - Caso Barrera 2002, p. 19. - Fialko Coxemans 2003, pp. 72–73. - Sharer and Traxler 2006, p. 774. Jones 1998, p. 46. Chuchiak IV 2005, p. 131. - Jones 1998, pp. 42, 47. - Chuchiak IV 2005, p. 132. - Means 1917, p. 79. - Means 1917, p. 80. - Means 1917, p. 81. - Sharer and Traxler 2006, p. 774. Means 1917, p. 81. - Means 1917, p. 81. Jones 1998, pp. 47–48. - Sharer and Traxler 2006, p. 774. Jones 1998, p. 48. - Jones 1998, p. 48. - Jones 1998, pp. 48–49. - Feldman 2000, p. 151. - Jones 1998, pp. 111, 132–133, 145. - Jones 1998, pp. 129–130. - Jones 1998, pp. 130–131. - Jones 1998, p. 131. - Jones 1998, pp. 132, 134. Means 1917, p. 97. - Jones 1998, pp. 135–136, 139–140. - Jones 1998, p. 141. - Jones 1998, p. 140. - Jones 1998, p. 142. - Jones 1998, p. 143. - Jones 1998, pp. 130, 144. - Jones 1998, pp. 148–149. - Jones 1998, p. 147. - Jones 1998, p. 154. Means 1917, pp. 117–118. - Jones 1998, p. 154. - Jones 1998, p. 163. - Jones 1998, p. 162. - Jones 1998, pp. 148, 150. - Jones 1998, pp. 130, 151–152. - Jones 1998, p. 152. - Jones 1998, pp. 150, 154. - Jones 1998, pp. 154–155. - Jones 1998, p. 155. - Jones 1998, p. 156. - Jones 1998, pp. 148, 157. Quezada 2011, p. 23. - Jones 1998, p. 157. - Jones 1998, p. 148. - Jones 1998, p. 158. - Jones 1998, pp. 158–159. - Jones 1998, pp. 159–160. - Jones 1998, p. 160. - Jones 1998, pp. 160–161. - Jones 1998, pp. 187, 189. - Jones 1998, pp. 189–190. Means 1917, p. 128. - Jones 1998, p. 190. - Sharer and Traxler 2006, p. 775. Jones 1998, p. 192. - Jones 1998, p. 205. - Jones 1998, p. 207. - Jones 1998, pp. 209–210. - Sharer and Traxler 2006, p. 775. Jones 1998, pp. 214–215. - Vayhinger-Scheer 2011, p. 383. - Sharer and Traxler 2006, pp. 775–776. Jones 1998, pp. 218–219. - Jones 1998, pp. 189, 226. - Jones 1998, p. 226. - Jones 1998, p. 227. Sharer and Traxler 2006, p. 776. - Jones 1998, p. 227. - Jones 1998, p. 228. Sharer and Traxler 2006, p. 776. - Jones 1998, p. 228. - Jones 1998, p. 229. - Jones 1998, pp. 232–233. - Jones 1998, p. 233. - Jones 1998, pp. 233–234. - Jones 1998, p. 479n59. - Jones 1998, p. 234-235. - Jones 1998, pp. 237–238. - Jones 1998, pp. 238–239. - Jones 1998, p. 240. - Jones 1998, pp. 241–242. - Jones 2000, p. 362. - Jones 2009, p. 59. Jones 1998, pp. 253, 265–266. - Jones 1998, pp. 268–269. - Jones 1998, pp. 252, 268. - Jones 1998, pp. 269–270. - Sharer and Traxler 2006, p. 777. Jones 1998, p. 295. - Jones 1998, p. 297. - Jones 1998, pp. 298–299. - Jones 2009, p. 59. - Sharer and Traxler 2006, pp. 777–778. - Sharer and Traxler 2006, p. 778. Jones 2009, p. 59. - Jones 1998, p. 295. - Jones 1998, p. 306. - Jones 1998, p. xix. |Wikisource has original text related to this article:| - Andrews, Anthony P. (Winter 1984). "The Political Geography of the Sixteenth Century Yucatan Maya: Comments and Revisions". Journal of Anthropological Research (Albuquerque, New Mexico, US: University of New Mexico) 40 (4): 589–596. JSTOR 3629799. Retrieved 2013-12-19. (subscription required) - Athena Review (1999a). "The Spanish Conquest of Yucatán (1526–46)". Athena Review 2 (1). Retrieved 2006-07-25. - Athena Review (1999b). "The Valdivia Shipwreck (1511)". Athena Review 2 (1). Retrieved 2006-07-25. - Caso Barrera, Laura (2002). Caminos en la selva: migración, comercio y resistencia: Mayas yucatecos e itzaes, siglos XVII–XIX [Roads in the Forest: Migration, Commerce and Resistance: Yucatec and Itza Maya, 17th–19th Centuries] (in Spanish). Mexico City, Mexico: El Colegio de México, Fondo de Cultura Económica. ISBN 978-968-16-6714-6. 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Skip to main content Get your Wikispaces Classroom now: the easiest way to manage your class. Pages and Files Anticipatory Guide for Reading Instruction Assessment and Evaluation of ELLs Assessment and Placement Procedures for Incoming ELLs Choral Reading to Promote Oral Language Development Developing Literacy in the Classroom and Home Environments ESL Laws, Legislation and Legal Decisions ESL Professional Development Homepage Evaluating Materials for ESL Reading and Content Instruction Idividualized Education Plan Input, Output and Comprehension Integrated View of SLA Krashen and the Natural Approach Language and Language Learning Language and Linguistics - Resources for ESL Teachers and Students Homepage Language Learning Resources for Students Add "All Pages" Input, Output and Comprehension Input, Output and Comprehension: Input: What is input in the ESL classroom? Input is the information that students receive from each other and the teacher. It is very important in second language acquisition and must be “comprehensible, developmentally appropriate, redundant, and accurate” (Kagan, 1995). So, students must be able to understand the basic message of the information they are receiving. “It is especially critical for them to receive comprehensible input from their teachers and classmates” (Haynes, 2005). The input must be received from many different sources for the information to move from “short-term comprehension to long-term acquisition” (Kagan, 1995). Also, the messages must be “slightly above their current English level” (Haynes, 2005). Krashen, working off of Vygotsky’s theory of the zone of proximal development, “devised a similar notion for the kind of input that an ESL student needs in order to make progress in acquiring English. He called this gap is the current level of proficiency” (ESL Workshop: Scaffolding Theory). Output: What does "output mean and how is it important to the language acquisition process? “Output takes an active role in the acquisition process in that output serves as an avenue for experimenting with language and for testing language hypotheses” (Module 6 Lecture Notes). It can be used as “feedback for intake because it assists the second language learner to check and correct language use,” or it “forces the second language learner to examine a syntactic rather than semantic analysis and examination of language” (Module 2 Lecture Notes). “Language acquisition is fostered by output that is functional and communicative,” as well as comprehensible (Kagan, 1995). For the output to be comprehensible, it must be ordered, structured, coherent and understandable. So, the learner must “draw on L2 syntax” (Module 6 Lecture Notes). It is important to note that “students, to a large extent, learn to speak by speaking” (Kagan, 1995). To become fluent they must be given multiple opportunities to speak about and discuss the same topic, express their own ideas and practice speaking on their own level and terms (Kagan, 1995; Excerpt from Teaching English Language Learners with Learning Difficulties; Haynes, 2005). “There are four ways in which output provides learners with opportunities for authentic language learning: testing hypotheses, receiving feedback, developing automaticity and shifting from meaning-based processing” (Module 6 Lecture Notes). The Role of Interaction: What is the role of interaction and how does it relate to language use and acquisition? Interaction provides ample opportunity for plenty of input and output. “Teachers must constantly involve students, ask many questions, encourage students to express their ideas and thoughts in the new language” (Excerpt from Teaching English Language Learners with Learning Difficulties, 2010). Small groups also help facilitate interaction and adaptation to different levels of proficiency. “You cannot learn a language without interaction” (Foppoli, 2008). Of course, there is the fear of miscommunication. However, miscommunication can encourage “conversational negotiation as well as negotiation of meaning – both beneficial to the process of language acquisition” (Module 6 Lecture Notes). Interaction enables the learner in developing an understanding of grammar and grammatical structures and developing knowledge of L2 syntax (Module 6 Lecture Notes). Comprehension: How do input, output and interaction encourage and promote comprehension? Comprehension is the result of language acquisition. The best way for a second language learner to acquire a new language is through receiving lots of input and having the opportunity to produce a lot of output. The best way for a student to have ample opportunities for input and output is through interaction. To help make input comprehensible to ELL students, make sure to give visual clues along with auditory and involve them in hands-on activities. To encourage output it is important that the students talk more than the teacher. Discussions and small group activities are great ways to make sure the students have ample opportunity for output. To foster interaction the teacher must make sure to have a lot of activities for the students to participate in and to limit the amount of individual and silent work they must participate in. To aid in comprehension it is important that the teacher make sure instructions and information are given clearly and multiple times. Repetition can be very helpful in comprehension. These tips and many more can be accessed by clicking the following link: Quick Tips to Help English Language Learner (ELL) Students in the Classroom Fopolli, J. (2008). The Role of Interaction to Acquire a Second Language. Retrieved from Free Language Lessons Board: Haynes, J. (2005). Comprehensible Input and Output. Retrieved from everythingESL.net: Kagan, S. (1995). We Can Talk: Cooperative Learning in the Elementary ESL Classroom. Retrieved from CAL Digests: What is Comprehensible Input? Teaching English-Language Learners with Learning Difficulties) . (2010). Retrieved from ESL Workshops: Scaffolding Theory. (n.d.). Retrieved from Ms. Cristina Hudgins Back to Professional Development Homepage help on how to format text Turn off "Getting Started"

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Watching this resources will notify you when proposed changes or new versions are created so you can keep track of improvements that have been made. Favoriting this resource allows you to save it in the “My Resources” tab of your account. There, you can easily access this resource later when you’re ready to customize it or assign it to your students. The femur, or thigh bone, is the most proximal (closest to the center of the body) bone of the leg in tetrapod vertebrates capable of walking or jumping, such as most land mammals, birds, many reptiles such as lizards, and amphibians such as frogs . In vertebrates with four legs such as dogs and horses, the femur is found only in the rear legs. The head of the femur articulates with the acetabulum in the pelvic bone to form the hip joint, while the distal part of the femur articulates with the tibia and patella to form the knee joint. By most measures, the femur is the strongest bone in the body. The femur is the longest bone of the human skeleton and is located between the hip bone and the knee. It is the only bone in the thigh. This bone is also one of the strongest bones in the human skeleton. It functions in supporting the weight of the body and allowing motion of the lower extremity. The head (at the proximal extremity) of the femur articulates with the acetabulum of the pelvis to form the hip joint . The lower extremity of the femur (or distal extremity), which is larger, is somewhat cuboid in form and consists of two oblong eminences known as the condyles. The articular surface of the lower end of the femur occupies the anterior, inferior, and posterior surfaces of the condyles. The front or anterior portion is the patellar surface and articulates with the patella. The lower and posterior parts articulate with the corresponding condyles of the tibia to form the knee joint. the femur head articulates with the acetabulum to form the hip joint, the femur is the sole bone in the leg, the femur is the longest bone in the body, or the distal femur articulates with the proximal tibia to form the knee joint

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|Location||Colony of Vancouver Island| |Parties||First Nations of Vancouver Island and the Colony of Vancouver Island| The Douglas Treaties, also known as the Vancouver Island Treaties or the Fort Victoria Treaties, were a series of treaties signed between certain indigenous groups on Vancouver Island and the Colony of Vancouver Island. With the signing of the Oregon Treaty in 1846, the Hudson's Bay Company (HBC) determined that its trapping rights in the Oregon Territory were tenuous. Thus in 1849, it moved its western headquarters from Fort Vancouver on the Columbia River (present day Vancouver, Washington) to Fort Victoria. Fort Vancouver's Chief Factor, James Douglas, was relocated to the young trading post to oversee the Company's operations west of the Rockies. This development prompted the British colonial office to designate the territory a crown colony on January 13, 1849. The new colony, Colony of Vancouver Island, was immediately leased to the HBC for a ten-year period, and Douglas was charged with encouraging British settlement. Richard Blanshard was named the colony's governor. Blanshard discovered that the hold of the HBC over the affairs of the new colony was all but absolute, and that it was Douglas who held all practical authority in the territory. There was no civil service, no police, no militia, and virtually every British colonist was an employee of the HBC. As the colony expanded the HBC started buying up lands for colonial settlement and industry from aboriginal peoples on Vancouver Island. For four years the governor, James Douglas, made a series of fourteen land purchases from aboriginal peoples. To negotiate the terms, Douglas met first in April 1850 with leaders of the Songhees nation, and made verbal agreements. Each leader made an X at the bottom of a blank ledger. The actual terms of the treaty were only incorporated in August, and modelled on the New Zealand Company's deeds of purchase for Maori land, used after the signing of Treaty of Waitangi. The Douglas Treaties cover approximately 930 square kilometres (360 sq mi) of land around Victoria, Saanich, Sooke, Nanaimo and Port Hardy, all on Vancouver Island that were exchanged for cash, clothing and blankets. They were able to retain existing village lands and fields for their use, and also were allowed to hunt and fish on the surrendered lands. These fourteen land purchases became the fourteen Treaties that make up the Douglas Treaties. Douglas didn't continue buying land due to lack of money and the slow growth of the Vancouver Island colony. |Treaty Group Name||Modern First Nation (band government)||Land covered by Treaty||Money exchanged for land||Ref| |Teechamitsa||Esquimalt First Nation||Country lying between Esquimalt and Point Albert||£27 10 shillings (UK £2,626 in 2016)||| |Kosampson||Esquimalt First Nation||Esquimalt Peninsula and Colquitz Valley||£52 10 shillings (UK £5,014 in 2016)||| |Whyomilth||Esquimalt First Nation||Northwest of Esquimalt Harbour||£30 (UK £2,865 in 2016)||| |Chewhaytsum||Becher Bay Band||Sooke||£45 ten shillings (UK £4,345 in 2016)||| |Chilcowitch||Songhees First Nation||Point Gonzales||£45 (UK £4,298 in 2016)||| |Che-ko-nein||Songhees First Nation||Point Gonzales to Cedar Hill||£79 10 shillings (UK £7,593 in 2016)||| |Sooke||T'sou-ke Nation||North-west of Sooke Inlet||£48 6 shillings 8 pence (UK £4,622 in 2016)||| |Ka-ky-aakan||Becher Bay Band||Metchosin||£43 6 shillings 8 pence (UK £4,145 in 2016)||| |Saanich Tribe (South)||Tsawout First Nation and Tsartlip First Nation First Nations||South Saanich||£41 13 shillings 4 pence (UK £3,973 in 2016)||| |Saanich Tribe (North)||Pauquachin First Nation and Tseycum First Nations||North Saanich||[amount not stated]||| |Saalequun||Snuneymuxw First Nation (Former Nanaimo Band)||[area not stated]||[amount not stated]||| |Swengwhung||Songhees First Nation||[area not stated]||[amount not stated]||| |Queackar||Kwakiutl (Kwawkelth) Band||Fort Rupert.||£64 (UK £6,112 in 2016)||| |Quakiolth||Kwakiutl (Kwawkelth) Band||Fort Rupert.||£86 (UK £8,213 in 2016)||| - "Douglas Treaties: 1850-1854". Executive Council of British Columbia. 2009. Retrieved July 28, 2009. - B.C. Archives seeks world heritage status for Douglas treaties, Victoria News, August 08, 2013 8:21 AM - Robin Fisher , 'With or Without Treaty : Indian Land Claims in Western Canada' , in Renwick , ed. . Sovereignty & Indigenous Rights, pp.53 - "1811 - 1867: Pre-Confederation Treaties II". canadiana.org. 2009. Retrieved July 28, 2009. - "Douglas Treaty Payments" (PDF). Executive Council of British Columbia. llbc.leg.bc.ca. 2009. Retrieved July 28, 2009. - British Columbia Indian Treaties In Historical Perspective, Dennis F. K. Madill, Research Branch, Corporate Policy, Department of Indian and Northern Affairs Canada, 1981

https://en.wikipedia.org/wiki/Douglas_Treaties

How the values of A and B affect the shape of the graph y = A sin(Bx). How to graph y = tan(theta) for 0 <= theta < pi/2. How to graph y = tan(q) for one or more periods. How to find the x-intercepts and vertical asymptotes of the graph of y = tan(q). How to recognize when y = 0 is the horizontal asymptote of a rational function. Concept of slope and graphing lines using the slope and y intercept How to graph a quadratic equation by hand. How to graph inequalities in the xy plane. How to write the equation of a graphed line. The description of sex chromosomes. How the value of h affects the shape of the graph y = A sin(B(x-h)). How to recognize the graph of an even or odd function. All about ellipses.

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Bleeding most often occurs due to injury, and depending upon the circumstances, the amount of force required to cause bleeding can be quite variable. Most people understand that falling from a height or being involved in a car accident can inflict great force and trauma upon the body. If blunt force is involved, the outside of the body may not necessarily be damaged, but enough compression may occur to internal organs to cause injury and bleeding. - Imagine a football player being speared by a helmet to the abdomen. The spleen or liver may be compressed by the force and cause bleeding inside the organ. If the hit is hard enough, the capsule or lining of the organ can be torn, and the bleeding can spill into the peritoneum (the space in the abdominal cavity that contains abdominal organs such as the intestines, liver, and spleen). - If the injury occurs in the area of the back or flank, where the kidney is located, retroperitoneal bleeding (retro=behind; behind the abdominal cavity) may occur. - The same mechanism causes bleeding due to crush injuries. For example, when a weight falls on a foot, the weight doesn't give, nor does the ground. The force needs to be absorbed by either the bone or the muscles of the foot. This can cause the bone to break and/or the muscle fibers to tear and bleed. - Other structures are compressible and may cause internal bleeding. For example, the eye can be compressed in the orbit when it is hit by a fist or a ball. The globe deforms and springs back to its original shape. Intraorbital hemorrhage may occur. Deceleration may cause organs in the body to be shifted inside the body. This may shear blood vessels away from the organ and cause bleeding to occur. This is often the mechanism for intracranial bleeding such as epidural or subdural hematomas. Force applied to the head causes an acceleration/deceleration injury to the brain, causing the brain to "bounce around" inside the skull. This can tear some of the small veins on the surface of the brain and cause bleeding. Since the brain is encased in the skull, which is a solid structure, even a small amount of blood can increase pressure inside the skull and decrease brain function. Bleeding may occur with broken bones. Bones contain the bone marrow in which blood production occurs. They have rich blood supplies, and significant amounts of blood can be lost with fractures. The break of a long bone such as the femur (thigh bone) can result in the loss of one unit (350-500cc) of blood. Flat bones such as the pelvis require much more force to cause a fracture, and many blood vessels that surround the structure can be torn by the trauma and cause massive bleeding. Bleeding in pregnancy is never normal, though not uncommon in the first trimester, and is a sign of a potential miscarriage. Early on, the concern is a potential ectopic or tubal pregnancy, in which the placenta and the fetus implant in the Fallopian tube or another location outside of the uterine cavity. As the placenta grows, it erodes through the tube or other involved organs and may cause fatal bleeding. Bleeding after 20 weeks of pregnancy may be due to placenta previa or placental abruption, and urgent medical care should be accessed. Placenta previa describes the situation in which the placenta attaches to the uterus close to the opening of the cervix and may cause painless vaginal bleeding. Abruption occurs when the placenta partially separates from the uterine wall and causes significant pain with or without bleeding from the vagina. Internal bleeding may occur spontaneously, especially in those people who take anticoagulation medications or who have inherited bleeding disorders. Routine bumps that occur in daily life may cause significant bleeding issues. Internal bleeding may be caused as a side effect of medications (most often from nonsteroidal anti-inflammatory drugs such as ibuprofen and aspirin) and alcohol. These substances can cause inflammation and bleeding of the esophagus, stomach, and duodenum, the first part of the small intestine as it leaves the stomach. Long-term alcohol abuse can also cause liver damage, which can cause bleeding problems through a variety of mechanisms. This answer should not be considered medical advice...This answer should not be considered medical advice and should not take the place of a doctor’s visit. Please see the bottom of the page for more information or visit our Terms and Conditions. Archived: March 20, 2014 Thanks for your feedback. 44 of 91 found this helpful Read the Original Article: Internal Bleeding

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3 Answers | Add Yours For historical background, Old English is one of the many precursors to the Modern English language, and was spoken and written between the 5th and 12th centuries C.E. (Wikipedia). It originated with the entrance of Germanic Anglo-Saxons. Latin influence left from the Roman Britain period is not clearly discernible (OED). Old English was a non-standardized collection of regional dialects, so there is no single dictionary for translation as there was no single language. The Old English literary Period started sometime in the 5th century, but there are no surviving documents from that time to serve as examples (runic texts and carvings allow the generalization of the time-frame). The fluxtuating dialect emphases continued throughout the centuries until the 11th century, when it began to change into Middle English based on the London dialect. Middle English held dominance until the standardization of Modern English in the 16th and 17th centuries (the works of Shakespeare and his contemporaries like Spenser and Philips are considered the first properly documented works of Modern English). Therefore, the Old English Period would start sometime in the 5th century and last until the end of the 11th century, when Old English became obsolete. The most famous work written in Old English is the epic poem Beowulf, of unknown author, which is still translated and performed today. The oldest surviving Old English document is Cædmon's Hymn, from the 7th century, which was originally a verbal poem and was never written down by the author. The last surviving document in Old English is a historical record, the Anglo-Saxon Chronicle dated 1154, and shows the beginning influence of Middle English. Middle English was Chaucer's period. The glee—wood rang, a song uprose When Hrothgar’s scop gave the hall good cheer. We’ve answered 300,944 questions. We can answer yours, too.Ask a question

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An oxo-acid is an acid that contains oxygen. To be more specific, it is a compound that contains hydrogen, oxygen, and at least one other element, with at least one hydrogen atom bound to oxygen that can dissociate to produce the H+ cation and the anion of the acid. Under Lavoisier's original theory, all acids contained oxygen, which was named from the Greek ὀξύς (oxys) (acid, sharp) and the root –γενής (–genes) (engender). It was later discovered that some acids, notably hydrochloric acid, did not contain oxygen and so acids were divided into oxoacids and these new hydracids. All oxy-acids have the acidic hydrogen bound to an oxygen atom, so bond strength (length) is not a factor, as it is with binary nonmetal hydrides. Rather, the electronegativity of the central atom (E) and the number of O atoms determine oxy-acid acidity. With the same "central atom" E to which the O is attached, acid strength increases as the number of oxygen attached to E increases. With the same number of oxygens around E, acid strength increases with the electronegativity of E. An oxy-acid molecule contains the structure M-O-H, where other atoms or atom groups can be connected to the central atom M. In a solution, such a molecule can be dissociated to ions in two distinct ways: - M-O-H <=> (M-O)− + H+ - M-O-H <=> (M)+ + OH− If the central atom M is strongly electronegative, then it attracts strongly the electrons of the oxygen atom. In that case, the bond between the oxygen and hydrogen atom is weak, and the compound ionizes easily in the way of the former of the two chemical equations above. In this case, the compound MOH is thus an acid, because it releases a proton, that is, a hydrogen ion. For example, nitrogen, sulfur and chlorine are strongly electronegative elements, and therefore nitric acid, sulfuric acid, and perchloric acid, are strong acids. If, however, the electronegativity of M is weak, then the compound is dissociated to ions according to the latter chemical equation, and MOH is an alkaline hydroxide. Examples of such compounds are sodium hydroxide NaOH and calcium hydroxide Ca(OH)2. If the electronegativity of M is somewhere in between, the compound can even be amphoteric, and in that case, it can dissociate to ions in both ways, in the former case when reacting with bases, and in the latter case when reacting with acids. Inorganic oxy-acids typically have a chemical formula of type HmXOn, where X is some atom functioning as a central atom, whereas parameters m and n depend on the oxidation state of the element X. In most cases, the element X is a nonmetal, but even some metals, for example chromium and manganese, can form oxy-acids when occurring at their highest oxidation state. When oxy-acids are heated, many of them dissociate to water and the anhydride of the acid. In most cases, such anhydrides are oxides of nonmetals. For example, carbon dioxide, CO2, is the anhydride of carbonic acid, H2CO3, and sulfur trioxide, SO3, is the anhydride of sulfuric acid, H2SO4. These anhydrides react quickly with water and form those oxy-acids again. Most acids are oxoacids. Indeed, in the 18th century, Lavoisier assumed that all acids contain oxygen and that oxygen causes their acidity. Because of this, he gave to this element its name, oxygenium, derived from Greek and meaning sharp-maker, which is still, in a more or less modified form, used in most languages. Later, however, Humphry Davy showed that the so-called muriatic acid did not contain oxygen, despite its being a strong acid; instead, it is a solution of hydrogen chloride, HCl. Such acids which do not contain oxygen are nowadays known as hydracids. Names of inorganic oxoacids Many inorganic oxoacids are traditionally called with names ending with the word acid and which also contain, in a somewhat modified form, the name of the element they contain in addition to hydrogen and oxygen. Well-known examples of such acids are sulfuric acid, nitric acid and phosphoric acid. This practice is fully well-established, and even IUPAC has accepted such names. In light of the current chemical nomenclature, this practice is, however, very exceptional, because systematic names of all other compounds are formed only according to what elements they contain and what is their molecular structure, not according to what other properties (for example, acidity) they have. IUPAC, however, does not recommend to call future compounds not yet discovered with a name ending with the word acid. Indeed, acids can even be called with names formed by adding the word hydrogen in front of the corresponding anion; for example, sulfuric acid could just as well be called hydrogen sulfate (or dihydrogen sulfate). In fact, the fully systematic name of sulfuric acid, according to IUPAC's rules, would be dihydroxidodioxidosulfur and that of the sulfate ion, tetraoxidosulfate(2-), Such names, however, are almost never used. However, the same element can form more than one acid when compounded with hydrogen and oxygen. In such cases, the English practice to distinguish such acids is to use the suffix -ic in the name of the element in the name of the acid containing more oxygen atoms, and the suffix -ous in the name of the element in the name of the acid containing fewe oxygen atoms. Thus, for example, sulfuric acid is H2SO4, and sulfurous acid, H2SO3. Analogously, nitric acid is HNO3, and nitrous acid, HNO2. If there are more than two oxoacids having the same element as the central atom, then, in some cases, acids are distinguished by adding the prefix per- or hypo- to their names. The prefix per-, however, is used only when the central atom is a halogen or a group 7 element. For example, chlorine has the four following oxoacids: The suffix -ite occurs in names of anions and salts derived from acids whose names end to the suffix -ous. On the other hand, the suffix -ate occurs in names of anions and salts derived from acids whose names end to the suffix -ic. Prefixes hypo- and per- occur even in name of anions and salts; for example the ion ClO4− is called perchlorate. In a few cases, even prefixes ortho- and para- occur in names of some oxoacids and their derivative anions. In such cases, the para acid is what can be thought as remaining of the ortho acid if a water molecule is separated from the ortho acid molecule. For example, phosphoric acid,H3PO4, has sometimes even be called as orthophosphoric acid, in order to distinguish it from metaphosphoric acid, HPO3. However, according to IUPAC' s current rules, the prefix ortho- should only be used in names of orthotelluric acid and orthoperiodic acid, and their corresponding anions and salts. In the following table, the formula and the name of the anion refer to what remains of the acid when it cedes all hydrogen atoms as protons. Many of these acids, however, are polyprotic, and in such cases, there exists also one or more intermediate anions. In name of such anions, the prefix hydro-, is added if needed, with some numeral prefixes. For example, SO42− is the sulfate anion, and HSO4−, the hydrosulfate anion. In a similar way, PO43− is the phosphate, H2PO42−, the dihydrophosphate, and HPO4−, the hydrophosphate ion. - Kivinen, Antti; Mäkitie, Osmo (1988). Kemia (in Finnish). Helsinki, Finland: Otava. ISBN 951-1-10136-6. - Nomenclature of Inorganic Compounds, IUPAC Recommendations 2005 (Red Book 2005). International Union of Pure and Applied Chemistry. 2005. ISBN 0-85404-438-8.[dead link] - Otavan suuri ensyklopedia, volume 2 (Cid-Harvey) (in Finnish). Helsinki, Finland: Otava. 1977. ISBN 951-1-04170-3. - Kivinen, Mäkitie: Kemia, p. 202-203, chapter=Happihapot - "Hapot". Otavan iso Fokus, Part 2 (El-Io). Otava. 1973. p. 990. ISBN 951-1-00272-4. - Otavan suuri Ensyklopedia, s. 1606, art. Happi - Otavan suuri Ensyklopedia, s. 1605, art. Hapot ja emäxet - Red Book 2005, s. 124, chapter IR-8: Inorganic Acids and Derivatives - Kivinen, Mäkitie: Kemia, p. 459-461, chapter Kemian nimistö: Hapot - Red Book 2005, p. 129-132, table IR-8-1 - Red Book 2005, p. 132, note a

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Electrons and Electric Fields Please edit this page to improve it. See this module's talk page for discussion. Lua error in package.lua at line 80: module 'Module:Yesno' not found. Electrons and Electric Fields Welcome to the lesson. If this is your first exposure to electrical theory, you need some foundational concepts in order to proceed to more difficult lessons. The nuts and bolts of electricity are the electrons, proton-packed atomic nuclei, atoms, molecules, void spaces and the forces that extend across the voids between the solid particles. First, consider those particles. Know subatomic nuclear physics? Electrons have a negative charge. Call it negative 1 electron volt. Neutrons are neutral. Protons have a plus 1 electron volt charge. So, what is a charge? Charge is the potential to attract or repel a charged particle. Positively charged particles are mutually attracted to negatively charged particles. Positively charged particles repel positively charged particles. Negatively charged particles repel negatively charged particles. Unlike charges attract. Like charges repel. It is a quirk of the universe that electrons and protons carry charges of identical magnitude. Use simple addition and subtraction to add or take away the charges of electrons or protons. There are other physical effects that bias this behavior but the description above will suffice in almost every circumstance you will encounter. We have to cover one immediately The nuclear force is short ranged. Something about the order of the diameter of four or five protons is where it begins to drop off. However, it is quite strong enough over its short range to keep the nucleus of an atom together. That is, until you get past Lead on the periodic table of the elements. The reason we have to brush on this topic is so that we can move past the problem of positively charged particles sticking together in the nuclei of our atoms. Take it as given, until the nucleus gets pretty sizeable, it does not want to split. Attraction and repulsion Atoms are formed with a central, positively charged nucleus. This nucleus contains one or many protons and some neutrons. The protons each carry a charge that attracts nearby electrons. An atomic nucleus typically holds electrons in orbit by this positive charge. However, an atom is defined by its nucleus, not by its electrons. So an atom has a cloud of electrons spinning at some distance from the nucleus. In a molecule or compound, the amount of attractive force in one atom may be different from its neighbor. An electron will migrate toward the more attractive atom. The one with the most apparent positive charge. Electrons can change their orbit. With the application of energy, an electron can be induced to speed up, rise to a higher orbit, then stay there until induced to change orbit again. As the electron's orbit gets higher and higher, the positive charge of the nucleus to the orbiting electron is diminished. Taken to its conclusion, the electron can be forced to leave its atom. In conductors, which include most metals, electrons are frequently shared between adjacent atoms. Should one electron be forced off its neighbor, it may in turn displace the neighbor's shared electron, which displaces the next and so on. Long distance relationship Force felt is related to charge by the inverse square of the distance. At a distance of 1 unit, electron A feels 1 unit of repulsion from electron B. At twice that distance, or 2 units, the repulsion would be the inverse of 2 squared or 1/2^2 or 1/4 the force at 1 unit of distance. Electrons orbit atoms, are shared by atoms, and atoms that fly off on their own. Happily, their behavior is predictable. Before an electron can accelerate to a higher orbit, it must be supplied with the force to speed it to the new orbit. In another galactic convenience, it has been demonstrated that electrons that move to a higher shell do so only in fixed units. Heard of quantum mechanics? Quantum physics? Well, a specific, exact, measurable quantum (quanta?) of energy is required to raise the electron exactly one orbital level. No more and no less will suffice. What happens when that orbit is allowed to decay and the electron loses energy? Can you guess? Put down your hand, brainiac. Each orbital unit of decay releases a quantum of energy precisely equal to the one that raised the electron in the first place. Where this becomes important is in things like light bulbs. Remember the loose treatment of electrons by nuclei in metals? Well, when the tungsten filament in a light bulb gets way, way hot, the electrons actually leave and bounce around inside the bulb! When they happen to fall back in toward the filament, they give off those quanta of energy in the form of light and heat. Since we understand the forces involved between charged particles, we now have the mental tools to understand a bit about magnetism. Magnets posses mechanically and chemically fixed electron biases. That is, electrons would appear to be trapped toward one side of the atoms and molecules that make up fixed magnets. This fixing means that one end of the magnet has a relative surplus of negatively charged electrons. The converse holds true, where positively charged atomic nuclei are lain relatively bare. As this is at the molecular level, one can split a fixed magnet and the polarization remains for both fractions of the original magnet. With every electron to the back of its attendant atom, the total charge across the magnet remains neutral. However the near surface is absolutely full of protons (for example) while the rear surface is absolutely full of protons. By our distance squared rule, we can easily see how the near positive charge outweighs the relatively distant negative charge. Splitting a fixed magnet should produce diminishing magnet strength consistent with the decreasing inverse distance squared. [I have not seen this theory presented in any textbook so I must bear the responsibility if this is just plain wrong] There is a way to create a magnetic field artificially. By coiling a conductor as you might coil thread on a spool, then applying an electric current, a magnetic field is induced. There is a left-hand rule for electromagnets. Look at the electromagnet. Put your left thumb in parallel with the core of the coil. If you can close your hand in the direction the electrons flow, then your thumb is pointing toward the North pole of the magnetic field. Otherwise, reverse your hand.. North and South The north pole of a magnet attracts the south pole of another while like poles repel. Magnetic field forces electrons to rotate around magnetic field lines. There are units of measure for magnetic field strength. They are beyond the scope of this lesson. I may have the polar electron attraction/repulsion backwards. Same with the left-hand rule. Experiment among yourselves. Document your results. For extra credit, visit Earth's magnetic north and south poles. Report your observations.

https://en.wikiversity.org/wiki/Electrons_and_Electric_Fields

On the basis of observations of many equilibrium reactions, two Norwegian chemists Goldberg and Waage suggested (1864) a quantitative relationship between the rates of reactions and the concentration of the reacting substances. This relationship is known as law of mass action. It states that “The rate of a chemical reaction is directly proportional to the product of the molar concentrations of the reactants at a constant temperature at any given time.” The molar concentration i.e. number of moles per litre is also called active mass. It is expressed by enclosing the symbols of formulae of the substance in square brackets. For example, molar concentration of A is expressed as [A]. Consider a simple reversible reaction aA + bB ? cC + dD (At a certain temperature) According to law of mass action Rate of forward reaction ∝ [A]a[B]b = kf[A]a[B]b Rate of backward reaction ∝ [C]c[D]d = kb[C]c[D]d Rate of forward reaction = Rate of backward reaction Kf[A]a[B]b = kb[C]c[D]d kf/kb = Kc = [C]c[D]d/[A]a[B]b Where, Kc is called equilibrium constant. In terms of partial pressures, equilibrium constant is denoted by Kp and In terms of mole fraction, equilibrium constant is denoted by and Relation between Kp, Kc and Kx Kp = Kc(RT)Δn Kp = Kk(P)Δn Δn = number of moles of gaseous products – number of moles of gaseous reactants in chemical equation. As a general rule, the concentration of pure solids and pure liquids are not included when writing an equilibrium equation. Characteristics of equilibrium constant (1) The value of equilibrium constant is independent of the original concentration of reactants. (2) The equilibrium constant has a definite value for every reaction at a particular temperature.However, it varies with change in temperature. (3) For a reversible reaction, the equilibrium constant for the forward reaction is inverse of the equilibrium constant for the backward reaction. In general, Kforwardreaction = 1/K'backwardreaction (4) The value of an equilibrium constant tells the extent to which a reaction proceeds in the forward or reverse direction. (5) The equilibrium constant is independent of the presence of catalyst. (6) The value of equilibrium constant changes with the change of temperature. Thermodynamically, it can be shown that if K1 and K2 be the equilibrium constants of a reaction at absolute temperatures T1 and T2. If ?H is the heat of reaction at constant volume, then log K2 – log K1 = –ΔH/2.303R [1/T2 – 1/T1] (Van’t Hoff equation) Transtutors is the best place to get answers to all your doubts regarding law of mass action, equilibrium constant, characteristics of equilibrium constant and Vant Hoff equation. You can submit your school, college or university level homework or assignment to us and we will make sure that you get the answers you need which are timely and also cost effective. Our tutors are available round the clock to help you out in any way with chemistry. Transtutors has a vast panel of experienced chemistry tutors who specialize in law of mass action and can explain the different concepts to you effectively. You can also interact directly with our chemistry tutors for a one to one session and get answers to all your problems in your school, college or university level chemical equilibrium homework. Our tutors will make sure that you achieve the highest grades for your chemistry assignments. We will make sure that you get the best help possible for exams such as the AP, AS, A level, GCSE, IGCSE, IB, Round Square etc. 1 sec ago Science/Math, Math, Calculus, College ask similar question Show transcribed image text Marvin the Miner is quickly walking north on... 58 secs ago Show transcribed image text Marvin the Miner is quickly walking north on 16th street at 5 feet per second while watching Blaster the Burro trotting west on Illinois street a Show transcribed image text Evaluate each of the following without using a... 1 min ago Show transcribed image text Evaluate each of the following without using a calculator. tan^-1 (-3/sqquareroot3 cos^-1(cos 7pi/6) Find the equation of the line tangent to t Show transcribed image text Determine whether each of the statements that... 2 mins ago Show transcribed image text Determine whether each of the statements that follow are TRUE or FALSE. Circle your selection. I Show transcribed image text The figure below show a telephone wire hanging... 3 mins ago Show transcribed image text The figure below show a telephone wire hanging between two poles at x = - b and x = b. It takes the shape of a catenary with equation y = c + a Show transcribed image text Sketch the circle of radius 1 centered at the... 4 mins ago Show transcribed image text Sketch the circle of radius 1 centered at the origin and the circle of radius 1 centered at the point (1,0) both on the same axis. Write an integral which represents the area of the intersection of... Show transcribed image text Write an integral for the arc length of the... 5 mins ago Show transcribed image text Write an integral for the arc length of the curve y = 2/5Square root 25 - x^2 from x = 0 to x = 5. You need not evaluate the integral. Approximate your answer with LEFT(2) and RIGHT(2). Write an... 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Finding the surface area of cones is not that hard. But it can require some patience and ingenuity, depending on what information is available at the beginning of the problem. Below are some suggested steps to keep track of everything. 1Identify the radius of the cone's base circle. If you have the diameter, cut it in half to get the radius. If you have the slant height and perpendicular height, use the Pythagorean theorem (see "Tips" below). 2Write the radius somewhere off to the side, where it's labelled and easy to find, because you will need it several times in several different calculations.It also helps to just find the radius instead of having to look at all your notes and look for the radius. 3Find the area of the base circle by squaring the radius and multiplying by pi. - If the instructions say anything like "exact value", it means that you write the Greek letter for pi and leave it. So a radius of 3 gives an area of 9pi. - Otherwise, use 3.14 or your calculator's pi button to finish the multiplication and get a decimal version for the area. - You can round, but keep at least 3 digits after the decimal point for now. 4Write that answer off to one side, somewhere where it is labelled "base area" and easy to find. 5Identify the slant height of the cone. This refers to the height along the slanted side of the cone, not the height from the tip of the cone to the center of the circle. - The radius, the perpendicular height (from tip to center), and the slant height are related by the Pythagorean theorem. See the "tips" section below. 6Multiply the slant height times the radius times pi. Again, "exact value" means write pi as pi; otherwise, use 3.14 to get the decimal approximation. 7Write that answer off to one side, somewhere where it is labelled "lateral area" and easy to find. 8Add the "base area" from step 4 with the "lateral area" from step 7. 9Round, as needed. This is your final answer. Questions and Answers Give us 3 minutes of knowledge! - The Pythagorean theorem applies to the radius, perpendicular height, and slant height, with the slant height acting as the hypotenuse: (radius)2 + (perpendicular height)2 = (slant height)2 - General rounding rules: any answer under 20 needs at least 2 decimal places. Any answer between 20 and 100 needs only 1 decimal place. Any answer over 100 can be rounded to the nearest whole number. - If either your radius or your slant height has a square root, you will not be able to finish the addition on step 8. Categories: Surface Area In other languages: Español: encontrar el área de la superficie de un cono, Português: Descobrir a Área Superficial de um Cone, Italiano: Calcolare l'Area Totale di un Cono, Deutsch: Die Oberfläche eines Kegels berechnen, Français: calculer la surface d’un cône, 中文: 求出圆锥体的表面积, Русский: найти площадь конуса, Bahasa Indonesia: Menghitung Luas Permukaan Kerucut Thanks to all authors for creating a page that has been read 107,843 times.

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Motion sickness is a sensation of wooziness. It usually occurs when someone is traveling by car, boat, plane, or train. The body's sensory organs send mixed messages to the brain, causing dizziness, lightheadedness, and/or nausea. Some people learn early in their lives that they are prone to the condition. Motion sickness frequently causes vomiting. There are many prevention and treatment measures that can prevent or treat motion sickness. A person maintains balance with the help of signals sent by many parts of the body—for instance, the eyes and inner ears. Other sensory receptors in the legs and feet let the nervous system know what parts of the body are touching the ground. Conflicting signals can cause motion sickness. For example, an airplane traveler cannot see turbulence, but his or her body can feel it. The resulting confusion can cause nausea or even vomiting. Any form of travel, on land, in the air, or on the water, can bring on the uneasy feeling of motion sickness. Sometimes, amusement rides and children’s playground equipment can induce motion sickness. Children between the ages of 3 and 12 are most likely to suffer from motion sickness (Medscape, 2013). Motion sickness usually causes an upset stomach. Other symptoms include a cold sweat and dizziness. A person with motion sickness may become pale or complain of a headache. Motion sickness resolves itself quickly and does not usually require a professional diagnosis. Most people know the feeling when it's coming on because the illness only occurs during travel or other specific activities. Sometimes, pregnant women or people suffering from migraines are misdiagnosed as having motion sickness (Medscape, 2013). Several medications exist for people afflicted with motion sickness. Most prevent only the onset of symptoms. Also, many induce sleepiness, so someone who is operating a vehicle cannot take them. Most people who are susceptible to motion sickness are aware of the fact. If you are prone to motion sickness, the following preventive measure may help: Plan ahead when booking a trip. If traveling by air, ask for a window or wing seat. On trains, sit toward the front. On a ship, ask for a cabin close to the front or the middle of the vessel at water level (Mayo Clinic, 2011). Sitting at the front of a car or bus, or doing the driving yourself, often helps. Many people who experience motion sickness in a vehicle find that they don't have the symptoms when they're driving. It is important to get plenty of rest the night before traveling and avoid drinking alcohol. Dehydration, headache, and anxiety all lead to poorer outcomes for people prone to motion sickness. Eat well so that your stomach is settled. Stay away from greasy or acidic foods during and prior to travel. Have a home remedy on hand or try alternative therapies. Many experts say peppermint can help, as well as ginger and black horehound. Although not very well proved by science, homeopathic remedies do exist. For pilots, astronauts, or others who experience motion sickness regularly or as part of their profession, cognitive therapy and biofeedback have been solutions. Breathing exercises have also been found to help. These treatments also work for people who feel unwell when they merely think about traveling. Written by: David Heitz Published on Dec 18, 2013 Medically reviewed on Dec 18, 2013 by [Ljava.lang.Object;@4b616074

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Native Americans full Moon names were created to help different tribes track the seasons. Think of it as a “nickname” for the Moon! See our list of other full Moon names for each month of the year and their meanings. Why Native Americans Named the Moons The early Native Americans did not record time by using the months of the Julian or Gregorian calendar. Many tribes kept track of time by observing the seasons and lunar months, although there was much variability. For some tribes, the year contained 4 seasons and started at a certain season, such as spring or fall. Others counted 5 seasons to a year. Some tribes defined a year as 12 Moons, while others assigned it 13. Certain tribes that used the lunar calendar added an extra Moon every few years, to keep it in sync with the seasons. Each tribe that did name the full Moons (and/or lunar months) had its own naming preferences. Some would use 12 names for the year while others might use 5, 6, or 7; also, certain names might change the next year. A full Moon name used by one tribe might differ from one used by another tribe for the same time period, or be the same name but represent a different time period. The name itself was often a description relating to a particular activity/event that usually occurred during that time in their location. Colonial Americans adopted some of the Native American full Moon names and applied them to their own calendar system (primarily Julian, and later, Gregorian). Since the Gregorian calendar is the system that many in North America use today, that is how we have presented the list of Moon names, as a frame of reference. The Native American names have been listed by the month in the Gregorian calendar to which they are most closely associated. Native American Full Moon Names and Their Meanings The Full Moon Names we use in the Almanac come from the Algonquin tribes who lived in regions from New England to Lake Superior. They are the names the Colonial Americans adapted most. Note that each full Moon name was applied to the entire lunar month in which it occurred. Link on the names below for your monthly Full Moon Guide! |January||Full Wolf Moon||This full Moon appeared when wolves howled in hunger outside the villages. It is also known as the Old Moon. To some Native American tribes, this was the Snow Moon, but most applied that name to the next full Moon, in February.| |February||Full Snow Moon||Usually the heaviest snows fall in February. Hunting becomes very difficult, and hence to some Native American tribes this was the Hunger Moon.| |March||Full Worm Moon||At the time of this spring Moon, the ground begins to soften and earthworm casts reappear, inviting the return of robins. This is also known as the Sap Moon, as it marks the time when maple sap begins to flow and the annual tapping of maple trees begins.| |April||Full Pink Moon||This full Moon heralded the appearance of the moss pink, or wild ground phlox—one of the first spring flowers. It is also known as the Sprouting Grass Moon, the Egg Moon, and the Fish Moon.| |May||Full Flower Moon||Flowers spring forth in abundance this month. Some Algonquin tribes knew this full Moon as the Corn Planting Moon or the Milk Moon.| |June||Full Strawberry Moon||The Algonquin tribes knew this Moon as a time to gather ripening strawberries. It is also known as the Rose Moon and the Hot Moon.| |July||Full Buck Moon||Bucks begin to grow new antlers at this time. This full Moon was also known as the Thunder Moon, because thunderstorms are so frequent during this month.| |August||Full Sturgeon Moon||Some Native American tribes knew that the sturgeon of the Great Lakes and Lake Champlain were most readily caught during this full Moon. Others called it the Green Corn Moon.| |September||Full Corn Moon||This full Moon corresponds with the time of harvesting corn. It is also called the Barley Moon, because it is the time to harvest and thresh the ripened barley. The Harvest Moon is the full Moon nearest the autumnal equinox, which can occur in September or October and is bright enough to allow finishing all the harvest chores.| |October||Full Hunter’s Moon||This is the month when the leaves are falling and the game is fattened. Now is the time for hunting and laying in a store of provisions for the long winter ahead. October’s Moon is also known as the Travel Moon and the Dying Moon.| |November||Full Beaver Moon||For both the colonists and the Algonquin tribes, this was the time to set beaver traps before the swamps froze, to ensure a supply of warm winter furs. This full Moon was also called the Frost Moon.| |December||Full Cold Moon||This is the month when the winter cold fastens its grip and the nights become long and dark. This full Moon is also called the Long Nights Moon by some Native American tribes.| Note: The Harvest Moon is the full Moon that occurs closest to the autumnal equinox. It can occur in either September or October. At this time, crops such as corn, pumpkins, squash, and wild rice are ready for gathering.

http://www.almanac.com/content/full-moon-names

Universal Declaration of Human Rights |Universal Declaration of Human Rights| |Ratified||16 December 1948| |Location||Palais de Chaillot, Paris| |Rights by claimant| |Other groups of rights| The Universal Declaration of Human Rights (UDHR) is a declaration adopted by the United Nations General Assembly on 10 December 1948 at the Palais de Chaillot, Paris. The Declaration arose directly from the experience of the Second World War and represents the first global expression of rights to which all human beings are inherently entitled. The full text is published by the United Nations on its website. The Declaration consists of thirty articles which have been elaborated in subsequent international treaties, regional human rights instruments, national constitutions, and other laws. The International Bill of Human Rights consists of the Universal Declaration of Human Rights, the International Covenant on Economic, Social and Cultural Rights, and the International Covenant on Civil and Political Rights and its two Optional Protocols. In 1966, the General Assembly adopted the two detailed Covenants, which complete the International Bill of Human Rights. In 1976, after the Covenants had been ratified by a sufficient number of individual nations, the Bill took on the force of international law. - 1 History - 2 Structure - 3 International Human Rights Day - 4 Significance and legal effect - 5 Reaction - 6 Organizations promoting the UDHR - 7 See also - 8 Notes - 9 References - 10 Further reading - 11 External links |Problems playing this file? See media help.| During World War II, the Allies adopted the Four Freedoms—freedom of speech, freedom of religion, freedom from fear, and freedom from want—as their basic war aims. The United Nations Charter "reaffirmed faith in fundamental human rights, and dignity and worth of the human person" and committed all member states to promote "universal respect for, and observance of, human rights and fundamental freedoms for all without distinction as to race, sex, language, or religion". When the atrocities committed by Nazi Germany became apparent after the war, the consensus within the world community was that the United Nations Charter did not sufficiently define the rights to which it referred. A universal declaration that specified the rights of individuals was necessary to give effect to the Charter's provisions on human rights. Creation and drafting The Declaration was commissioned in 1946 and was drafted over two years by the Commission on Human Rights. The Commission consisted of 18 members from various nationalities and political backgrounds. The Universal Declaration of Human Rights Drafting Committee was chaired by Eleanor Roosevelt, who was known for her human rights advocacy. Canadian John Peters Humphrey was called upon by the United Nations Secretary-General to work on the project and became the Declaration's principal drafter. At the time, Humphrey was newly appointed as Director of the Division of Human Rights within the United Nations Secretariat. The Commission on Human Rights, a standing body of the United Nations, was constituted to undertake the work of preparing what was initially conceived as an International Bill of Rights. British representatives were extremely frustrated that the proposal had moral but no legal obligation. (It was not until 1976 that the International Covenant on Civil and Political Rights came into force, giving a legal status to most of the Declaration.) The membership of the Commission was designed to be broadly representative of the global community, served by representatives from the following countries: Australia, Belgium, Byelorussian Soviet Socialist Republic, Chile, Republic of China, Egypt, France, India, Iran, Lebanon, Panama, Philippines, United Kingdom, United States, Union of Soviet Socialist Republics, Uruguay, and Yugoslavia. Well-known members of the Commission included Eleanor Roosevelt of the United States (who was the Chairperson), René Cassin of France, Charles Malik of Lebanon, P. C. Chang of the Republic of China, and Hansa Mehta of India. Humphrey provided the initial draft which became the working text of the Commission. The draft was further discussed by the Commission on human rights, the Economic and Social Council, the Third Committee of the General Assembly before being put to vote. During these discussions many amendments and propositions were made by UN Member States. On 10 December 1948, the Universal Declaration was adopted by the General Assembly by a vote of 48 in favor, none against, and eight abstentions (the Soviet Union, Ukrainian SSR, Byelorussian SSR, People's Federal Republic of Yugoslavia, People's Republic of Poland, Union of South Africa, Czechoslovakia, and the Kingdom of Saudi Arabia). Honduras and Yemen—both members of UN at the time—failed to vote or abstain. South Africa's position can be seen as an attempt to protect its system of apartheid, which clearly violated any number of articles in the Declaration. The Saudi Arabian delegation's abstention was prompted primarily by two of the Declaration's articles: Article 18, which states that everyone has the right "to change his religion or belief"; and Article 16, on equal marriage rights. The six communist nations abstentions centered around the view that the Declaration did not go far enough in condemning fascism and Nazism. Eleanor Roosevelt attributed the abstention of the Soviet bloc nations to Article 13, which provided the right of citizens to leave their countries. The following countries voted in favor of the Declaration: - Costa Rica - Dominican Republic - El Salvador - New Zealand - United Kingdom - United States Despite the central role played by the Canadian John Peters Humphrey, the Canadian Government at first abstained from voting on the Declaration's draft, but later voted in favor of the final draft in the General Assembly. The underlying structure of the Universal Declaration was introduced in its second draft, which was prepared by René Cassin. Cassin worked from a first draft, which was prepared by John Peters Humphrey. The structure was influenced by the Code Napoléon, including a preamble and introductory general principles. Cassin compared the Declaration to the portico of a Greek temple, with a foundation, steps, four columns, and a pediment. Articles 1 and 2 are the foundation blocks, with their principles of dignity, liberty, equality, and brotherhood. The seven paragraphs of the preamble—setting out the reasons for the Declaration—represent the steps. The main body of the Declaration forms the four columns. The first column (articles 3–11) constitutes rights of the individual such as the right to life and the prohibition of slavery. Articles 6 through 11 refer to the fundamental legality of human rights with specific remedies cited for their defense when violated. The second column (articles 12–17) constitutes the rights of the individual in civil and political society (including such things as freedom of movement). The third column (articles 18–21) is concerned with spiritual, public, and political freedoms such as freedom of association, thought, conscience, and religion. The fourth column (articles 22–27) sets out social, economic, and cultural rights. In Cassin's model, the last three articles of the Declaration provide the pediment which binds the structure together. These articles are concerned with the duty of the individual to society and the prohibition of use of rights in contravention of the purposes of the United Nations Organisation. International Human Rights Day The adoption of the Universal Declaration is a significant international commemoration marked each year on 10 December, and is known as Human Rights Day or International Human Rights Day. The commemoration is observed by individuals, community and religious groups, human rights organizations, parliaments, governments, and the United Nations. Decadal commemorations are often accompanied by campaigns to promote awareness of the Declaration and human rights. 2008 marked the 60th anniversary of the Declaration, and was accompanied by year-long activities around the theme "Dignity and justice for all of us". Significance and legal effect The Guinness Book of Records describes the Declaration as the world's "Most Translated Document" (464 different translations). In its preamble, governments commit themselves and their people to progressive measures which secure the universal and effective recognition and observance of the human rights set out in the Declaration. Eleanor Roosevelt supported the adoption of the Declaration as a declaration rather than as a treaty because she believed that it would have the same kind of influence on global society as the United States Declaration of Independence had within the United States. In this, she proved to be correct. Even though it is not legally binding, the Declaration has been adopted in or has influenced most national constitutions since 1948. It has also served as the foundation for a growing number of national laws, international laws, and treaties, as well as for a growing number of regional, sub national, and national institutions protecting and promoting human rights. While not a treaty itself, the Declaration was explicitly adopted for the purpose of defining the meaning of the words "fundamental freedoms" and "human rights" appearing in the United Nations Charter, which is binding on all member states. For this reason, the Universal Declaration is a fundamental constitutive document of the United Nations. In addition, many international lawyers believe that the Declaration forms part of customary international law and is a powerful tool in applying diplomatic and moral pressure to governments that violate any of its articles. The 1968 United Nations International Conference on Human Rights advised that the Declaration "constitutes an obligation for the members of the international community" to all persons. The Declaration has served as the foundation for two binding UN human rights covenants: the International Covenant on Civil and Political Rights and the International Covenant on Economic, Social and Cultural Rights. The principles of the Declaration are elaborated in international treaties such as the International Convention on the Elimination of All Forms of Racial Discrimination, the International Convention on the Elimination of Discrimination Against Women, the United Nations Convention on the Rights of the Child, the United Nations Convention Against Torture, and many more. The Declaration continues to be widely cited by governments, academics, advocates, and constitutional courts, and by individuals who appeal to its principles for the protection of their recognised human rights. The Universal Declaration has received praise from a number of notable people. The Lebanese philosopher and diplomat Charles Malik called it "an international document of the first order of importance", while Eleanor Roosevelt—first chairwoman of the Commission on Human Rights (CHR) that drafted the Declaration—stated that it "may well become the international Magna Carta of all men everywhere." In a speech on 5 October 1995, Pope John Paul II called the Declaration "one of the highest expressions of the human conscience of our time". In a statement on 10 December 2003 on behalf of the European Union, Marcello Spatafora said that the Declaration "placed human rights at the centre of the framework of principles and obligations shaping relations within the international community." However, in 1948, Saudi Arabia abstained from the ratification vote on the Declaration, claiming that it violated Sharia law. Pakistan—which had signed the declaration—disagreed and critiqued the Saudi position. In 1982, the Iranian representative to the United Nations, Said Rajaie-Khorassani, said that the Declaration was "a secular understanding of the Judeo-Christian tradition" which could not be implemented by Muslims without conflict with Sharia. On 30 June 2000, members of the Organisation of the Islamic Conference (now the Organisation of Islamic Cooperation) officially resolved to support the Cairo Declaration on Human Rights in Islam, an alternative document that says people have "freedom and right to a dignified life in accordance with the Islamic Shari'ah", without any discrimination on grounds of "race, colour, language, sex, religious belief, political affiliation, social status or other considerations". Turkey—a secular state with an overwhelmingly Muslim population—signed the Declaration in 1948. A number of scholars in different fields have expressed concerns with the Declaration's alleged Western bias. These include Irene Oh, Abdulaziz Sachedina, Riffat Hassan, and Faisal Kutty. Hassan has argued: What needs to be pointed out to those who uphold the Universal Declaration of Human Rights to be the highest, or sole, model, of a charter of equality and liberty for all human beings, is that given the Western origin and orientation of this Declaration, the "universality" of the assumptions on which it is based is – at the very least – problematic and subject to questioning. Furthermore, the alleged incompatibility between the concept of human rights and religion in general, or particular religions such as Islam, needs to be examined in an unbiased way. Kutty writes: "A strong argument can be made that the current formulation of international human rights constitutes a cultural structure in which western society finds itself easily at home ... It is important to acknowledge and appreciate that other societies may have equally valid alternative conceptions of human rights." On the other hand, others[who?] have written that some of these "cultural arguments" can go so far as to undermine the very nature of human freedom and choice, the protection of which is the purpose of the UN declaration. For example, typical versions of Sharia law forbid Muslims from leaving Islam under the penalty of capital punishment. Islamic legal scholar Faisal Kutty argues that existing blasphemy laws in Muslim countries are actually un-Islamic and are a legacy of colonial rule. Mohsen Haredy, an Islamic scholar, states that Muslim countries have their own views of Sharia and blasphemies are the internal issues of those countries. Ironically, a number of Islamic countries that as of 2014[update] are among the most resistant to UN intervention in domestic affairs, played an invaluable role in the creation of the Declaration, with countries such as Syria and Egypt having been strong proponents of the universality of human rights and the right of countries to self-determination. "The Right to Refuse to Kill" Groups such as Amnesty International and War Resisters International have advocated for "The Right to Refuse to Kill" to be added to the Universal Declaration. War Resisters International has stated that the right to conscientious objection to military service is primarily derived from—but not yet explicit in—Article 18 of the UDHR: the right to freedom of thought, conscience, and religion. Steps have been taken within the United Nations to make this right more explicit, but —to date (2015)—[update] those steps have been limited to less significant United Nations documents. Sean MacBride—Assistant Secretary-General of the United Nations and Nobel Peace Prize laureate—has said: "To the rights enshrined in the Universal Declaration of Human Rights one more might, with relevance, be added. It is 'The Right to Refuse to Kill'." American Anthropological Association The American Anthropological Association criticized the UDHR while it was in its drafting process. The AAA warned that the document would be defining universal rights from a Western paradigm which would be unfair to countries outside of that scope. They further argued that the West's history of colonialism and Evangelicalism made them a problematic moral representative for the rest of the world. They proposed three notes for consideration with underlying themes of cultural relativism: "1. The individual realizes his personality through his culture, hence respect for individual differences entails a respect for cultural differences", "2. Respect for differences between cultures is validated by the scientific fact that no technique of qualitatively evaluating cultures has been discovered.", and "3. Standards and values are relative to the culture from which they derive so that any attempt to formulate postulates that grow out of the beliefs or moral codes of one culture must to that extent detract from the applicability of any Declaration of Human Rights to mankind as a whole." During the lead up to the World Conference on Human Rights held in 1993, ministers from Asian states adopted the Bangkok Declaration, reaffirming their governments' commitment to the principles of the United Nations Charter and the Universal Declaration of Human Rights. They stated their view of the interdependence and indivisibility of human rights and stressed the need for universality, objectivity, and non-selectivity of human rights. However, at the same time, they emphasized the principles of sovereignty and non-interference, calling for greater emphasis on economic, social, and cultural rights—in particular, the right to economic development over civil and political rights. The Bangkok Declaration is considered to be a landmark expression of the Asian values perspective, which offers an extended critique of human rights universalism. Organizations promoting the UDHR International Federation for Human Rights The International Federation for Human Rights (FIDH) is nonpartisan, nonsectarian, and independent of any government, and its core mandate is to promote respect for all the rights set out in the Universal Declaration of Human Rights, the International Covenant on Civil and Political Rights, and the International Covenant on Economic, Social and Cultural Rights. In 1988, director Stephen R. Johnson and 41 international animators, musicians, and producers created a 20-minute video for Amnesty International to celebrate the 40th Anniversary of the Universal Declaration. The video was to bring to life the Declaration's 30 articles. Amnesty International celebrated Human Rights Day and the 60th anniversary of the Universal Declaration all over the world by organizing the "Fire Up!" event. Unitarian Universalist Service Committee The Unitarian Universalist Service Committee (UUSC) is a non-profit, nonsectarian organization whose work around the world is guided by the values of Unitarian Universalism and the Universal Declaration of Human Rights. It works to provide disaster relief and promote human rights and social justice around the world. Quaker United Nations Office and American Friends Service Committee The Quaker United Nations Office and the American Friends Service Committee work on many human rights issues, including improving education on the Universal Declaration of Human Rights. They have developed a Curriculum to help introduce High School students to the Universal Declaration of Human Rights. American Library Association In 1997, the council of the American Library Association (ALA) endorsed Article 19 from the Universal Declaration of Human Rights. Along with Article 19, Article 18 and 20 are also fundamentally tied to the ALA Universal Right to Free Expression and the Library Bill of Rights. Censorship, the invasion of privacy, and interference of opinions are human rights violations according to the ALA. - Human rights - Non-binding agreements - Cairo Declaration on Human Rights in Islam (1990) - Vienna Declaration and Programme of Action (1993) - United Nations Millennium Declaration (2000) - International human rights law - Fourth Geneva Convention (1949) - European Convention on Human Rights (1952) - Convention Relating to the Status of Refugees (1954) - Convention on the Elimination of All Forms of Racial Discrimination (1969) - International Covenant on Civil and Political Rights (1976) - International Covenant on Economic, Social and Cultural Rights (1976) - Convention on the Elimination of All Forms of Discrimination Against Women (1981) - Convention on the Rights of the Child (1990) - Charter of Fundamental Rights of the European Union (2000) - Convention on the Rights of Persons with Disabilities (2007) - Thinkers influencing the Declaration - Charles Malik - Jacques Maritain - John Peters Humphrey - Tommy Douglas - John Sankey, 1st Viscount Sankey - Wu Teh Yao - Peng Chun Chang - Slavery in international law - Slave Trade Acts - Human rights in China (PRC) - Command responsibility - Declaration on Great Apes, an as-yet unsuccessful effort to extend some human rights to great apes. - "Consent of the governed" - Racial equality proposal (1919) - The Farewell Sermon (632 CE) - Youth for Human Rights International - "The Universal Declaration of Human Rights". un.org. - Williams 1981. This is the first book edition of the Universal Declaration of Human Rights, with a foreword by Jimmy Carter. - "United Nations Charter, preamble and article 55". United Nations. Retrieved 2013-04-20. - Cataclysm and World Response in Drafting and Adoption : The Universal Declaration of Human Rights, udhr.org. - "UDHR50: Didn't Nazi tyranny end all hope for protecting human rights in the modern world?". Udhr.org. 1998-08-28. Retrieved 2012-07-07. - "UDHR – History of human rights". Universalrights.net. Retrieved 2012-07-07. - Morsink 1999, p. 5 - Morsink 1999, p. 133 - Morsink 1999, p. 4 - Universal Declaration of Human Rights. Final authorized text. The British Library. September 1952. Retrieved 16 August 2015. - The Declaration was drafted during the Chinese Civil War. P.C. Chang was appointed as a representative by the Republic of China, then the recognised government of China, but which was driven from mainland China and now administers only Taiwan and nearby islands (). - "Drafting of the Universal Declaration of Human Rights". Research Guides. United Nations. Dag Hammarskjöld Library. Retrieved 2015-04-17. - Carlson, Allan: Globalizing Family Values, 12 January 2004. - CCNMTL. "default". Center for New Media Teaching and Learning (CCNMTL). Columbia University. Retrieved 2013-07-12. - UNAC. "Questions and answers about the Universal Declaration of Human Rights". United Nations Association in Canada (UNAC). p. "Who are the signatories of the Declaration?". Archived from the original on 2012-09-12. External link in - Jost Müller-Neuhof (2008-12-10). "Menschenrechte: Die mächtigste Idee der Welt". Der Tagesspiegel (in German). Retrieved 2013-07-12. - Peter Danchin. "The Universal Declaration of Human Rights: Drafting History - 10. Plenary Session of the Third General Assembly Session". Retrieved 2015-02-25. - Glendon 2002, pp. 169–70 - "Yearbook of the United Nations 1948–1949 p 535" (PDF). Retrieved 24 July 2014. - Schabas, William (1998). "Canada and the Adoption of Universal Declaration of Human Rights" (PDF). McGill Law Journal 43: 403. - Glendon 2002, pp. 62–64. - Glendon 2002, Chapter 10. - "The Universal Declaration of Human Rights: 1948–2008". United Nations. Retrieved 15 February 2011. - "Universal Declaration of Human Rights". United Nations Office of the High Commissioner for Human rights. - Humphrey JP, "The Universal Declaration of Human Rights: Its History, Impact and Juridical Character", in Ramcharan BG (ed), Human Rights: Thirty Years After the Universal Declaration (1979) pp. 2l, 37; Sohn 1, "The Human Rights Law of the Charter" (1977) 12 Texas Int LJ 129, 133; McDougal MS, Lasswell H and Chen I, Human Rights and World Public Order (1980) pp. 273–274, 325–327; D'Amato A, International Law: Process and Prospect( 1986) pp. 123–147. - Office of the High Commissioner for Human Rights. "Digital record of the UDHR". United Nations. - "Statement by Charles Malik as Representative of Lebanon to the Third Committee of the UN General Assembly on the Universal Declaration". 6 November 1948. Archived from the original on 28 September 2008. - Michael E. Eidenmuller (1948-12-09). "Eleanor Roosevelt: Address to the United Nations General Assembly". Americanrhetoric.com. Retrieved 2012-07-07. - "John Paul II, Address to the U.N., October 2, 1979 and October 5, 1995". Vatican.va. Retrieved 2012-07-07. - Nisrine Abiad (2008). Sharia, Muslim states and international human rights treaty obligations: a comparative study. BIICL. pp. 60–65. ISBN 978-1-905221-41-7. - Price 1999, p. 163 - Littman, D (February–March 1999). "Universal Human Rights and Human Rights in Islam". Midstream. Archived from the original on 2006-05-12. - "Resolution No 60/27-P". Organisation of the Islamic Conference. 2000-06-27. Retrieved 2011-06-02. - "Universal Declaration of Human Rights". Retrieved 2015-10-30. - "Are Human Rights Compatible with Islam?". religiousconsultation.org. Retrieved 2012-11-12. - "The Rights of God". Georgetown University Press, 2007. - "Non-Western Societies Have Influenced Human Rights". in Jacqueline Langwith (ed.), Opposing Viewpoints: Human Rights, Gale/Greenhaven Press: Chicago, 2007. - "Why Blasphemy Laws Are Actually Anti-Islamic". The Huffington Post. - "Why Should Blasphemy Be Punishable at All?". onislam.net. - Professor Susan Waltz: Universal Rights Group, Syria calls for greater UN intervention in domestic human rights situations…. - Out of the margins: the right to conscientious objection to military service in Europe: An announcement of Amnesty International's forthcoming campaign and briefing for the UN Commission on Human Rights, 31 March 1997. Amnesty International. - A Conscientious Objector's Guide to the UN Human Rights System, Parts 1, 2 & 3, Background Information on International Law for COs, Standards which recognise the right to conscientious objection, War Resisters' International. - Sean MacBride, The Imperatives of Survival, Nobel Lecture, 12 December 1974, The Nobel Foundation – Official website of the Nobel Foundation. (English index page; hyperlink to Swedish site.) From Nobel Lectures in Peace 1971–1980. - "Statement on Human Rights" (PDF). Retrieved 2015-10-30. - "Final Declaration Of The Regional Meeting For Asia Of The World Conference On Human Rights". Law.hku.hk. Retrieved 2012-07-07. - Contribution to the EU Multi-stakeholder Forum on CSR (Corporate Social Responsibility), 10 February 2009; accessed on 9 November 2009 - Information Partners, web site of the UNHCR, last updated 25 February 2010, 16:08 GMT (web retrieval 25 February 2010, 18:11 GMT) - "UDHR film". Amnesty International. Retrieved 2013-07-19. - "Fire Up!". Amnesty International. Retrieved 2013-07-19. - "UNHCR Partners". UNHCR. Retrieved 11 November 2014. - "AFSC Universal Declaration of Human Rights web page". American Friends Service Committee. Retrieved 11 November 2014. - "Resolution on IFLA, Human Rights and Freedom of Expression". ala.org. - "The Universal Right to Free Expression:". ala.org. - Glendon, Mary Ann (2002). A world made new: Eleanor Roosevelt and the Universal Declaration of Human Rights. Random House. ISBN 978-0-375-76046-4. - Hashmi, Sohail H. (2002). Islamic political ethics: civil society, pluralism, and conflict. Princeton University Press. ISBN 978-0-691-11310-4. - Morsink, Johannes (1999). The Universal Declaration of Human Rights: origins, drafting, and intent. University of Pennsylvania Press. ISBN 978-0-8122-1747-6. - Price, Daniel E. (1999). Islamic political culture, democracy, and human rights: a comparative study. Greenwood Publishing Group. ISBN 978-0-275-96187-9. - Williams, Paul (1981). The International bill of human rights. United Nations General Assembly. Entwhistle Books. ISBN 978-0-934558-07-5. - Feldman, Jean-Philippe. "Hayek's Critique Of The Universal Declaration Of Human Rights". Journal des Economistes et des Etudes Humaines, Volume 9, Issue 4 (December 1999): 1145-6396. - Nurser, John. "For All Peoples and All Nations. Christian Churches and Human Rights.". (Geneva: WCC Publications, 2005). - Universal Declaration of Human Rights pages at Columbia University (Centre for the Study of Human Rights), including article by article commentary, video interviews, discussion of meaning, drafting and history. - Introductory note by Antônio Augusto Cançado Trindade and procedural history on the Universal Declaration of Human Rights in the Historic Archives of the United Nations Audiovisual Library of International Law ||This article's use of external links may not follow Wikipedia's policies or guidelines. (November 2015)| |Wikisource has original text related to this article:| |Wikimedia Commons has media related to Universal Declaration of Human Rights.| - Text of the UDHR - Official translations of the UDHR - Resource Guide on the Universal Declaration of Human Rights at the UN Library, Geneva. - Drafting of the Universal Declaration of Human Rights - documents and meetings records - United Nations Dag Hammarskjöld Library - Questions and answers about the Universal Declaration - Text, Audio, and Video excerpt of Eleanor Roosevelt's Address to the United Nations on the Universal Declaration of Human Rights - UDHR – Education - Revista Envío – A Declaration of Human Rights For the 21st Century - Introductory note by Antônio Augusto Cançado Trindade and procedural history note on the Universal Declaration of Human Rights in the Historic Archives of the United Nations Audiovisual Library of International Law - DHpedia: Universal Declaration of Human Rights - The Laws of Burgos: 500 Years of Human Rights from the Law Library of Congress blog. - Librivox: Human-read audio recordings in several Languages - Text, Audio, and Video excerpt of Eleanor Roosevelt's Address to the United Nations on the Universal Declaration of Human Rights - Animated presentation of the Universal Declaration of Human Rights by Amnesty International on YouTube (in English duration 20 minutes and 23 seconds). - Audio: Statement by Charles Malik as Representative of Lebanon to the Third Committee of the UN General Assembly on the Universal Declaration, 6 November 1948 - UN Department of Public Information introduction to the drafters of the Declaration - Audiovisual material on the Universal Declaration of Human Rights in the Historic Archives of the United Nations Audiovisual Library of International Law

https://en.wikipedia.org/wiki/Universal_Declaration_of_Human_Rights

Nowadays many pupils, when given a research task, immediately might think to themselves, “I’ll just Google that.” Internet search engines (of which Google is only one of many) are powerful tools but many pupils use only a fraction of the power of them, and then can also have difficulty finding the information specific to the task. There are many resources now available to help in developing pupil skills in searching more effectively using online search engines. And, of course, when they do find information how do pupils know it is appropriate for the task? Or how do they evaluate what is suitable, and how do they present it and show where the information was found. Tools to Help Teach Research Skills The Big 6 One method of teaching information skills for investigating sources of information from databases, encyclopedias and the Internet is that known as “the Big Six.” This process sets out the steps as follows: 1. Define the task – what needs to be done? 2. Information Seeking Strategies – what resources can I use? 3. Location and Access – where can I find these resources? 4. Use of information – what can I use from these resources? 5. Synthesis – what can I make to finish the job? 6. Evaluation – how will I know I did my job well? The Kentucky Virtual Library How to Do Research The Kentucky Virtual Library has an online poster-style How to Do Research site for guiding younger pupils through the steps to finding the information they need on any topic, whether in print form, multimedia or online. Presented in a visual comic/game style it explains in child-friendly language the process to find the information being sought. And each page of advice is presented as a set of easy to digest straightforward steps, breaking down each task (whether finding the information, recording it, evaluating it, or presenting it) in cartoon-style visual interactive style making it attractive to primary users. Finding Dulcinea – How to Search the Internet – aimed at older pupils, this provides a host of helpful tips and links to a variety of resources about searching and using information from the Internet. It includes sections on What Is the Internet, Web Site Credibility, How Search Engines Work, Choosing a Search Engine, Online Databases, Social Bookmarking Tools, How to Cite a Source. Ergo – Teaching Research Skills Ergo – Teaching Research Skills from the State Library of Victoria, Australia, is a guide for pupils to finding the information they need for a school assignment. The guide provides helpful explanations, hints, tips and further resources for each of the steps: Define the task, Locate information, Select resources, Organise notes, Present the ideas, Evaluate your work. Common Sense Media Digital Curriculum Common Sense Media Digital Curriculum has a section on teaching online research for various age groups. Each section has lesson plans with ideas and resources for teaching different aspects of research online with pupils. Teachers can select resources according to age group or stage (all stages in primary school are included, and resources are age-appropriate), resources to support the research topics which best suit the needs of pupils. All About Explorers All About Explorers has been designed as an interactive Internet search task to guide pupils through making more discerning use of information presented online. The task includes a range of spoof material to help show primary pupils how to evaluate what they read online, and how to be selective about the information they find. The tasks are presented to pupils as an interactive Webquest. There is a section for teachers which includes a series of lessons and explanations of what the pupils are learning about better online searching as they complete each webquest. Save the Tree Octopus – an example of a spoof website which could be used to show pupils that, even though the site looks very well put together and with a host of features to make it look authoritative, websites can provide completely fictitious information. Ten Tips for Teaching How to Research and Filter Information Ten Tips for Teaching Students How to Research and Filter Information – a post by Kathleen Morris which details advice for ten steps for showing primary school pupils how to find and use information: Search, Delve, Source, Validity, Purpose, Background, Teach, Justify, Path, Cite. Google Tools for Better Searching Google has produced a series of posters for educators to help support pupils use the Internet search engine more effectively. Google A Day is a daily-changing search challenge which could be used by a class to make better use of a search engine. Each day a new challenge is presented (and you can go back to previous challenges if you wish). Each challenge is presneted as a question which pupils are challenged to answer by using the search engine. If not sure how to get started pupils can click on the hint to get a bit of help to guide how to make a better search to find the answer. And the answer itself is provided. In addition there are links to tips and techniques for better Internet searching. Google for Educators is a collection of resources collated by David Andrade on his Educational Technology Guy blog. This brings together a series of resources providing tips, ideas and guides to how the vast array of Google tools can be used in schools, including how to find what you’re wanting using the search engine. Interesting Ways to Use Google Search in the Classroom is a collection of ideas collected by Tom Barrett shared by many teachers – like others in the “Interesting ideas” series it grows as more teachers contribute ideas. So if you have a way you have used Google Search in your classroom then you too can add yours there too. Google Guide is an online interactive tutorial and reference for experienced users, novices, and everyone in between. Nancy Blachman developed Google Guide to provide more information about Google’s capabilities, features, and services. There are hints and ideas, a printable sheet of tips, and interactive exercises teachers can use with pupils to guide them them through making use of different techniques for more effective searching for information. Entire Guide to Google Search Features for teachers and Students by Mohamed Kharbach details steps, tips and tools to make better use of the Google search engine, from the basics to advanced searching to using a variety of features of the search engine in many different contexts. Google Search Education Evangelism is a site with lessons to download for free, including Powerpoint presentations and guides for printing about making the best of Google search tools. These are arranged in categories and for different audiences, whether teachers self-study or for use with pupils. Google Search Education has lesson plans and video tutorials in categories of various search skills in using Google. Within each category there are then tutorials presneted to suit different skill levels. 10 Google Search Tips by Catlin Tucker provides 10 Questions & 10 Answers to Help You or your pupils Search Smarter! 12 Ways to use Google Search by Degree of Difficulty is a series of lessons by Jeff Dunn providing graded techniques for being better at using the search facility with Google – for each step there are three levels of difficulty so you just choose which best suits your need for your class. PhD in Googling! An Animated presentation on tips to using Google search engine. Thanks to David Andrade for sharing this. PhD in Googling presents a series of graphically interesting screens with nugget-sized tips on each page, and with animated text appearing, explaining the tip. Update Your Search methods – a blog post by Chris Betcher explaining how in 2013 Google changed the way the seacrh engibe works to better interpret plain English questions, the way someone would ask a question if speaking, rather than relying on keywords. Get More Out of Google is a poster with advice and practical tips for making more eficient use of Google search engine.

https://blogs.glowscotland.org.uk/fa/ICTFalkirkPrimaries/tag/google/

Threatened Species of Shark Bay Shark Bay World Heritage Area is a refuge for some of the world’s most endangered animals and plants. Its isolated islands and peninsulas have been largely spared the feral predators and habitat destruction that wreaked havoc on mainland Australia. The importance of these habitats in protecting vulnerable wildlife, and providing scientific information on the impact of habitat change, was a major factor in Shark Bay being declared a World Heritage site. Two of Shark Bay’s islands, Bernier and Dorre Islands, are the last stronghold for five critically endangered land mammals – four of which occur in the wild nowhere else on Earth. An ecological restoration project is underway on the largest of Shark Bay's islands, Dirk Hartog Island. This project aims to restore habitats, remove feral cats and reintroduce ten species lost from the island during its pastoral era. It will also introduce two native species not previously known to occur on the island. In other parts of Shark Bay, cats, foxes and grazing stock have been removed in order to allow the ecosystem to rejuvenate. Captive-bred animals have been introduced to places such as Francois Peron National Park as part of Project Eden, a local conservation initiative. The Australian Wildlife Conservancy has established a wildlife sanctuary on Faure Island. After removing cats and goats they successfully introduced a number of species. Find out more here. Shark bay mouse (Pseudomys fieldi) Spiny-tailed skink (Egernia stokesii) Small dragon orchid Caladenia barbarella Shark Bay also features many endemic plants, including two threatened species. For more information about Western Australian wildlife, check out the WA Museum Fauna Base website. Learn more about Western Australia’s plants at the West Australian Herbarium’s FloraBase website.

http://www.sharkbay.org.au/nature-of-shark-bay-threatened-species.aspx

May 31, 2002 Antarctica is home to more than 70 lakes that lie thousands of meters under the surface of the continental ice sheet, including one under the South Pole itself. Lake Vostok, beneath Russia's Vostok Station, is one of the largest of these subglacial lakes, comparable in size and depth to Lake Ontario, one of the North American Great Lakes. There is some evidence that Vostok's waters may contain microbial life. Exploration of the lake to confirm that life exists will be an international effort and will require the development of ultra-clean technologies to prevent contaminating the waters. The National Science Foundation, as manager of the U.S. Antarctic Program, coordinates nearly all U.S. research in Antarctica and would lead U.S. participation in any international effort to explore the lake. NSF's Office of Polar Programs has established a steering committee to study the possible scientific exploration of Antarctic subglacial lakes. See: http://www.nsf.gov/od/opp/antarct/subglclk.jsp Vostok Station is located in one of the world's most inaccessible places, near the South Geomagnetic Pole, at the center of the East Antarctic Ice Sheet. The station is 3.5 kilometers (11,484 feet) above sea level. The coldest temperature ever recorded on Earth, -89.2 degree Celsius (-128.6 degrees Fahrenheit), was measured at Vostok Station on July 21, 1983. Lake Vostok's physical characteristics have led scientists to argue that it might serve as an earthbound analog for Europa, a moon of Jupiter. Confirming that life can survive in Lake Vostok might strengthen the argument for the presence of life on Europa. Russian and British scientists confirmed the lake's existence in 1996 by integrating a variety of data, including airborne ice-penetrating radar observations and spaceborne radar altimetry. Researchers working at Vostok Station have already contributed greatly to climatology by producing one of the world's longest ice cores in 1998. A joint Russian, French and U.S, team drilled and analyzed the core, which is 3,623 meters (11,886 feet) long. The core contains layers of ice deposited over millennia, representing a record of Earth's climate stretching back more than 420,000 years. Drilling of the core was deliberately halted roughly 150 meters (492 feet) above the suspected boundary where the ice sheet and the liquid waters of the lake are thought to meet to prevent contamination of the lake. It is from samples of this ice core, specifically from ice that is thought to have formed from lake water freezing onto the base of the ice sheet, that NSF-funded scientists believe they have found evidence that the lake water supports life. Their research was published in Science in 1999. For more information, see: http://www.nsf.gov/od/lpa/news/press/99/pr9972.htm More recently, NSF-funded researchers from the Lamont-Doherty Earth Observatory at Columbia University, using data gathered by the University of Texas Institute for Geophysics, published a paper in Nature suggesting that the hydrodynamics of Lake Vostok may make it possible to search for evidence of life in the layers of ice that accumulate on the lake's eastern shore. Scientists say such a possibility would provide another avenue for exploring the lake's potential as a harbor for microscopic life, in addition to exploring the lake itself. For more information, see: http://www.nsf.gov/od/lpa/news/02/pr0219.htm International consensus building Discussions at an NSF workshop for U.S. researchers held in 1998, a subsequent international meeting held in Cambridge, England in 1999, as well as other international meetings about subglacial lakes, have formed the basis for a developing scientific consensus on whether, and how, to proceed with exploring the lake's waters. To read a report from the 1998 NSF workshop "Lake Vostok: A Curiosity or a Focus for Interdisciplinary Study?" see: http://www.ldeo.columbia.edu/vostok/ The Scientific Committee on Antarctic Research (SCAR) hosts a Web site on subglacial lake exploration with links to several reports from various international workshops. See: http://salegos-scar.montana.edu/ Peter West, NSF, (703) 292-8070, [email protected] The National Science Foundation (NSF) is an independent federal agency that supports fundamental research and education across all fields of science and engineering. In fiscal year (FY) 2016, its budget is $7.5 billion. NSF funds reach all 50 states through grants to nearly 2,000 colleges, universities and other institutions. Each year, NSF receives more than 48,000 competitive proposals for funding and makes about 12,000 new funding awards. NSF also awards about $626 million in professional and service contracts yearly. Get News Updates by Email Useful NSF Web Sites: NSF Home Page: http://www.nsf.gov NSF News: http://www.nsf.gov/news/ For the News Media: http://www.nsf.gov/news/newsroom.jsp Science and Engineering Statistics: http://www.nsf.gov/statistics/ Awards Searches: http://www.nsf.gov/awardsearch/

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