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June 7, 2013 Underwater Earthquake Sounds May Help Predict Tsunamis Brett Smith for redOrbit.com - Your Universe Online While seismologists have the ability to detect undersea earthquakes, they don´t know which ones will turn into tsunamis. Some ocean-based sensors are able to detect an approaching tsunami, but they can only provide a few minutes of advanced warning.According to a new report in The“¯Bulletin of the Seismological Society of America, two Stanford University researchers have identified key acoustic characteristics of the 2011 Japan earthquake that“¯could be used to significantly improve tsunami warning systems. The researchers stumbled onto the acoustic signature somewhat accidentally while working with advanced computer models designed to simulate the catastrophic quake. While investigating the quake, the seismologists were puzzled over the reason behind why the earthquake rupture propagated from the fault deep below the ocean all the way up to the seafloor, creating a massive upward thrust and a resulting tsunami. The two geophysicists“¯Eric Dunham and Jeremy Kozdon began using the powerful supercomputers at Stanford's“¯Center for Computational Earth and Environmental Science (CEES) to model the tremors as they moved through the Earth´s crust and through the ocean. Although direct observations of the earthquake were scare, the researchers created a high-resolution model that used known geologic features of the Japan Trench and CEES simulations to identify previous earthquake events that matched the available data. The Stanford models were able to accurately recreate the seafloor uplift in the 2011 earthquake, which caused the gigantic tsunami wave heights. The model also revealed sound waves that propagated within the ocean during the event. The researchers found that surface-breaking ruptures, like the 2011 earthquake, create higher amplitude ocean acoustic waves than those that do not. These sound waves would have traveled through the water and reached shore 15 to 20 minutes before the tsunami, according to the report. "We've found that there's a strong correlation between the amplitude of the sound waves and the tsunami wave heights," Dunham said. "Sound waves propagate through water 10 times faster than the tsunami waves, so we can have knowledge of what's happening a hundred miles offshore within minutes of an earthquake occurring. We could know whether a tsunami is coming, how large it will be and when it will arrive." According to the researchers, their model could be used on fault zones around the world that have the potential to generate a massive tsunami, including faults near Japan, the Pacific Northwest and Chile. "The ideal situation would be to analyze lots of measurements from major events and eventually be able to say, 'this is the signal'," Kozdon said. "Fortunately, these catastrophic earthquakes don't happen frequently, but we can input these site specific characteristics into computer models — such as those made possible with the CEES cluster — in the hopes of identifying acoustic signatures that indicates whether or not an earthquake has generated a large tsunami." The researchers said underwater microphones would still need to be deployed on the seafloor or using buoys to establish an early warning system. Any acoustic signals would then need to be analyzed to confirm a threat. They expect policymakers would coordinate with scientists to determine a course of action in the event that a tsunami is detected.

http://www.redorbit.com/news/science/1112868200/earthquake-sounds-help-predict-tsunamis-060713/

If you like us, please share us on social media, tell your friends, tell your professor or consider building or adopting a Wikitext for your course. Some of the most important carbon-carbon bond-forming and bond-breaking processes in biological chemistry involve the gain or loss, by an organic molecule, of a single carbon in the form of CO2. You undoubtedly have seen this chemical equation before in an introductory biology or chemistry class: 6CO2 + 6H2O + energy → C6H12O6 + 6O2 This of course represents the photosynthetic process, by which plants (and some bacteria) harness energy from sunlight to build glucose from individual carbon dioxide molecules. The key chemical step in which carbon dioxide is 'fixed' (in other words, condensed with an existing organic molecule) is called a carboxylation reaction. It is catalyzed by the enzyme ribulose 1,5-bisphosphate carboxylase, commonly known as Rubisco, in the 'Calvin cycle' of carbon fixation. The reverse chemical equation is also probably familiar to you: C6H12O6 + 6O2 → 6CO2 + 6H2O + energy This equation expresses what happens in respiration: the oxidative breakdown of glucose to form carbon dioxide, water, and energy. In the course of this transformation, each of the carbon atoms of glucose is eventually converted to individual CO2 molecules. The actual chemical step by which a carbon atom, in the form of carbon dioxide, breaks off from a larger organic molecule is called a decarboxylation reaction. The key decarboxylation steps in the conversion of glucose to carbon dioxide occur in the citric acid (Krebs) cycle and the pentose phosphate pathway. Let's now look at the organic mechanisms of some carboxylation and decarboxylation reactions. Carboxylation reactions are essentially just aldol condensations, except that the carbonyl electrophile is CO2 rather than a ketone or aldehyde. The mechanism for Rubisco, the key carbon-fixing enzyme in plants and photosynthetic bacteria (and the most abundant enzyme on earth!), is shown below. Magnesium ion plays a key charge-stabilizing role throughout the reaction. Step 1, not surprisingly, is deprotonation of an alpha-carbon to form an enolate. Step 2 is simply an intramolecular proton transfer, which has the effect of creating a different enolate intermediate and making C2 into the nucleophile for an aldol-like attack on CO2 (step 3). Carbon dioxide has now been 'fixed' into organic form - it has become a carboxylate group on a six-carbon sugar derivative. To follow the Rubisco mechanism through to its endpoint: Step 4 is a retro-Claisen mechanism, with a water nucleophile and enolate leaving group. After protonation of this enolate, we are left with two molecules of 3-phosphoglycerate, which are incorporated into the 'gluconeogenesis' pathway of glucose synthesis. Propose a complete mechanism for the retro-Claisen reaction in the figure above. Mechanistically, a decarboxylation has parallels to retro-aldol cleavage reactions: Just as in retro-aldol reactions, the electrons from the broken carbon-carbon bond have to have some place to go - they must, in other words, be stabilized - for the decarboxylation step to take place. Quite often, the electrons are stabilized by the formation of an enolate, as is the case in the general mechanism pictured above. This of course means that a carbonyl group must be positioned beta to (i.e. two carbons down from) the carboxylate carbon. If there is no stable place for the electrons in the carbon-carbon bond to go, then a decarboxylation is very unlikely. Be especially careful, when drawing decarboxylation mechanisms, to resist the temptation to treat the CO2 molecule as the leaving group: This is not what a decarboxylation looks like! In a decarboxylation step, it is the organic part of the molecule that is, in fact, the leaving group, 'pushed off' by the electrons on the carboxylate. Below are two important key b-carbonyl decarboxylation steps in glucose metabolism, each representing a point at which a carbon derived from glucose is released as CO2. Exercise 13.9: Draw mechanistic arrows showing the first step in each of the reactions shown above. Exercise 13.10: Draw a mechanism for the decarboxylation step in the reaction above. An interesting variation on decarboxylation in the synthesis of tyrosine in bacteria is shown below. a) Draw a mechanism for the reaction above. b) Draw hypothetical decarboxylation mechanisms showing the formation of alternate products A and B from the same starting compound: c) How would you describe in words the relationship between the actual product and hypothetical products A and B? Which of the three would you expect to be most thermodynamically stable, and why? Recall from section 6.5B that many enzymes are dependent upon the assistance of coenzymes, which are small (relative to protein) organic molecules that bind - covalently or non-covalently - in an enzyme's active site and help it to catalyze its reaction. S-adenosylmethionine (SAM, section 9.1A) and ATP (section 10.2) are two examples that we have encountered so far, and we will see several more in the chapters ahead. Although Rubisco (described in part B of this section) is an exception, most enzymes that catalyze carboxylation reactions are dependent upon a coenzyme called biotin, which serves as a temporary carrier of carbon dioxide. Pyruvate carboxylase, the enzyme catalyzing the first step of the gluconeogensis pathway, is a good example of a biotin-dependent carboxylation reaction. Notice that the CO2 in this reaction is derived from bicarbonate, unlike the Rubisco reaction in which CO2 is 'fixed' directly from the atmosphere. Biotin is covalently attached to the enzyme through an amide linkage to an active site lysine. The exact mechanism by which biotin-dependent carboxylation reactions operate is still not completely understood, however the following is a likely picture. First, the bicarbonate ion is phosphorylated by ATP (step 1, see section 10.2), and thus is activated for decarboxylation, which generates free CO2 (step 2). Biotin's job is to hold on to the carbon dioxide molecule until pyruvate comes into the active site. Carboxylation of biotin involves deprotonation of the amide nitrogen to form an enolate-like intermediate (step 3 - amides have a pKa of approximately 17, and this is lowered by the presence of an active site acid near the oxygen). This step is followed by attack of the nucleophilic nitrogen on carbon dioxide to form carboxybiotinylated enzyme (step 4). When a pyruvate molecule binds, rearrangement of the active site architecture causes the previous step to go in reverse (step 5), freeing the CO2 and generating a biotin base to deprotonate the alpha-carbon of pyruvate so that it can condense, in an aldol-like fashion, with CO2 to form oxaloacetate (steps 6-7). If you have studied some biochemistry, you may have heard about biotin in a somewhat different context that what is discussed in this section. A protein called avidin, found in abundance in egg white, binds non-covalently and extremely tightly to biotin (in fact, avidin-biotin is the tightest protein-ligand binding pair known to science). Biochemists often make use of this property by covalently linking a biomolecule of interest to biotin. The 'biotinylated' species can then be easily isolated from a complex mixture by running the mixture through an 'affinity column' that is coated with avidin.

http://chemwiki.ucdavis.edu/Organic_Chemistry/Organic_Chemistry_With_a_Biological_Emphasis/Chapter_13%3A_Reactions_with_stabilized_carbanion_intermediates_I/Section_13.5%3A_Carboxylation_and_decarboxylation_reactions

Plasma is a form (also called state) of matter. The three other common states of matter are solids, liquids and gases, so plasma is sometimes called the fourth state of matter. Plasma is created by adding energy to a gas so that some of its electrons leave its atoms. This is called ionization. It results in negatively charged electrons, and positively charged ions. Unlike the other states of matter, the charged particles in a plasma will react strongly to electric and magnetic fields (i.e. electromagnetic fields). If a plasma loses heat, the ions will re-form into a gas, emitting the energy which had caused them to ionize. Over 99% of the matter in the visible universe is believed to be plasma. When the atoms in a gas are broken up, the pieces are called electrons and ions. Because they have an electric charge, they are pulled together or pushed apart by electric fields and magnetic fields. This makes a plasma act different from a gas. For example, magnetic fields can be used to hold a plasma, but not to hold a gas. Plasma is a better conductor of electricity than copper. Plasma is usually very hot, because it takes very high temperatures to break the bonds between electrons and the nuclei of the atoms. Sometimes plasmas can have very high pressure, like in stars. Stars (including the Sun) are mostly made of plasma. Plasmas can also have very low pressure, like in outer space. On earth lightning and aurora make plasma. Artificial (man-made) uses of plasma include fluorescent lightbulbs, neon signs, and plasma displays used for television or computer screens, as well as plasma lamps and globes which are a popular children's toy and room decoration. Scientists are experimenting with plasma to make a new kind of nuclear power, called fusion, which would be much better and safer than ordinary nuclear power, and would produce much less radioactive waste. Related pages[change | change source] Other websites[change | change source] - Plasmas: the Fourth State of Matter - Plasma Science and Technology - Plasma on the Internet comprehensive list of plasma related links. - Introduction to Plasma Physics: Graduate course given by Richard Fitzpatrick | M.I.T. Introduction by I.H.Hutchinson - NRL Plasma Formulary online (or an html version) - Plasma Coalition page - Plasma Material Interaction - How to make a glowing ball of plasma in your microwave with a grape | More (Video) - How to make plasma in your microwave with only one match (video) - U.S. Dept of Agriculture research project "Decontamination of Fresh Produce with Cold Plasma" - (French) CNRS LAEPT "Electric Arc Thermal Plasmas - "Phases of Matter". NASA. http://www.grc.nasa.gov/WWW/K-12/airplane/state.html. Retrieved 2011-05-04.

https://simple.wikipedia.org/wiki/Plasma_(physics)

Civil Rights Act of 1964 Teacher Resources Find Civil Rights Act of 1964 educational ideas and activities Showing 81 - 100 of 155 resources Language Recognition and Language Families Learners debate the pros and cons of bilingualism in the United States and in the classroom. Students investigate how language reflects and influences culture, and focus on how to make language acquisition easier for learners. 9th - 12th Social Studies & History Shirley Chisholm: Unbought, Unbossed, and Unforgotten Students examine Chisholm's life and career, and explore her thoughts on some of the issues of her time. They read a speech that Chisholm delivered on the Equal Rights Amendment and conduct research on the Women's Rights Movement of the... 7th - 12th Social Studies & History Remember: The Journey to School Integration Young scholars explore equality by reading a classic Civil Rights story. In this desegregation lesson, students read the book Remember: The Journey to School Integration by Toni Morrison while discussing the Civil Rights Movement. Young... 4th - 9th Visual & Performing Arts Concept Formation/Life in a Box: The Turbulent 60's and 70's Students examine 1960's and 1970's America. In this contemporary history lesson, students examine images, texts, and documents that provide a glimpse into American politics during the decades. Students then conduct further research and... 8th - 11th Social Studies & History "I Am Woman, Hear Me Roar" Students explore U.S. history by examining the role women played in the development of the country. In this women's rights lesson, students read several books with their classmates which discuss the fight women had to go through to get... 5th - 10th Visual & Performing Arts Do the Research! Civil Rights Act Passed in U.S. In this research worksheet, learners find information on the Internet about the Civil Rights Act passed in the United States. Using the information they have found, students answer four general short essay questions. 5th - 6th Social Studies & History Adding to the Picture: The 1963 March on Washington Who do your scholars imagine when they think about the civil rights movement? If only a few faces come to mind, this instructional activity will expand their concepts of the movement's leaders. Learners examine an image of the 1963 March... 7th - 11th Social Studies & History Does It Looks All Right to Me? Learners explore the concept of philanthropy. In this service learning lesson, students examine the accomplishments of Civil Rights leaders' as works of philanthropy. Learners read literature regarding diversity and study the Selma to... 6th - 8th English Language Arts

http://www.lessonplanet.com/lesson-plans/civil-rights-act-of-1964/5

Polymer Exchange Membrane Fuel Cells The polymer exchange membrane fuel cell (PEMFC) is one of the most promising fuel cell technologies. This type of fuel cell will probably end up powering cars, buses and maybe even your house. The PEMFC uses one of the simplest reactions of any fuel cell. First, let's take a look at what's in a PEM fuel cell: In Figure 1 you can see there are four basic elements of a PEMFC: - The anode, the negative post of the fuel cell, has several jobs. It conducts the electrons that are freed from the hydrogen molecules so that they can be used in an external circuit. It has channels etched into it that disperse the hydrogen gas equally over the surface of the catalyst. - The cathode, the positive post of the fuel cell, has channels etched into it that distribute the oxygen to the surface of the catalyst. It also conducts the electrons back from the external circuit to the catalyst, where they can recombine with the hydrogen ions and oxygen to form water. - The electrolyte is the proton exchange membrane. This specially treated material, which looks something like ordinary kitchen plastic wrap, only conducts positively charged ions. The membrane blocks electrons. For a PEMFC, the membrane must be hydrated in order to function and remain stable. - The catalyst is a special material that facilitates the reaction of oxygen and hydrogen. It is usually made of platinum nanoparticles very thinly coated onto carbon paper or cloth. The catalyst is rough and porous so that the maximum surface area of the platinum can be exposed to the hydrogen or oxygen. The platinum-coated side of the catalyst faces the PEM. This content is not compatible on this device. Figure 2. Animation of a working fuel cell Figure 2 shows the pressurized hydrogen gas (H2) entering the fuel cell on the anode side. This gas is forced through the catalyst by the pressure. When an H2 molecule comes in contact with the platinum on the catalyst, it splits into two H+ ions and two electrons (e-). The electrons are conducted through the anode, where they make their way through the external circuit (doing useful work such as turning a motor) and return to the cathode side of the fuel cell. Meanwhile, on the cathode side of the fuel cell, oxygen gas (O2) is being forced through the catalyst, where it forms two oxygen atoms. Each of these atoms has a strong negative charge. This negative charge attracts the two H+ ions through the membrane, where they combine with an oxygen atom and two of the electrons from the external circuit to form a water molecule (H2O). This reaction in a single fuel cell produces only about 0.7 volts. To get this voltage up to a reasonable level, many separate fuel cells must be combined to form a fuel-cell stack. Bipolar plates are used to connect one fuel cell to another and are subjected to both oxidizing and reducing conditions and potentials. A big issue with bipolar plates is stability. Metallic bipolar plates can corrode, and the byproducts of corrosion (iron and chromium ions) can decrease the effectiveness of fuel cell membranes and electrodes. Low-temperature fuel cells use lightweight metals, graphite and carbon/thermoset composites (thermoset is a kind of plastic that remains rigid even when subjected to high temperatures) as bipolar plate material. In the next section, we'll see how efficient fuel-cell vehicles can be.

http://auto.howstuffworks.com/fuel-efficiency/alternative-fuels/fuel-cell2.htm

Human genetic variation Human genetic variation is the genetic differences both within and among populations. There may be multiple variants of any given gene in the human population (genes), leading to polymorphism. Many genes are not polymorphic, meaning that only a single allele is present in the population: the gene is then said to be fixed. On average, in terms of DNA sequence all humans are 99.5% similar to any other humans. No two humans are genetically identical. Even monozygotic twins, who develop from one zygote, have infrequent genetic differences due to mutations occurring during development and gene copy-number variation. Differences between individuals, even closely related individuals, are the key to techniques such as genetic fingerprinting. Alleles occur at different frequencies in different human populations, with populations that are more geographically and ancestrally remote tending to differ more. Causes of differences between individuals include the exchange of genes during meiosis and various mutational events. There are at least two reasons why genetic variation exists between populations. Natural selection may confer an adaptive advantage to individuals in a specific environment if an allele provides a competitive advantage. Alleles under selection are likely to occur only in those geographic regions where they confer an advantage. The second main cause of genetic variation is due to the high degree of neutrality of most mutations. Most mutations do not appear to have any selective effect one way or the other on the organism. The main cause is genetic drift, this is the effect of random changes in the gene pool. In humans, founder effect and past small population size (increasing the likelihood of genetic drift) may have had an important influence in neutral differences between populations. The theory that humans recently migrated out of Africa supports this. The study of human genetic variation has both evolutionary significance and medical applications. It can help scientists understand ancient human population migrations as well as how different human groups are biologically related to one another. For medicine, study of human genetic variation may be important because some disease-causing alleles occur more often in people from specific geographic regions. New findings show that each human has on average 60 new mutations compared to their parents. Apart from mutations, many genes that may have aided humans in ancient times plague humans today. For example, it is suspected that genes that allow humans to more efficiently process food are those that make people susceptible to obesity and diabetes today. - 1 Measures of variation - 2 History and geographic distribution - 3 Categorization of the world population - 4 Health - 5 Genome projects - 6 See also - 7 References - 8 Further reading - 9 External links Measures of variation Nucleotide diversity is the average proportion of nucleotides that differ between two individuals. The human nucleotide diversity is estimated to be 0.1% to 0.4% of base pairs. A difference of 1 in 1,000 amounts to approximately 3 million nucleotide differences, because the human genome has about 3 billion nucleotides. Single nucleotide polymorphisms A single nucleotide polymorphism (SNP) is difference in a single nucleotide between members of one species that occurs in at least 1% of the population. It is estimated that there are 10 to 30 million SNPs in humans. SNPs are the most common type of sequence variation, estimated to account for 90% of all sequence variation. Other sequence variations are single base exchanges, deletions and insertions. SNPs occur on average about every 100 to 300 bases and so are the major source of heterogeneity. A functional, or non-synonymous, SNP is one that affects some factor such as gene splicing or messenger RNA, and so causes a phenotypic difference between members of the species. About 3% to 5% of human SNPs are functional (see International HapMap Project). Neutral, or synonymous SNPs are still useful as genetic markers in genome-wide association studies, because of their sheer number and the stable inheritance over generations. A coding SNP is one that occurs inside a gene. There are 105 Human Reference SNPs that result in premature stop codons in 103 genes. This corresponds to 0.5% of coding SNPs. They occur due to segmental duplication in the genome. These SNPs result in loss of protein, yet all these SNP alleles are common and are not purified in negative selection. Structural variation is the variation in structure of an organism's chromosome. Structural variations, such as copy-number variation and deletions, inversions, insertions and duplications, account for much more human genetic variation than single nucleotide diversity. This was concluded in 2007 from analysis of the diploid full sequences of the genomes of two humans: Craig Venter and James D. Watson. This added to the two haploid sequences which were amalgamations of sequences from many individuals, published by the Human Genome Project and Celera Genomics respectively. Copy number variation A copy-number variation (CNV) is a difference in the genome due to deleting or duplicating large regions of DNA on some chromosome. It is estimated that 0.4% of the genomes of unrelated humans differ with respect to copy number. When copy number variation is included, human-to-human genetic variation is estimated to be at least 0.5% (99.5% similarity). Copy number variations are inherited but can also arise during development. Epigenetic variation is variation in the chemical tags that attach to DNA and affect how genes get read. The tags, "called epigenetic markings, act as switches that control how genes can be read." At some alleles, the epigenetic state of the DNA, and associated phenotype, can be inherited across generations of individuals. Genetic variability is a measure of the tendency of individual genotypes in a population to vary (become different) from one another. Variability is different from genetic diversity, which is the amount of variation seen in a particular population. The variability of a trait is how much that trait tends to vary in response to environmental and genetic influences. In biology, a cline is a continuum of species, populations, races, varieties, or forms of organisms that exhibit gradual phenotypic and/or genetic differences over a geographical area, typically as a result of environmental heterogeneity. In the scientific study of human genetic variation, a gene cline can be rigorously defined and subjected to quantitative metrics. In the study of molecular evolution, a haplogroup is a group of similar haplotypes that share a common ancestor with a single nucleotide polymorphism (SNP) mutation. Haplogroups pertain to deep ancestral origins dating back thousands of years. The most commonly studied human haplogroups are Y-chromosome (Y-DNA) haplogroups and mitochondrial DNA (mtDNA) haplogroups, both of which can be used to define genetic populations. Y-DNA is passed solely along the patrilineal line, from father to son, while mtDNA is passed down the matrilineal line, from mother to both daughter and son. The Y-DNA and mtDNA may change by chance mutation at each generation. Variable number tandem repeats A variable number tandem repeat (VNTR) is the variation of length of a tandem repeat. A tandem repeat is the adjacent repetition of a short nucleotide sequence. Tandem repeats exist on many chromosomes, and their length varies between individuals. Each variant acts as an inherited allele, so they are used for personal or parental identification. Their analysis is useful in genetics and biology research, forensics, and DNA fingerprinting. History and geographic distribution The Out of Africa theory (more precisely called "recent African origin of modern humans") is the most widely accepted explanation of the origin and early dispersal of anatomically modern humans, Homo sapiens sapiens. The theory states that archaic Homo sapiens evolved into modern humans solely in Africa, 200,000 to 100,000 years ago; around that time, one African subpopulation of hominins among several was the subpopulation ancestral to all human beings today. Some members of that subpopulation left Africa by 60,000 years ago and over time replaced earlier hominin populations such as Homo erectus and Neanderthals on Earth. Alternative theories include the multiregional origin of modern humans hypothesis. The theory is supported by both genetic and fossil evidence. The hypothesis originated in the 19th century, with Darwin's Descent of Man, but remained speculative until the 1980s when it was supported by study of present-day mitochondrial DNA, combined with evidence from physical anthropology of archaic specimens. A large study published in 2009 found that modern humans probably originated near the border of Namibia and South Africa (reported as Namibia and Angola by BBC), and some left Africa through East Africa. Observations consistent with this are that Africa contains the most human genetic diversity anywhere on Earth, and the genetic structure of Africans traces to 14 ancestral population clusters that correlate with ethnicity and culture or language. The study lasted ten years and analyzed variations at 1,327 DNA markers of 121 African populations, 4 African American populations, and 60 non-African populations. According to a 2000 study of Y-chromosome sequence variation, human Y-chromosomes trace ancestry to Africa, and the descendants of the derived lineage left Africa and eventually were replaced by archaic human Y-chromosomes in Eurasia. The study also shows that a minority of contemporary East Africans and Khoisan are the descendants of the most ancestral patrilineages of anatomically modern humans that left Africa 35,000 to 89,000 years ago. Other evidence supporting the theory is that variations in skull measurements decrease with distance from Africa at the same rate as the decrease in genetic diversity. Human genetic diversity decreases in native populations with migratory distance from Africa, and this is thought to be due to bottlenecks during human migration, which are events that temporarily reduce population size. In the field of population genetics, it is believed that the distribution of neutral polymorphisms among contemporary humans reflects human demographic history. It has been theorized that humans passed through a population bottleneck before a rapid expansion coinciding with migrations out of Africa leading to an African-Eurasian divergence around 100,000 years ago (ca. 5,000 generations), followed by a European-Asian divergence about 40,000 years ago (ca. 2,000 generations). Richard G. Klein, Nicholas Wade and Spencer Wells, among others, have postulated that modern humans did not leave Africa and successfully colonize the rest of the world until as recently as 60,000 - 50,000 years B.P., pushing back the dates for subsequent population splits as well. The rapid expansion of a previously small population has two important effects on the distribution of genetic variation. First, the so-called founder effect occurs when founder populations bring only a subset of the genetic variation from their ancestral population. Second, as founders become more geographically separated, the probability that two individuals from different founder populations will mate becomes smaller. The effect of this assortative mating is to reduce gene flow between geographical groups, and to increase the genetic distance between groups. The expansion of humans from Africa affected the distribution of genetic variation in two other ways. First, smaller (founder) populations experience greater genetic drift because of increased fluctuations in neutral polymorphisms. Second, new polymorphisms that arose in one group were less likely to be transmitted to other groups as gene flow was restricted. Our history as a species also has left genetic signals in regional populations. For example, in addition to having higher levels of genetic diversity, populations in Africa tend to have lower amounts of linkage disequilibrium than do populations outside Africa, partly because of the larger size of human populations in Africa over the course of human history and partly because the number of modern humans who left Africa to colonize the rest of the world appears to have been relatively low (Gabriel et al. 2002). In contrast, populations that have undergone dramatic size reductions or rapid expansions in the past and populations formed by the mixture of previously separate ancestral groups can have unusually high levels of linkage disequilibrium (Nordborg and Tavare 2002). Many other geographic, climatic, and historical factors have contributed to the patterns of human genetic variation seen in the world today. For example, population processes associated with colonization, periods of geographic isolation, socially reinforced endogamy, and natural selection all have affected allele frequencies in certain populations (Jorde et al. 2000b; Bamshad and Wooding 2003). In general, however, the recency of our common ancestry and continual gene flow among human groups have limited genetic differentiation in our species. Distribution of variation The distribution of genetic variants within and among human populations are impossible to describe succinctly because of the difficulty of defining a "population," the clinal nature of variation, and heterogeneity across the genome (Long and Kittles 2003). In general, however, an average of 85% of genetic variation exists within local populations, ~7% is between local populations within the same continent, and ~8% of variation occurs between large groups living on different continents,. (Lewontin 1972; Jorde et al. 2000a). The recent African origin theory for humans would predict that in Africa there exists a great deal more diversity than elsewhere, and that diversity should decrease the further from Africa a population is sampled. Sub-Saharan Africa has the most human genetic diversity and the same has been shown to hold true for phenotypic diversity. Phenotype is connected to genotype through gene expression. Genetic diversity decreases smoothly with migratory distance from that region, which many scientists believe to be the origin of modern humans, and that decrease is mirrored by a decrease in phenotypic variation. Skull measurements are an example of a physical attribute whose within-population variation decreases with distance from Africa. The distribution of many physical traits resembles the distribution of genetic variation within and between human populations (American Association of Physical Anthropologists 1996; Keita and Kittles 1997). For example, ~90% of the variation in human head shapes occurs within continental groups, and ~10% separates groups, with a greater variability of head shape among individuals with recent African ancestors (Relethford 2002). A prominent exception to the common distribution of physical characteristics within and among groups is skin color. Approximately 10% of the variance in skin color occurs within groups, and ~90% occurs between groups (Relethford 2002). This distribution of skin color and its geographic patterning — with people whose ancestors lived predominantly near the equator having darker skin than those with ancestors who lived predominantly in higher latitudes — indicate that this attribute has been under strong selective pressure. Darker skin appears to be strongly selected for in equatorial regions to prevent sunburn, skin cancer, the photolysis of folate, and damage to sweat glands. Understanding how genetic diversity in the human population impacts various levels of gene expression is an active area of research. While earlier studies focused on the relationship between DNA variation and RNA expression, more recent efforts are characterizing the genetic control of various aspects of gene expression including chromatin states, translation, and protein levels. A study published in 2007 found that 25% of genes showed different levels of gene expression between populations of European and Asian descent. The primary cause of this difference in gene expression was thought to be SNPs in gene regulatory regions of DNA. Another study published in 2007 found that approximately 83% of genes were expressed at different levels among individuals and about 17% between populations of European and African descent. There is a hypothesis that anatomically modern humans interbred with Neanderthals during the Middle Paleolithic. In May 2010, the Neanderthal Genome Project presented genetic evidence that interbreeding did likely take place and that a small but significant[how?] portion of Neanderthal admixture is present in the DNA of modern Eurasians and Oceanians, and nearly absent in sub-Saharan African populations. Between 4% and 6% of the genome of Melanesians (represented by the Papua New Guinean and Bougainville Islander) are thought to derive from Denisova hominins - a previously unknown species which shares a common origin with Neanderthals. It was possibly introduced during the early migration of the ancestors of Melanesians into Southeast Asia. This history of interaction suggests that Denisovans once ranged widely over eastern Asia. Thus, Melanesians emerge as the most archaic-admixed population, having Denisovan/Neanderthal-related admixture of ~8%. In a study published in 2013, Jeffrey Wall from University of California studied whole sequence-genome data and found higher rates of introgression in Asians compared to Europeans. Hammer et al. tested the hypothesis that contemporary African genomes have signatures of gene flow with archaic human ancestors and found evidence of archaic admixture in African genomes, suggesting that modest amounts of gene flow were widespread throughout time and space during the evolution of anatomically modern humans. Categorization of the world population New data on human genetic variation has reignited the debate about a possible biological basis for categorization of humans into races. Most of the controversy surrounds the question of how to interpret the genetic data and whether conclusions based on it are sound. Some researchers argue that self-identified race can be used as an indicator of geographic ancestry for certain health risks and medications. Although the genetic differences among human groups are relatively small, these differences in certain genes such as duffy, ABCC11, SLC24A5, called ancestry-informative markers (AIMs) nevertheless can be used to reliably situate many individuals within broad, geographically based groupings. For example, computer analyses of hundreds of polymorphic loci sampled in globally distributed populations have revealed the existence of genetic clustering that roughly is associated with groups that historically have occupied large continental and subcontinental regions (Rosenberg et al. 2002; Bamshad et al. 2003). Some commentators have argued that these patterns of variation provide a biological justification for the use of traditional racial categories. They argue that the continental clusterings correspond roughly with the division of human beings into sub-Saharan Africans; Europeans, Western Asians, Central Asians, Southern Asians and Northern Africans; Eastern Asians, Southeast Asians, Polynesians and Native Americans; and other inhabitants of Oceania (Melanesians, Micronesians & Australian Aborigines) (Risch et al. 2002). Other observers disagree, saying that the same data undercut traditional notions of racial groups (King and Motulsky 2002; Calafell 2003; Tishkoff and Kidd 2004). They point out, for example, that major populations considered races or subgroups within races do not necessarily form their own clusters. Furthermore, because human genetic variation is clinal, many individuals affiliate with two or more continental groups. Thus, the genetically based "biogeographical ancestry" assigned to any given person generally will be broadly distributed and will be accompanied by sizable uncertainties (Pfaff et al. 2004). In many parts of the world, groups have mixed in such a way that many individuals have relatively recent ancestors from widely separated regions. Although genetic analyses of large numbers of loci can produce estimates of the percentage of a person's ancestors coming from various continental populations (Shriver et al. 2003; Bamshad et al. 2004), these estimates may assume a false distinctiveness of the parental populations, since human groups have exchanged mates from local to continental scales throughout history (Cavalli-Sforza et al. 1994; Hoerder 2002). Even with large numbers of markers, information for estimating admixture proportions of individuals or groups is limited, and estimates typically will have wide confidence intervals (Pfaff et al. 2004). Genetic data can be used to infer population structure and assign individuals to groups that often correspond with their self-identified geographical ancestry. Recently, Lynn Jorde and Steven Wooding argued that "Analysis of many loci now yields reasonably accurate estimates of genetic similarity among individuals, rather than populations. Clustering of individuals is correlated with geographic origin or ancestry." Forensic anthropologists can determine aspects of geographic ancestry (i.e. Asian, African, or European) from skeletal remains with a high degree of accuracy by analyzing skeletal measurements. According to some studies, individual test methods such as mid-facial measurements and femur traits can identify the geographic ancestry and by extension the racial category to which an individual would have been assigned during their lifetime, with over 80% accuracy, and in combination can be even more accurate. However, the skeletons of persons who have recent ancestry in different geographical regions, can exhibit characteristics of more than one ancestral group, and hence cannot be identified as belonging to any single ancestral group. Gene flow and admixture Gene flow between two populations reduces the average genetic distance between the populations, only totally isolated human populations experience no gene flow and most populations have continuous gene flow with other neighboring populations which create the clinal distribution observed for moth genetic variation. When gene flow takes place between well-differentiated genetic populations the result is referred to as "genetic admixture". Admixture mapping is a technique used to study how genetic variants cause differences in disease rates between population. Recent admixture populations that trace their ancestry to multiple continents are well suited for identifying genes for traits and diseases that differ in prevalence between parental populations. African-American populations have been the focus of numerous population genetic and admixture mapping studies, including studies of complex genetic traits such as white cell count, body-mass index, prostate cancer and renal disease. An analysis of phenotypic and genetic variation including skin color and socio-economic status was carried out in the population of Cape Verde which has a well documented history of contact between Europeans and Africans. The studies showed that pattern of admixture in this population has been sex-biased and there is a significant interactions between socio economic status and skin color independent of the skin color and ancestry. Another study shows an increased risk of graft-versus-host disease complications after transplantation due to genetic variants in human leukocyte antigen (HLA) and non-HLA proteins. Differences in allele frequencies contribute to group differences in the incidence of some monogenic diseases, and they may contribute to differences in the incidence of some common diseases (Risch et al. 2002; Burchard et al. 2003; Tate and Goldstein 2004). For the monogenic diseases, the frequency of causative alleles usually correlates best with ancestry, whether familial (for example, Ellis-van Creveld syndrome among the Pennsylvania Amish), ethnic (Tay-Sachs disease among Ashkenazi Jewish populations), or geographical (hemoglobinopathies among people with ancestors who lived in malarial regions). To the extent that ancestry corresponds with racial or ethnic groups or subgroups, the incidence of monogenic diseases can differ between groups categorized by race or ethnicity, and health-care professionals typically take these patterns into account in making diagnoses. Even with common diseases involving numerous genetic variants and environmental factors, investigators point to evidence suggesting the involvement of differentially distributed alleles with small to moderate effects. Frequently cited examples include hypertension (Douglas et al. 1996), diabetes (Gower et al. 2003), obesity (Fernandez et al. 2003), and prostate cancer (Platz et al. 2000). However, in none of these cases has allelic variation in a susceptibility gene been shown to account for a significant fraction of the difference in disease prevalence among groups, and the role of genetic factors in generating these differences remains uncertain (Mountain and Risch 2004). Neil Risch of Stanford University has proposed that self-identified race/ethnic group could be a valid means of categorization in the USA for public health and policy considerations. 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Science 318 (5858): 1842–1843. doi:10.1126/science.318.5858.1842. PMID 18096770. - Ramachandran, S; Tang, H; Gutenkunst, RN; Bustamante, CD (2010). "Genetics and Genomics of Human Population Structure". In Speicher, MR; Antonarakis, SE; Motulsky, AG. Vogel and Motulsky’s Human Genetics: Problems and Approaches (4th ed.). Springer. ISBN 3-540-37653-4.

https://en.wikipedia.org/wiki/Human_genetic_variation

Teaching Word Meanings / Edition 1by Steven A. Stahl, William E. Nagy Pub. Date: 09/26/2005 Publisher: Taylor & Francis Learning new words is foundational to success in school and life. Researchers have known for years that how many word meanings a student knows is one of the strongest predictors of how well that student will understand text and be able to communicate through writing. This book is about how children learn the meanings of new words (and the concepts they convey) and… See more details below Learning new words is foundational to success in school and life. Researchers have known for years that how many word meanings a student knows is one of the strongest predictors of how well that student will understand text and be able to communicate through writing. This book is about how children learn the meanings of new words (and the concepts they convey) and how teachers can be strategic in deciding which words to teach, how to teach them, and which words not to teach at all. This book offers a comprehensive approach to vocabulary instruction. It offers not just practical classroom activities for teaching words (though plenty of those are included), but ways that teachers can make the entire curriculum more effective at promoting students' vocabulary growth. It covers the 'why to' and 'when to' as well as the 'how to' of teaching word meanings. Key features of this exciting new book include: *A variety of vocabulary activities. Activities for teaching different kinds of words such as high frequency words, high utility words, and new concepts, are explained and illustrated. *Guidelines for choosing words. A chart provides a simple framework built around seven basic categories of words that helps teachers decide which words to teach and how to teach them. *Word learning strategies. Strategies are offered that will help students use context, word parts, and dictionaries more effectively. *Developing Word Consciousness. Although specific vocabulary instruction is fully covered, the primary goal of this book is to develop students' independent interest in words and their motivation to learn them. *Integrated Vocabulary Instruction. Teachers are encouraged to improve the reading vocabularies of their students by looking for opportunities to integrate vocabulary learning into activities that are undertaken for other purposes. Table of Contents Contents: Preface. Part I: The Lay of the Land. The Importance of Vocabulary. Vocabulary Knowledge, Reading Comprehension, and Readability. Problems and Complexities. A Comprehensive Approach to Vocabulary Learning. Part II: Teaching Specific Words. Teaching Words for Ownership. Teaching Concepts. Teaching High-Frequency Words. Talking About Words. Part III: Independent Word Learning. Exposure to Rich Language. Promoting Word Consciousness. Teaching Word Learning Strategies: Word Parts. Teaching Word Learning Strategies: Context. Teaching Word Learning Strategies: Definitions. Conclusion: Matching Instructional Approaches to Students and Words. and post it to your social network Most Helpful Customer Reviews See all customer reviews >

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From Diversifying Economic Quality: A Wiki for Instructors and Departments Active learning is an educational approach in which educators enable students to construct their understanding, teaching them to become problem solvers and critical thinkers. In contrast to a classic 'chalk and talk' presentation in which an instructor delivers information to students, students learn how to gather, analyze, and evaluate information themselves. Studies show that active learning can help encourage women and racial minority groups in fields where they are underrepresented. Multiple studies have looked at the effects of active learning in the natural sciences, areas in which women and racial minorities are also underrepresented. Interactive engagement methods such as Peer Instruction have been reported to increase understanding for all students and to decrease the gender gap in Physics. Other work in the sciences, such as that motivating the The National Academies Summer Institutes on Undergraduate Education , similarly suggests that active learning improves all students' comprehension of materials and may be particularly beneficial to female, African American, and Latino students. Students of teachers who emphasized interest in science, further study in science, and experimental methods had higher scores; this benefit was significantly greater for underrepresented minority students. According to an overview of past research about the effectiveness of inquiry based learning , past projects to increase gender ratios in the sciences have found improved success rates with an inquiry-based teaching compared to traditional lecture formats. How to Incorporate Active Learning Incorporating active learning into the classroom requires changing the environment from one of passive information reception to one of inquiry and desire to understand. This shift in thought can be achieved by adopting several related practices and attitudes. - Become familiar with Bloom's Taxonomy and help your students move up the pyramid. - Use Think-pair-share. - Use One-minute papers. - "The one-minute paper is a "modest, relatively simple and low-tech" innovation designed to obtain regular feedback from students. In the final minute or two of class, the teacher asks students to respond to the following two questions: - 1. What is the most important thing you learned today? - 2. What is the muddiest point still remaining at the conclusion of today's class?" - Using an experimental design, John F Chizmar and Anthony L. Ostrosky (1998) report an approximate 6.6 percent increase in economic knowledge relative to pre-treatment levels. - Flip your classroom. - Participate in the Wikipedia Education Program. - Use Peer Instruction. - Peer Instruction is an instructional strategy that works even in large classes; it engages students through a structured questioning process involving every student. Harvard researchers implemented and evaluated the method and found "increased student mastery of both conceptual reasoning and quantitative problem solving upon implementing PI." - Pressley, McDaniel, Turnure, Wood, and Ahmad (1987) presented undergraduate students with a list of sentences, each describing the action of a particular man (e.g., “The hungry man got into the car”). Students in the treatment group were prompted to explain “Why did that particular man do that?” Another group of students was instead provided with an explanation for each sentence, and a third group simply read each sentence. On a final test in which participants were cued to recall which man performed each action (e.g., “Who got in the car?”), the treatment group substantially outperformed the other two groups. (Summary from Dunlosky, Rawson, Marsh, Nathan, and Willingham, ) - Translating this result into the economics classroom is feasible and desirable, but it requires a bit more nuance. Generally economics requires students to retain not only the base information but also a particular explanation of that information. Economists could provide follow-up research to identify the efficacy of questioning techniques that lead students to develop and retain this higher order learning. - Emphasize the "how" rather than the "what" of knowledge. - Explain the methods economists used to learn the causes of increased income inequality rather than simply reporting the casues.Thirteen.org By placing an emphasis on the knowledge-creation process, students learn basic concepts and begin to learn how to generate knowledge themselves. - Don't emphasize that there is "one right answer." - An emphasis on a single correct answer to a question discourages student involvement and discourages critical thinking. When students contribute to classroom discussions, identify the value in their comments. Then, clearly explain the generally accepted answer and why that answer is valuable. - Read The Art of Questioning. - Teach using the case method. - Providing students with a case representative of the lesson's educational objective effectively engages them beyond pure memorization. With this method, students develop a solid understanding of the underlying concepts through analysis of the case. Consult this guide to implementing the case method in the economics classroom. - Use problem sets with context-rich problems. - Problem sets effectively engage and challenge students by requiring them to comprehend and use concepts from the lesson. For a guide on using context-rich problems in the Economics classroom read here. - Schedule periodic recitation sessions with students. Other Examples of Active, Inquiry-Based Learning Stephen D. Morris (Department of Economics, University of California, San Diego) presents research-based suggestions for improving the teaching of AS/AD in his paper, Teaching General Equilibrium to Undergraduates: A Graphical Approach. See "Focus on Inquiry: A Teacher's Guide to Implementing Inquiry-Based Learning" by the Alberta Ministry of Learning , and a similar, shorter document from Penn State. A Dream Experiment in Development Economics by Prakarsh Singh & Alexa Russo, The Journal of Economic Education (Volume 44, Issue 2, 2013) Partial-immersion language programs promote language acquisition through active use rather than through memorization of vocabulary and verb conjugations. See Thirteen.org. Additional Evidence and Research Becker & Watts, 2008. This paper examines how economics was taught in four different undergraduate classes in colleges and universities. U.S. academic economists filled out a survey in 1995, 2000, and 2005, and researchers compared the responses to see how teaching methods changed throughout this decade. During this decade, there was nationally a greater focus on encouraging instructors to spend more time, attention, and effort on teaching, especially through active, student-centered teaching methods (i.e., less use of direct instruction, known colloquially as 'chalk-and-talk'). By 2005, more instructors were using other teaching methods beyond chalk-and-talk, such as classroom discussions, lecture notes provided in hard-copy or online, and computer lab assignments in econometrics/ statistics courses. Additionally, a small but growing minority of instructors used internet database searches, classroom experiments, or assignments referencing current financial news, sports, literature, drama, and music. Cooperative learning methods were used much less frequently. Click here to access it. This study can be found in The Journal of Economic Education. Becker & Watts, 2001. In this article, the authors compare the results of surveys on teaching style conducted in 1995 and then again in 2000. They found that although higher-education institutions have effectively shifted from professors' focus from being more research-oriented to being more focused on their teaching, outdated teaching methods still permeate the discipline. From the surveys conducted, the authors see that classroom presentations are still dominated by the "chalk and talk" method. The authors also find that teacher-student discussion does not occur until until upper level courses, and student-student discussion is rare for the discipline as a whole. On a similar note, it is observed that the use of multiple-choice test formats seems to be excessive--especially in introductory theory courses. Click here to access the article. Major & Palmer, 2001. "Problem‑Based Learning (PBL) is an innovative educational approach that is gaining prominence in higher education. A review of the literature of PBL outcomes summarizes, across multiple studies, the positive effects of problem‑based learning. Since PBL brings with it unique challenges to traditional assessment, however, this study suggests alternative approaches. Alternative assessment may provide additional insight into the effectiveness of PBL and other alternative pedagogies." Click here to view it. Crouch, Watkins, Fagen, and Mazur, 2007. "Peer Instruction is an instructional strategy for engaging students during class through a structured questioning process that involves every student. We describe Peer Instruction (hereafter PI) and report data from more than ten years of teaching with PI in the calculus- and algebra-based introductory physics courses for non-majors at Harvard University, where this method was developed. Our results indicate increased student mastery of both conceptual reasoning and quantitative problem solving upon implementing PI." See link provided above. The rate of information dissemination has dramatically increased, due to technological development and global interconnection. An educational system that focuses on memorization and lower order cognitive skills is obsolete and inefficient. Instead, curricula can be reorganized to emphasize problem-solving and other higher order skills. Through inquiry-based teaching practices, educators create an environment of inquiry, helping students to seek more than simple answers, to explore the mechanisms underlying what is known, and to learn how to create knowledge themselves. von Secker, Clare. 2002. "Effects of Inquiry-Based Teacher Practices on Science Excellence and Equity" The Journal of Educational Research, 95, 3. 151-160. "Inquiry-based Learning: Explanation." THIRTEEN - New York Public Media. Web. 03 June 2011. <http://www.thirteen.org/edonline/concept2class/inquiry/index_sub7.html>. "Inquiry-based Learning." Printable Worksheets for Teachers and Students. Web. 01 June 2011. <http://www.worksheetlibrary.com/teachingtips/inquiry.html>. "World Language - Partial Immersion." FCPS Home Page Redirect Page. Web. 10 June 2011. <http://www.fcps.edu/DIS/OHSICS/forlang/partial.htm>. Becker, William E., and Michael Watts. "Teaching Methods in U. S. Undergraduate Economics Courses." The Journal of Economic Education 32.3 (2001): 269-79. Web. <http://www.jstor.org/pss/1183384>. Kinkead, Joyce. 2003. "Learning Through Inquiry: An Overview of Undergraduate Research" NEW DIRECTIONS FOR TEACHING AND LEARNING, 93. 5-17. <ftp://charmian.sonoma.edu/pub/references/Kinkead.pdf>

http://www.diversifyingecon.org/index.php/Inquiry-based_learning

planning a fair test ... representing data in line graphs and interpreting what these show. Work in this unit ... Belair Lesson 1 uses sand, flour and water. Letts pp ... This printable Math lesson plan outline named "Year 6 Science Unit 6C More about dissolving" provides more information about lesson 1, graphs, etc. To make sure that this file is what you need, before you download this Math lesson plan outline, try to read this file first by click the following link. On the other hand, if you want to save this file directly into your computer, you can download this pdf Math lesson plan outline through the following download link. Students of pre-kindergarten through 8th grade can be identified as young learners. Producing teaching note for young learners should be imaginative in selecting the learning practices because of their characteristics. It is better for educator to understand the characteristics of young learners before designing the lesson plan for them. Wendy A. Scott and Lisbeth H. Ytreberg describe the characteristics of young learners. Here we will talk it in related to the class concept. First, Young learners like to co-operate and physically active. That is why teachers have to design the practices which involve the students to participate for example: desigining a work in group or individually, like interesting quizes which involved the physical practices. Second, young learners are happy to play. They study well when they are enjoying themselves. In related in playing, we now have many improvements in teaching Math strategy for children through games or even the colored and entertaining worksheet design. It is easy to find where we can find the fun Math worksheet on the online resources. Third, young learners cannot focus for a long time. Teacher should have a great mannerin dividing teaching time from the beginning till the end of the class. Young learners are happier with different materials and they cannot remember things for a long time if it is not repeated. So keep repeating the lesson with different fun ways. Five, seven or twelve years old Learners will grow as thinkers who can be trustworthy and take responsibility for class practices and routines. They will also study how to play and organize the best way to maintain an activity, work with others and study from others. Moreover, young Learners still depend on teacher. They should be guided and accompanied well. So keep guiding students with the best way. Teaching Math is fun and let us make them happy to learn Math.

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An allosome (also referred to as a sex chromosome, heterotypical chromosome, heterochromosome, or idiochromosome) is a chromosome that differs from an ordinary autosome in form, size, and behavior. The human sex chromosomes, a typical pair of mammal allosomes, determine the sex of an individual created in sexual reproduction. Autosomes differ from allosomes because autosomes appear in pairs whose members have the same form but differ from other pairs in a diploid cell, whereas members of an allosome pair may differ from one another and thereby determine sex. In humans, each cell nucleus contains 23 pairs of chromosomes a total of 46 chromosomes. The first 22 pairs are called autosomes which look exactly the same in both males and females. The 23rd pair of chromosomes is called an allosome. These sex chromosomes usually differ between males and females. females have two copies of the X chromosome, while males have one X chromosome and one Y chromosome. The X chromosome is always present as the 23rd chromosome in the ovum, while either X or Y chromosomes can be present in an individual sperm. All diploid organisms with allosome-determined genders get half of their allosomes from each of their parents. In mammals, females are XX, they can pass along either of their X’s, and since the males are XY they can pass along either an X or a Y. For a mammal to be considered a female, the individual must receive an X chromosome from both parents, whereas to be considered a male, the individual must receive a X chromosome from their mother and a Y chromosome from their father. It is thus the male’s sperm that determines the sex of each offspring in humans. There is about a 51 percent chance of producing a male offspring and 49 percent a female, with Fisher's principle determining this sex ratio. Allosomes not only carry the genes that determine male and female traits, but also those for some other characteristics as well. Genes that are carried by either sex chromosome are said to be sex linked. Sex-linked diseases are passed down through families through one of the X or Y chromosomes. Since only men inherit Y chromosomes, they are the only ones to inherit Y-linked traits. Men and women can get the X-linked ones since both inherit X chromosomes. A gene is either said to be dominant or recessive. Dominant inheritance occurs when an abnormal gene from one parent causes disease even though the matching gene from the other parent is normal. The abnormal gene dominates. Recessive inheritance is when both matching genes must be abnormal to cause disease. If only one gene in the pair is abnormal, the disease does not occur, or is mild. Someone who has one abnormal gene (but no symptoms) is called a carrier. A carrier can pass this abnormal gene to his or her children. X chromosome carry about 1500 genes, more than any other chromosome in the human body. Most of them code for something other than female anatomical traits. Many of the non-sex determining X-linked genes are responsible for abnormal conditions. The Y chromosome carries about 78 genes. Most of the Y chromosome genes are involved with essential cell house-keeping activities and sperm production. Only one of the Y chromosome genes, the SRY gene, is responsible for male anatomical traits. When any of the 9 genes involved in sperm production are missing or defective the result is usually very low sperm counts and infertility. Examples of mutations on the X chromosome include more common diseases such as color blindness, hemophilia, and fragile-X syndrome. - Color blindness or color vision deficiency is the inability or decreased ability to see color, or perceive color differences, under normal lighting conditions. Color blindness affects many individuals in the population. There is no actual blindness, but there is a deficiency of color vision. The most usual cause is a fault in the development of one or more sets of retinal cones that perceive color in light and transmit that information to the optic nerve. This type of color blindness is usually a sex-linked condition. The genes that produce photopigments are carried on the X chromosome; if some of these genes are missing or damaged, color blindness will be expressed in males with a higher probability than in females because males only have one X chromosome. - Hemophilia refers to a group of bleeding disorders in which it takes a long time for the blood to clot. This is referred to as X-Linked recessive. Hemophilia is much more common in males than females because males are hemizygous. They only have one copy of the gene in question and therefore express the trait when they inherit one mutant allele. In contrast, a female must inherit two mutant alleles, a less frequent event since the mutant allele is rare in the population. X-linked traits are maternally inherited from carrier mothers or from an affected father. Each son born to a carrier mother has a 50% probability of inheriting the X-chromosome carrying the mutant allele. - Queen Victoria was a carrier of the gene for hemophilia. She passed on the harmful allele to one of her four sons and at least two of her five daughters. Her son Leopold had the disease and died at age 30. As a result of marrying into other European royal families, the princesses Alice and Beatrice spread hemophilia to Russia, Germany, and Spain. By the early 20th century, ten of Victoria's descendents had hemophilia. All of them were men, as expected. - Fragile X syndrome is a genetic condition involving changes in part of the X chromosome. It is the most common form of inherited intellectual disability (mental retardation) in males. It is caused by a change in a gene called FMR1. A small part of the gene code is repeated on a fragile area of the X chromosome. The more repeats, the more likely there is to be a problem. Males and females can both be affected, but because males have only one X chromosome, a single fragile X is likely to affect them more. Most fragile-X males have large testes, big ears, narrow faces, and sensory processing disorders that result in learning disabilities. Other complications include: - 46,XX testicular disorder of sex development, also called XX male syndrome, is a condition in which individuals with two X chromosomes in each cell, the pattern normally found in females, have a male appearance. People with this disorder have male external genitalia. In most people with 46,XX testicular disorder of sex development, the condition results from an abnormal exchange of genetic material between chromosomes (translocation). This exchange occurs as a random event during the formation of sperm cells in the affected person's father. The SRY gene (which is on the Y chromosome) is misplaced in this disorder, almost always onto an X chromosome. Anyone with an X chromosome that carries the SRY gene will develop male characteristics despite not having a Y chromosome. - Fisher's principle - Haldane's rule - XY sex-determination system - ZW sex-determination system - X0 sex-determination system - How many chromosomes do people have? - Genetics Home Reference - Biological Basis of Heredity: Sex Linked Genes - Sex-linked recessive - National Library of Medicine - PubMed Health - Sex-Linked Traits | Heredity and genetics | Khan Academy - Hemophilia - National Library of Medicine - PubMed Health - Fragile X Syndrome - Symptoms, Diagnosis, Treatment of Fragile X Syndrome - NY Times Health Information - 46,XX testicular disorder of sex development - Genetics Home Reference |Wikimedia Commons has media related to Sex chromosomes.|

https://en.wikipedia.org/wiki/Allosome

In computer language design, stropping is a method of explicitly marking letter sequences as having a special property such as being a keyword or certain type of variable or storage location, and thus inhabiting a different namespace from ordinary names ("identifiers"), avoiding clashes. Stropping is not used in most modern languages – instead, keywords are reserved words and cannot be used as identifiers. Stropping allows the same letter sequence to be used both as a keyword and as an identifier, and simplifies parsing in that case – for example allowing a variable named if without clashing with the keyword if. The method of stropping and the term "stropping" arose in the development of ALGOL in the 1960s, where it was used to represent typographical distinctions (boldface and underline) found in the publication language which could not directly be represented in the hardware language – a typewriter could have bold characters, but in encoding in punch cards there were no bold characters. The term "stropping" arose in ALGOL 60, from "apostrophe", as some implementations of ALGOL 60 used apostrophes around text to indicate boldface, such as 'if' to represent the keyword if. Stropping is also important in ALGOL 68, where multiple methods of stropping, known as "stropping regimes", are used; the original matched apostrophes from ALGOL 60 was not widely used, with a leading period or uppercase being more common, as in IF and the term "stropping" was applied to all of these. A range of different syntaxes for stropping have been used: - Algol 60 commonly used only the convention of single quotes around the word, generally as apostrophes, whence the name "stropping" (e.g. - Algol 68 in some implementations treat letter sequences prefixed by a single quote, ', as being keywords (e.g., In fact it was often the case that several stropping conventions might be in use within one language. For example, in ALGOL 68, the choice of stropping convention can be specified by a compiler directive (in ALGOL terminology, a "pragmat"), namely POINT, UPPER, QUOTE, or RES: - POINT for 6-bit (not enough characters for lowercase), as in .FOR– a similar convention is used in FORTRAN 77, where LOGICAL keywords are stropped as .EQ.etc. (see below) - UPPER for 7-bit, as in FOR– with lowercase used for ordinary identifiers - QUOTE as in ALGOL 60, as in - RES reserved words, as used in modern languages – foris reserved and not available to ordinary identifiers The various rules regimes are a lexical specification for stropped characters, though in some cases these have simple interpretations: in the single apostrophe and dot regimes, the first character is functioning as an escape character, while in the matched apostrophes regime the apostrophes are functioning as delimiters, as in string literals. - Atlas Autocode had the choice of three: keywords could be underlinedusing backspace and overstrike on a Flexowriter keyboard, they could be introduced by a %percent %symbol, or they could be typed in UPPER CASEwith no delimiting character ("uppercasedelimiters" mode, in which case all variables had to be in lower case). - ALGOL 68RS programs are allowed the use of several stropping variants, even within the one language processor. Note the leading pr (abbreviation of pragmat) directive, which is itself stropped in POINT or quote style, and the ¢ for comment (from "2¢") – see ALGOL 68: pr & co: Pragmats and Comments for details. as typically published ¢ underline or bold typeface ¢ mode xint = int; xint sum sq:=0; for i while sum sq≠70×70 do sum sq+:=i↑2 od 'pr' quote 'pr' 'mode' 'xint' = 'int'; 'xint' sum sq:=0; 'for' i 'while' sum sq≠70×70 'do' sum sq+:=i↑2 'od' |For a 7-bit character code compiler .PR UPPER .PR MODE XINT = INT; XINT sum sq:=0; FOR i WHILE sum sq/=70*70 DO sum sq+:=i**2 OD |For a 6-bit character code compiler .PR POINT .PR .MODE .XINT = .INT; .XINT SUM SQ:=0; .FOR I .WHILE SUM SQ .NE 70*70 .DO SUM SQ .PLUSAB I .UP 2 .OD |Algol68 using res stropping .PR RES .PR mode .xint = int; .xint sum sq:=0; for i while sum sq≠70×70 do sum sq+:=i↑2 od Most modern computer languages do not use stropping, with one notable exception. The use of many languages in Microsoft's .NET Common Language Infrastructure (CLI) requires a way to use variables in a different language that may be keywords in a calling language. This is sometimes done by prefixes, such as @ in C#, or enclosing the identifier in brackets, in Visual Basic.NET. There are other, more minor examples. For example, Web IDL uses a leading underscore _ to strop identifiers that otherwise collide with reserved words: the value of the identifier strips this leading underscore, making this stropping, rather than a naming convention. Unstropping by the compiler In a compiler frontend, unstropping originally occurred during an initial line reconstruction phase, which also eliminated whitespace. This was then followed by scannerless parsing (no tokenization); this was standard in the 1960s, notably for ALGOL. In modern use, unstropping is generally done as part of lexical analysis. This is clear if one distinguishes the lexer into two phases of scanner and evaluator: the scanner categorizes the stropped sequence into the correct category, and then the evaluator unstrops when calculating the value. For example, in a language where an initial underscore is used to strop identifiers to avoid collisions with reserved words, the sequence _if would be categorized as an identifier (not as the reserved word if) by the scanner, and then the evaluator would give this the value (Identifier, if) as the token type and value. A number of similar techniques exist, generally prefixing or suffixing an identifier to indicate different treatment, but the semantics are varied. Strictly speaking, stropping consists of different representations of the same name (value) in different namespaces, and occurs at the tokenization stage. For example, in ALGOL 60 with matched apostrophe stropping, 'if' is tokenized as (Keyword, if), while if is tokenized as (Identifier, if) – same value in different token classes. Using uppercase for keywords remains in use as a convention for writing grammars for lexing and parsing – tokenizing the reserved word if as the token class IF, and then representing an if-then-else clause by the phrase IF Expression THEN Statement ELSE Statement where uppercase terms are keywords and capitalized terms are nonterminal symbols in a production rule (terminal symbols are denoted by lowercase terms, such as integer, for an integer literal). Most loosely, one may use naming conventions to avoid clashes, commonly prefixing or suffixing with an underscore, as in _then. A leading underscore is often used to indicate private members in object-oriented programming. These names may be interpreted by the compiler and have some effect, though this is generally done at the semantic analysis phase, not the tokenization phase. For example, in Python, a single leading underscore is a weak private indicator, and affects which identifiers are imported on module import, while a double leading underscore (and no more than one trailing underscore) on a class attribute invokes name mangling. While modern languages generally use reserved words rather than stropping to distinguish keywords from identifiers – e.g., making if reserved – they also frequently reserve a syntactic class of identifiers as keywords, yielding representations which can be interpreted as a stropping regime, but instead have the semantics of reserved words. This is most notable in C, where identifiers that begin with an underscore are reserved, though the precise details of what identifiers are reserved at what scope are involved, and leading double underscores are reserved for any use; similarly in C++ any identifier that contains a double underscore is reserved for any use, while an identifier that begins with an underscore is reserved in the global space.[a] Thus one can add a new keyword foo using the reserved word __foo. While this is superficially similar to stropping, the semantics are different. As a reserved word, the string __foo represents the identifier __foo in the common identifier namespace. In stropping (by prefixing keywords by __), the string __foo represents the keyword foo in a separate keyword namespace. Thus using reserved words, the tokens for foo are (identifier, __foo) and (identifier, foo) – different values in the same category – while in stropping the tokens for foo are (keyword, foo) and (identifier, foo) – same values in different categories. These solve the same problem of namespace clashes in a way that is the same for a programmer, but which differs in terms of formal grammar and implementation. Name mangling also addresses name clashes by renaming identifiers, but does this much later in compilation, during semantic analysis, not during tokenization. This consists of creating names that include scope and type information, primarily for use by linkers, both to avoid clashes and to include necessary semantic information in the name itself. In these cases the original identifiers may be identical, but the context is different, as in the functions foo(int x) versus foo(char x), in both cases having the same identifier foo, but different signature. These names might be mangled to foo_c, for instance, to include the type information. A syntactically similar but semantically different phenomenon are sigils, which instead indicate properties of variables. These are common in Perl, Ruby, and various other languages to identify characteristics of variables/constants: Perl to designate the type of variable, Ruby to distinguish variables from constants and to indicate scope. Note that this affects the semantics of the variable, not the syntax of whether it is an identifier or keyword. - There are other restrictions, such as an identifier that begins with an underscore, followed by an uppercase letter. - Proceedings of an International Conference on ALGOL 68 Implementation: Department of Computer Science, University of Manitoba, Winnipeg, June 18–20, 1974, ed. Peter R. King, University of Manitoba. Dept. of Computer Science, p. 148 – More serious problems are posed by "stropping," the technique used to distinguish boldface text from roman text. Some implementations demand apostrophes around boldface (whence the name stropping); others require backspacing and underlining; ... - Revised Report, p. 123, footnote - van Wijngarten et al. (1976) Section 9.3 - Lindsey and van der Meulen (1977) pp.348-349 - Web IDL, "3.1. Names": For all of these constructs, the identifier is the value of the identifier token with any single leading U+005F LOW LINE ("_") character (underscore) removed. Note A leading "_" is used to escape an identifier from looking like a reserved word so that, for example, an interface named “interface” can be defined. The leading "_" is dropped to unescape the identifier. - PEP 008: Descriptive: Naming Styles - C99 standard, 7.1.3 Reserved identifiers - A. van Wijngaarden; et al. (1976). Revised Report on the Algorithmic Language ALGOL 68. Springer-Verlag. ISBN 0-387-07592-5. OCLC 1991170. - C. H. Lindsey; S. G. van der Meulen (1977). Informal Introduction to ALGOL 68. North-Holland. ISBN 0-7204-0726-5. OCLC 230034877. Cite uses deprecated parameter - W. J. Hansen; H. J. Boom (1978). "Report on the Standard Hardware Representation for Revised ALGOL 68". Acta Informatica 9: 105–119. doi:10.1007/BF00289072. - C.H. Lindsey, "An ISO-Code Representation for ALGOL 68" ALGOL Bulletin AB31.3.6, Issue 31, March 1970, pp. 37–60 ACM

https://en.wikipedia.org/wiki/Stropping_(syntax)

Fossil, remnant, impression, or trace of an animal or plant of a past geologic age that has been preserved in Earth’s crust. The complex of data recorded in fossils worldwide—known as the fossil record—is the primary source of information about the history of life on Earth. Only a small fraction of ancient organisms are preserved as fossils, and usually only organisms that have a solid and resistant skeleton are readily preserved. Most major groups of invertebrate animals have a calcareous skeleton or shell (e.g., corals, mollusks, brachiopods, bryozoans). Other forms have shells of calcium phosphate (which also occurs in the bones of vertebrates), or silicon dioxide. A shell or bone that is buried quickly after deposition may retain these organic tissues, though they become petrified (converted to a stony substance) over time. Unaltered hard parts, such as the shells of clams or brachiopods, are relatively common in sedimentary rocks, some of great age. The hard parts of organisms that become buried in sediment may be subject to a variety of other changes during their conversion to solid rock, however. Solutions may fill the interstices, or pores, of the shell or bone with calcium carbonate or other mineral salts and thus fossilize the remains, in a process known as permineralization. In other cases there may be a total replacement of the original skeletal material by other mineral matter, a process known as mineralization, or replacement. In still other cases, circulating acid solutions may dissolve the original shell but leave a cavity corresponding to it, and circulating calcareous or siliceous solutions may then deposit a new matrix in the cavity, thus creating a new impression of the original shell. By contrast, the soft parts of animals or plants are very rarely preserved. The embedding of insects in amber (a process called resin fossilization) and the preservation of the carcasses of Pleistocene mammoths in ice are rare but striking examples of the fossil preservation of soft tissues. Leaves, stems, and other vegetable matter may be preserved through the process of carbonization, where such parts are flattened between two layers of rock. The chemical reduction of the part produces a carbon film that occurs on one layer of rock, while an impression of that part occurs on the other layer of the rock. Fossils of hard and soft parts that are too small to be observed by the naked eye are called microfossils. Some fossils are completely devoid of plant and animal parts but show evidence of an organism’s activities. Such traces of organisms, which are appropriately known as “trace fossils,” include tracks or trails, preserved waste products, and borings. The great majority of fossils are preserved in a water environment because land remains are more easily destroyed. Anaerobic conditions at the bottom of the seas or other bodies of water are especially favourable for preserving fine details, since no bottom faunas, except for anaerobic bacteria, are present to destroy the remains. In general, for an organism to be preserved two conditions must be met: rapid burial to retard decomposition and to prevent the ravaging of scavengers; and possession of hard parts capable of being fossilized. In some places, such as the Grand Canyon in northern Arizona, one can observe a great thickness of nearly horizontal strata representing the deposition of sediment on the seafloor over many hundreds of millions of years. It is often apparent that each layer in such a sequence contains fossils that are distinct from those of the layers that are above and below it. In such sequences of layers in different geographic locations, the same, or similar, fossil floras or faunas occur in the identical order. By comparing overlapping sequences, it is possible to build up a continuous record of faunas and floras that have progressively more in common with present-day life forms as the top of the sequence is approached. The study of the fossil record has provided important information for at least four different purposes. The progressive changes observed within an animal group are used to describe the evolution of that group. Fossils also provide the geologist a quick and easy way of assigning a relative age to the strata in which they occur. The precision with which this may be done in any particular case depends on the nature and abundance of the fauna: some fossil groups were deposited during much longer time intervals than others. Fossils used to identify geologic relationships are known as index fossils. Fossil organisms may provide information about the climate and environment of the site where they were deposited and preserved (e.g., certain species of coral require warm, shallow water, or certain forms of deciduous angiosperms can only grow in colder climatic conditions). Fossils are useful in the exploration for minerals and mineral fuels. For example, they serve to indicate the stratigraphic position of coal seams. In recent years, geologists have been able to study the subsurface stratigraphy of oil and natural gas deposits by analyzing microfossils obtained from core samples of deep borings. Fossil collection as performed by paleontologists, geologists, and other scientists typically involves a rigorous excavation and documentation process. Unearthing the specimen from the rock is often painstaking work that includes labeling each part of the specimen and cataloging the location of each part within the rock. Those fossils slated for removal from the rock are slowly and carefully excavated using techniques designed to prevent or minimize damage to the specimen. Such fossils often become part of museum or university collections. Many other fossils, however, are collected by hobbyists and commercial entities. Often such specimens are not carefully documented or excavated, resulting in a loss of data from the site and risking potential damage to the specimen. For these reasons and the fact that it stimulates nonscientific collecting, the commercial exploitation of fossils is controversial among academic paleontologists.

http://www.britannica.com/science/fossil

Anatomical terms of location |This article is part of a series on| All vertebrates (including humans) have the same basic body plan — they are bilaterally symmetrical. That is, they have mirror-image left and right halves if divided down the centre. For these reasons, the basic directional terms can be considered to be those used in vertebrates. By extension, the same terms are used for many other (invertebrate) organisms as well. While these terms are standardized within specific fields of biology, there are unavoidable, sometimes dramatic, differences between some disciplines. For example, differences in terminology remain a problem that, to some extent, still separates the terminology of human anatomy from that used in the study of various other zoological categories. - 1 Introduction - 2 Main terms - 3 Other terms and special cases - 4 Specific animals and other organisms - 5 Citations - 6 Sources Standardized anatomical and zoological terms of location have been developed, usually based on Latin and Greek words, to enable all biological and medical scientists to precisely delineate and communicate information about animal bodies and their component organs, even though the meaning of some of the terms often is context-sensitive. The vertebrates and Craniata share a substantial heritage and common structure, so many of the same terms are used to describe location. To avoid ambiguities this terminology is based on the anatomy of each animal in a standard way. For humans, one type of vertebrate, anatomical terms may differ from other forms of vertebrates. For one reason, this is because humans have a different neuraxis and, unlike animals that rest on four limbs, humans are considered when describing anatomy as being in the standard anatomical position. Thus what is on "top" of a human is the head, whereas the "top" of a dog may be its back, and the "top" of a flounder could refer to either its left or its right side. For invertebrates, standard application of locational terminology often becomes difficult or debatable at best when the differences in morphology are so radical that common concepts are not homologous and do not refer to common concepts. For example, many species are not even bilaterally symmetrical. In these species, terminology depends on their type of symmetry (if any). Standard anatomical position Because animals can change orientation with respect to their environment, and because appendages like limbs and tentacles can change position with respect to the main body, positional descriptive terms need to refer to the animal as in its standard anatomical position. All descriptions are with respect to the organism in its standard anatomical position, even when the organism in question has appendages in another position. This helps avoid confusion in terminology when referring to the same organism in different postures. In humans, this refers to the body in a standing position with arms at the side and palms facing forward (thumbs out). While the universal vertebrate terminology used in veterinary medicine would work in human medicine, the human terms are thought to be too well established to be worth changing. Many anatomical terms can be combined, either to indicate a position in two axes simultaneously or to indicate the direction of a movement relative to the body. For example, "anterolateral" indicates a position that is both anterior and lateral to the body axis (such as the bulk of the pectoralis major muscle). In radiology, an X-ray image may be said to be "anteroposterior", indicating that the beam of X-rays pass from their source to patient's anterior body wall through the body to exit through posterior body wall. There is no definite limit to the contexts in which terms may be modified to qualify each other in such combinations. Generally the modifier term is truncated and an "o" is added in prefixing it to the qualified term. For example, a view of an animal from an aspect dorsal and lateral might be called "dorsolateral". Where desirable three or more terms may be agglutinated or concatenated, as in "anteriodorsolateral". Such terms sometimes used to be hyphenated, but the modern tendency is to omit the hyphen. There is however little basis for any strict rule to interfere with choice of convenience in such use. Three basic reference planes are used to describe location. - The sagittal plane is a plane parallel to the sagittal suture, All other sagittal planes (referred to as parasagittal planes) are parallel to it. It is also known as a "longitudinal plane". The plane is an Y-Z plane, perpendicular to the ground. - The median plane or midsagittal plane is in the midline of the body, and divides the body into left and right (sinister and dexter) portions. This passes through the head, spinal cord, navel and, in animals, the tail. The median plane can also refer to the midsagittal plane of other structures, such as a digit. - The frontal plane or coronal plane divides the body into dorsal and ventral (back and front, or posterior and anterior) portions. For post-embryonic humans a coronal plane is vertical and a transverse plane is horizontal, but for embryos and quadrupeds a coronal plane is horizontal and a transverse plane is vertical. A longitudinal plane is any plane perpendicular to the transverse plane. The coronal plane and the sagittal plane are examples of longitudinal planes. - A transverse plane, also known as a cross-section, divides the body into cranial and caudal (head and tail) portions. - In human anatomy - A transverse (also known as horizontal) plane is an X-Y plane, parallel to the ground, which (in humans) separates the superior from the inferior, or put another way, the head from the feet. - A coronal (also known as frontal) plane is a X-Z plane, perpendicular to the ground, which (in humans) separates the anterior from the posterior, the front from the back, the ventral from the dorsal. |Axis||Directional term||Directed towards| |Anteroposterior (rostrocaudal1,craniocaudal1, cephalocaudal2)||Anterior||Head end| |Dorsoventral||Dorsal||Back, spinal column| |Left-right (dextro-sinister2, sinistro-dexter2)||Left (sinister)||Left-hand side| |Right (dexter)||Right-hand side| |Lateral||Left and right| |Proximal/distal||Proximal||Point at which appendage joins the body| |Distal||Extremity of appendage| (1) Fairly common usage. (2) Uncommon usage. (3) Equivalent to one-half of the left-right axis. (The terms "intermediate", "ipsilateral", "contralateral", "superficial", and "deep", while indicating directions, are relative terms and thus do not properly define fixed anatomical axes. Also, while the "rostrocaudal" and anteroposterior directionality are equivalent in a significant portion of the human body, they are different directions in other parts of the body.) To begin with, distinct, polar-opposite ends of the organism are chosen. By definition, each pair of opposite points defines an axis. In a bilaterally symmetrical organism, there are 6 polar opposite points, giving three axes that intersect at right angles—the x, y, and z axes familiar from three-dimensional geometry. Superior and inferior In anatomical terminology superior (from Latin, meaning "above") is used to refer to what is above something, and inferior (from Latin, meaning "below") to what is below it. For example, in the anatomical position the most superior part of the human body is the head, and the most inferior is the feet. As a second example, in humans the neck is superior to the chest but inferior to the head. Anterior and posterior Anterior refers to what is in front (from Latin ante, meaning "before") and posterior, what is to the back of the subject (from Latin post, meaning "after"). For example, in a dog the nose is anterior to the eyes and the tail is considered the most posterior part; in many fish the gill openings are posterior to the eyes, but anterior to the tail. Medial and lateral Lateral (from Latin lateralis, meaning "to the side") refers to the sides of an animal, as in "left lateral" and "right lateral". The term medial (from Latin medius, meaning "middle") is used to refer to structures close to the centre of an organism, called the "median plane". For example, in a fish the gills are medial to the operculum, but lateral to the heart. The terms "left" and "right" are sometimes used, or their Latin alternatives (Latin: dexter; "right", Latin: sinister; "left"). However, as left and right sides are mirror images, using these words is somewhat confusing, as structures are duplicated on both sides. For example, it's very confusing to say the dorsal fin of a dolphin is "right of" the left pectoral fin, but is "left of" the right eye, but much easier and clearer to say "the dorsal fin is medial to the pectoral fins". Derived terms include: - Contralateral (from Latin contra, meaning "against"): on the side opposite to another structure. For example, the left arm is contralateral to the right arm, or the right leg. - Ipsilateral (from Latin ipse, meaning "same"): on the same side as another structure. For example, the left arm is ipsilateral to the left leg. Proximal and distal The terms proximal (from Latin proximus, meaning "nearest") and distal (from Latin distare, meaning "to stand away from") are used to describe parts of a feature that are close to or distant from the main mass of the body. Thus the upper arm in humans is proximal and the hand is distal. These terms are particularly useful when describing appendages such as fins, tentacles, limbs or indeed any structure that extends that can potentially move separately from the main body. Although the direction indicated by "proximal" and "distal" is always respectively towards or away from the point of attachment, a given structure can be either proximal or distal in relation to another point of reference. Thus the elbow is distal to a wound on the upper arm, but proximal to a wound on the lower arm. Superficial and deep These two terms relate to the distance of a structure from the surface of an animal. Deep (from Old English) refers to something further away from the surface of the organism. For example, the external oblique muscle of the abdomen is deep to the skin. "Deep" is one of the few anatomical terms of location derived from Old English rather than Latin - the anglicised Latin term would have been "profound" (Latin: profundus, "due to depth"). Dorsal and ventral These two terms refer to front/belly (ventral) and back (dorsal) of an organism. Cranial and caudal Specific terms exist to describe how close or far something is to the head or tail of an animals. To describe how close to the head of an animal something is, three distinct terms are used: - Rostral (from Latin rostrum, meaning "beak, nose"): situated toward the oral or nasal region, or in the case of the brain, toward the tip of the frontal lobe. - Cranial (from Greek κρανίον (kranion), meaning "skull") or cephalic (from Greek κεφάλι (kephalē), meaning "head"). To describe how close something is to the end of an organism, the term caudal is used (from Latin cauda, meaning "tail"). In the horse, for example, the eyes are caudal to the nose and rostral to the back of the head. These terms are generally preferred in veterinary medicine and not used as often in human medicine. In humans, "cranial" and "cephalic" are used to refer to the skull, with "cranial" being used more commonly. The term "rostral" is rarely used in human anatomy, apart from embryology, and refers more to the front of the face than the superior aspect of the organism. Similarly, the term "caudal" is only occasionally used in human anatomy. This is because the brain is situated at the superior part of the head whereas the nose is situated in the anterior part. Thus the "rostrocaudal axis" refers to a C shape (see image). Other terms and special cases The location of anatomical structures can also be described with relation to different anatomical landmarks. Structures may be described as being at the level of a specific spinal vertebra, depending on the section of the vertebral column the structure is at. The position is often abbreviated. For example, a structures at the level of the fourth cervical vertebra may be abbreviated as "C4", at the level of a thoracic vertebra "T4", at the level of a lumbar vertebra "L3". Because the sacrum is fused, it is not often used to provide location. References may also take origin from superficial anatomy, made to landmarks that are on the skin or visible underneath. For example, structures may be described relative to the anterior superior iliac spine, the medial malleolus or the medial epicondyle. Anatomical lines, theoretical lines drawn through structures, are also used to describe anatomical location. For example, the mid-clavicular line is used as part of the cardiac exam in medicine to feel the apex beat of the heart. Mouth and teeth Fields such as osteology, palaeontology and dentistry apply special terms of location to describe the mouth and teeth. This is because although teeth may be aligned with their main axes within the jaw, some different relationships require special terminology as well; for example teeth also can be rotated, and in such contexts terms like "anterior" or "lateral" become ambiguous. Terms such as "distal" and "proximal" are also redefined to mean the distance away or close to the mandibular symphysis. Terms used to describe structures include "buccal" (from Latin bucca, meaning "cheek") and "palatal" (from Latin) referring to structures close to the cheek and hard palate respectively. Hands and feet Several unique terms are used to describe the hands and feet For improved clarity, the directional term palmar (from Latin palma, meaning "palm of the hand") is usually used to describe the front of the hand, and dorsal is the back of the hand. For example, the top of a dog's paw is its dorsal surface; the underside, either the palmar (on the forelimb) or the plantar (on the hindlimb) surface. The palmar fascia is palmar to the tendons of muscles which flex the fingers, and the dorsal venous arch is so named because it is on the dorsal side of the foot. Volar can also be used to refer to the underside of the palm or sole, which are themselves also sometimes used to describe location as palmar and plantar. For example, volar pads are those on the underside of hands, fingers, feet, and toes. These terms are used to avoid confusion when describing the median surface of the hand and what is the "anterior" or "posterior" surface — "anterior" can be used to describe the palm of the hand, and "posterior" can be used to describe the back of the hand and arm. This confusion can arise because the forearm can pronate and supinate. Similarly, in the forearm, for clarity, the sides are named after the bones. Structures closer to the radius are radial, structures closer to the ulna are ulnar, and structures relating to both bones are referred to as radioulnar. Similarly, in the lower leg, structures near the tibia (shinbone) are tibial and structures near the fibula are fibular (or peroneal). Most terms of anatomical location are relative to linear motion (translation) along the X- Y- and Z-axes, but there are other degrees of freedom as well, in particular, rotation around any of those three axes. Anteversion and retroversion are complementary anatomical terms of location, describing the degree to which an anatomical structure is rotated forwards (towards the front of the body) or backwards (towards the back of the body) respectively, relative to some datum position. The terms also describe the positioning of surgical implants, such as in arthroplasty. - Anteversion refers to an anatomical structure being tilted further forward than normal, whether pathologically or incidentally. For example, there may be a need to measure the anteversion of the neck of a bone such as a femur. For example, a woman's uterus typically is anteverted, tilted slightly forward. A misaligned pelvis may be anteverted, that is to say tilted forward to some relevant degree. - Retroversion is rotation around the same axis as that of anteversion, but in the opposite sense, that is to say, tilting back. A structure so affected is described as being retroverted. As with anteversion, retroversion is a completely general term and can apply to a backward tilting of such hard structures as bones, soft organs such as uteri, or surgical implants. Other directional terms Several other terms are also used to describe location. These terms are not used to form the fixed axes. Terms include: - Axial (from Latin axōn, meaning "axle"): around the central axis of the organism or the extremity. Two related terms, "abaxial" and "adaxial", refer to locations away and toward the central axis of an organism, respectively. - Intermediate (from Latin inter, meaning "between", and Latin medius, meaning "middle"): between two other structures. For example, the navel is intermediate to the left arm and the contralateral (right) leg. - Parietal (from Latin paries, meaning "wall"): pertaining to the wall of a body cavity. For example, the parietal peritoneum is the lining on the inside of the abdominal cavity. Parietal can also refer specifically to the parietal bone of the skull or associated structures. - Visceral (from Latin viscus, meaning "internal organs"): associated with organs within the body's cavities. For example, the stomach is covered with a lining called the visceral peritoneum. Viscus can also be used to mean "organ". For example, the stomach is a viscus within the abdominal cavity. - Latin convention Commonly when, for example, one anatomical feature is nearer to something than another, one may use an expression such as "nearer the distal end" or "distal to". However, an unambiguous and concise convention is to use the Latin suffix -ad, meaning "towards", or sometimes "to". So for example, "distad" means "in the distal direction", and "distad of the femur" means "beyond the femur in the distal direction". The suffix may be used very widely, as in the following examples: anteriad (towards the anterior), apicad (towards the apex), basad (towards the basal end), caudad, centrad, cephalad (towards the cephalic end), craniad, dextrad, dextrocaudad, dextrocephalad, distad, dorsad, ectad (towards the ectal, or exterior, direction), entad (towards the interior), laterad, mediad, mesad, neurad, orad, posteriad, proximad, rostrad, sinistrad, sinistrocaudad, sinistrocephalad, ventrad. Specific animals and other organisms The large variety of body shapes present in invertebrates presents a difficult problem when attempting to apply standard directional terms. Depending on the organism, some terms are taken by analogy from vertebrate anatomy, and appropriate novel terms are applied as needed. Some such borrowed terms are widely applicable in most invertebrates; for example proximal, literally meaning "near" refers to the part of an appendage nearest to where it joins the body, and distal, literally meaning "standing away from" is used for the part furthest from the point of attachment. In all cases, the usage of terms is dependent on the body plan of the organism. As humans are approximately bilaterally symmetrical organisms, anatomical descriptions usually use the same terms as those for vertebrates and other members of the taxonomic group Bilateria. However, for historical and other reasons, standard human directional terminology has several differences from that used for other bilaterally symmetrical organisms. The terms of zootomy and anatomy came into use at a time when all scientific communication took place in Latin. In their original Latin forms the respective meanings of "anterior" and "posterior" are in front of (or before) and behind (or after), those of "dorsal" and "ventral" are toward the spine and toward the belly, and those of "superior" and "inferior" are above and below. Humans, however, have the rare property of having an upright torso. This makes their anterior/posterior and dorsal/ventral directions the same, and the inferior/superior directions necessary. Most animals, furthermore, are capable of moving relative to their environment. So while "up" might refer to the direction of a standing human's head, the same term ("up") might be used to refer to the direction of the belly of a supine human. It is also necessary to employ some specific anatomical knowledge in order to apply the terminology unambiguously: For example, while the ears would be superior to (above) the shoulders in a human, this fails when describing the armadillo, where the shoulders are above the ears. Thus, in veterinary terminology, the ears would be cranial to (i.e., "toward the head from") the shoulders in the armadillo, the dog, the kangaroo, or any other terrestrial vertebrate, including the human. Likewise, while the belly is considered anterior to (in front of) the back in humans, this terminology fails for the flounder, the armadillo, and the dog. In veterinary terms, the belly would be ventral ("toward the abdomen") in all vertebrates. While it would be possible to introduce a system of axes that is completely consistent between humans and other vertebrates by having two separate pairs of axes, one used exclusively for the head (e.g., anterior/posterior and inferior/superior) and the other exclusively for the torso (e.g., dorsal/ventral and caudal/rostral, or "toward the tail"/"toward the beak"), doing so would require the renaming of very many anatomical structures. Asymmetrical and spherical organisms In organisms with a changeable shape, such as amoeboid organisms, most directional terms are meaningless, since the shape of the organism is not constant and no distinct axes are fixed. Similarly, in spherically symmetrical organisms, there is nothing to distinguish one line through the centre of the organism from any other. An indefinite number of triads of mutually perpendicular axes could be defined, but any such choice of axes would be useless, as nothing would distinguish a chosen triad from any others. In such organisms, only terms such as superficial and deep, or sometimes proximal and distal, are usefully descriptive. In organisms that maintain a constant shape and have one dimension longer than the other, at least two directional terms can be used. The long or longitudinal axis is defined by points at the opposite ends of the organism. Similarly, a perpendicular transverse axis can be defined by points on opposite sides of the organism. There is typically no basis for the definition of a third axis. Usually such organisms are planktonic (free-swimming) protists, and are nearly always viewed on microscope slides, where they appear essentially two-dimensional. In some cases a third axis can be defined, particularly where a non-terminal cytostome or other unique structure is present. Some elongated protists have distinctive ends of the body. In such organisms, the end with a mouth (or equivalent structure, such as the cytostome in Paramecium or Stentor), or the end that usually points in the direction of the organism's locomotion (such as the end with the flagellum in Euglena), is normally designated as the anterior end. The opposite end then becomes the posterior end. Properly, this terminology would apply only to an organism that is always planktonic (not normally attached to a surface), although the term can also be applied to one that is sessile (normally attached to a surface). Organisms that are attached to a substrate, such as sponges, or some animal-like protists also have distinctive ends. The part of the organism attached to the substrate is usually referred to as the basal end (Latin: basis, "support/foundation"), whereas the end furthest from the attachment is referred to as the apical end (Latin: apex, "peak/tip"). Radially symmetrical organisms Radially symmetrical organisms include those in the group Radiata—primarily jellyfish, sea anemones and corals and the comb jellies. Adult echinoderms, such as starfish, sea urchins, sea cucumbers and others are also included, since they are pentaradial, meaning they have five discrete rotational symmetry. Echinoderm larvae are not included, since they are bilaterally symmetrical. Radially symmetrical organisms always have one distinctive axis. Cnidarians (jellyfish, sea anemones and corals) have an incomplete digestive system, meaning that one end of the organism has a mouth, and the opposite end has no opening from the gut (coelenteron). For this reason, the end of the organism with the mouth is referred to as the oral end (Latin: oris, "mouth"), and the opposite surface is the aboral end (Latin: ab-, prefix meaning "away from"). Unlike vertebrates, cnidarians have no other distinctive axes. "Lateral", "dorsal", and "ventral" have no meaning in such organisms, and all can be replaced by the generic term peripheral (Latin: peri-, "around"). Medial can be used, but in the case of radiates indicates the central point, rather than a central axis as in vertebrates. Thus, there are multiple possible radial axes and medio-peripheral (half-) axes. However, it is noteworthy that some biradially symmetrical comb jellies do have distinct "tentacular" and "pharyngeal" axes and are thus anatomically equivalent to bilaterally symmetrical animals. Figure 10: Aurelia aurita, another species of jellyfish, showing multiple radial and medio-peripheral axes. Two specialized terms are useful in describing views of arachnid legs and pedipalps. Prolateral refers to the surface of a leg that is closest to the anterior end of an arachnid's body. Retrolateral refers to the surface of a leg that is closest to the posterior end of an arachnid's body. Because of the unusual nature and positions of the eyes of the Araneae (spiders), and their importance in taxonomy, evolution and anatomy, special terminology with associated abbreviations has become established in arachnology. Araneae normally have eight eyes in four pairs. All the eyes are on the carapace of the prosoma, and their sizes, shapes and locations are characteristic of various spider families and other taxa. In some taxa not all four pairs of eyes are present, the relevant species having only three, two, or one pair of eyes. Some species (mainly troglobites) have no functional eyes at all. In what is seen as the likeliest ancestral arrangement of the eyes of the Araneae, there are two roughly parallel, horizontal, symmetrical, transverse rows of eyes, each containing two symmetrically placed pairs, respectively called: anterior and posterior lateral eyes (ALE) and (PLE); and anterior and posterior median eyes (AME) and (PME). As a rule it is not difficult to guess which eyes are which in a living or preserved specimen, but sometimes it can be. Apart from the fact that in some species one or more pairs may be missing, sometimes eyes from the posterior and anterior rows may be very close to each other, or even fused. Also, either one row or both might be so grossly curved that some of the notionally anterior eyes actually may lie posterior to some of the eyes in the posterior row. In some species the curve is so gross that the eyes apparently are arranged into two anteroposterior parallel rows of eyes. Aspects of spider anatomy; This aspect shows the mainly prolateral surface of the anterior femora, plus the typical horizontal eye pattern of the Sparassidae Typical arrangement of eyes in the Lycosidae, with PME being the largest In the Salticidae the AME are the largest ||This article cites its sources but does not provide page references. (February 2014)| - Wake 1992, p. 1. - Hickman, C. P., Jr., Roberts, L. S. and Larson, A. Animal Diversity. McGraw-Hill 2003 ISBN 0-07-234903-4 - David Craigie (1838). Elements of Anatomy, General, Special, and Comparative. A. and C. Black. - Dorland's Medical Dictionary for Health Consumers. 2007 by Saunders, an imprint of Elsevier, Inc. - "dorsolateral". Merriam-Webster. - Wake 1992, p. 6. - Marieb, E. N. (1995). Human Anatomy and Physiology. Benjamin Cummings. ISBN 0-8053-4281-8. - Tortora, G. J.; Derrickson, B. (2006). Principles of Anatomy and Physiology. John Wiley & Sons. ISBN 0-471-68934-3. - Wake 1992, p. 5. - "What Do Distal and Proximal Mean?". The Survival Doctor. Retrieved 2016-01-07. - GO 2014, "dorsal/ventral axis specification" (GO:0009950). - Miller, S. A. General Zoology Laboratory Manual McGraw-Hill, ISBN 0-07-252837-0 and ISBN 0-07-243559-3 - Ruppert, E. E., Fox, R. S. and Barnes, R. D. Invertebrate Zoology: A Functional Evolutionary Approach. Thomson, Belmont, 2004. ISBN 0-03-025982-7 - Pieter A. Folkens (2000). Human Osteology. Gulf Professional Publishing. pp. 558–. ISBN 978-0-12-746612-5. - Smith, J. B.; Dodson, P. (2003). "A proposal for a standard terminology of anatomical notation and orientation in fossil vertebrate dentitions". Journal of Vertebrate Paleontology 23 (1): 1–12. doi:10.1671/0272-4634(2003)23[1:APFAST]2.0.CO;2. - http://informahealthcare.com/doi/abs/10.3109/02841858909177461?journalCode=ard Evaluation of Three Methods for Measurement of Femoral Neck Anteversion Femoral neck anteversion, definition, measuring methods and errors 1989, Vol. 30, No. 1, Pages 69-73 by Arne Høiseth1†, O. Reikerås1 and E. Fønstelien - Gordh, Gordon; Headrick, David H (2011). A Dictionary of Entomology (2nd ed.). CABI. ISBN 978-1845935429. - Tucker, T. G. (1931). A Concise Etymological Dictionary of Latin. Halle (Saale): Max Niemeyer Verlag. - Valentine, James W. (2004). On the Origin of Phyla. Chicago: University of Chicago Press. ISBN 0-226-84548-6. - Ruppert et al. (2004), p. 184. - Kaston, B. J. (1972). How to Know the Spiders (2nd ed.). Dubuque, Iowa: W. C. Brown Co. p. 19. ISBN 0-697-04899-3.

https://en.wikipedia.org/wiki/Anatomical_terms_of_location

Definitions for pluralˈplʊər əl This page provides all possible meanings and translations of the word plural plural, plural form(adj) the form of a word that is used to denote more than one composed of more than one member, set, or kind grammatical number category referring to two or more items or units : a word in the form in which it potentially refers to something other than one person or thing; and other than two things if the language has a dual form. Consisting of or containing more than one of something. Origin: From plurelle, from plurel, from pluralis, from plus, pluris + -alis. relating to, or containing, more than one; designating two or more; as, a plural word the plural number; that form of a word which expresses or denotes more than one; a word in the plural form Origin: [L. pluralis, from plus, pluris, more; cf. F. pluriel, OF. plurel. See Plus.] The plural, in many languages, is one of the values of the grammatical category of number. Plural forms of nouns typically denote a quantity other than the default quantity represented by a noun, which is generally one. Most commonly, therefore, plurals are used to denote two or more of something, although they may also denote fractional, zero or negative amounts. An example of a plural is the English word cats, which corresponds to the singular cat. Plurality is a linguistic universal, represented variously among the languages as a separate word, an affix, or by other morphological indications such as stress or implicit markers/context. Words of other types, such as verbs, adjectives and pronouns, also frequently have distinct plural forms, which are used in agreement with the number of their associated nouns. Some languages also have a dual or other systems of number categories. However in English and many other languages, singular and plural are the only grammatical numbers, except for possible remnants of the dual in pronouns such as both and either. Sample Sentences & Example Usage The plural of anecdote is data. It is striking how much of the majority's reasoning would apply with equal force to the claim of a fundamental right to plural marriage. Yes there is a willingness to enter Raqqa because we want to cleanse all of Syria. We want to get rid of all of ISIS. We have sworn anywhere that there is ISIS, we must get rid of them for the sake of a free, democratic and plural Syria. Images & Illustrations of plural Translations for plural From our Multilingual Translation Dictionary - جمع, صيغة الجمعArabic - pluralCatalan, Valencian - množné čísloCzech - mehrfach, plural, pluralisch, Mehrzahl, Mehrzahl-German - monikko-, monikko, monikollinenFinnish - meartalWestern Frisian - iolraScottish Gaelic - לשון רביםHebrew - հոգնակի թիվ, հոգնակիArmenian - í fleirtölu, fleirtölu-, marg-, fleir-, fleirtalaIcelandic - 複数形, 複数Japanese - 복수, 복수형Korean - Pluriel, MéizuelLuxembourgish, Letzeburgesch - олон тооMongolian - meervoud, meervoudig, meervoudigeDutch - díkwíjíltʼéegoNavajo, Navaho - liczba mnoga, mnogiPolish - множественный, множественное числоRussian - pluralni, множински, množinski, плуралниSerbo-Croatian - množné číslo, plurálSlovak - flertal, plural, pluralisSwedish - akthari, wingiSwahili - множина, чисельнийUkrainian - plunum, plunumikVolapük Get even more translations for plural » Find a translation for the plural definition in other languages: Select another language:

http://www.definitions.net/definition/plural

Interactive Modeling is a straightforward, quickly paced, seven-step process that’s effective for teaching children any academic or social skill, routine, or procedure that you want them to do in a specific way (whether for safety, efficiency, or other reasons). One of the essential practices of the Responsive Classroom approach to teaching elementary school children, Interactive Modeling can be used by any adult anywhere in school at any time of year. How does Interactive Modeling differ from traditional modeling? In traditional modeling, the teacher shows children how to do a skill, routine, or procedure, tells them what to notice, and expects that they will learn it immediately. Interactive Modeling also shows children how to do skills, routines, or procedures, but it goes well beyond that basic step. Students also: - Learn exactly why the skill, routine, or procedure is important to their learning and the respectful, smooth functioning of the classroom. - Are asked what they noticed about the teacher’s modeling (rather than told by their teacher what to notice). - See a few classmates additionally model the routine or procedure after the teacher’s initial modeling. - Practice the routine or procedure right away. - Receive immediate feedback and coaching from their teacher while they practice. Why is Interactive Modeling more effective than traditional modeling? The distinctive steps of Interactive Modeling incorporate key elements of effective teaching: modeling positive behaviors, engaging students in active learning, and immediately assessing their understanding. Research shows that when we teach in this way, children achieve greater, faster, and longer-lasting success in meeting expectations and mastering skills. With Interactive Modeling, children create clear, positive mental images of what is expected of them. They do the noticing themselves, which builds up their powers of observation and their analysis and communication skills. In addition, because they get immediate practice, they gain quicker expertise and stronger mastery of the procedure or skill being taught. What are the seven steps of Interactive Modeling? - Briefly state what you will model, and why. - Model the behavior exactly as you expect students to do it (the right way, not the wrong way, and without describing what you’re doing unless you need to “show” a thinking process). - Ask students what they noticed. (You may need to do some prompting, but children soon notice every little detail, especially as they gain expertise with this practice.) - Invite one or more students to model the same way you did. - Again, ask students what they noticed the modelers doing. - Have all students model while you observe and coach them. - Provide feedback, naming specific, positive actions you notice and redirecting respectfully but clearly when students go off track. What can I teach with Interactive Modeling? Here are just a few examples: Academic & Social Skills - Listening and responding to questions - Working with a partner or small group - Using technology and other resources - Taking part in a whole-group discussion - Test-prep procedures Procedures & Routines - Arrival and dismissal routines - Cleaning up - Lunch, recess, and bathroom routines - Schoolwide assembly procedures - Transitions from one classroom/activity to another How long does Interactive Modeling take? An Interactive Modeling lesson to demonstrate lining up, for example, may take only three or four minutes. A more involved lesson, such as teaching children how to partner chat, might take twenty minutes. This modest investment saves you time in the long run. That’s because children gain mastery more quickly and are thus able to spend much more time on task. You’ll have less confusion in the classroom and fewer interruptions because children will not need to ask you or peers over and over what to do. As a result, you’ll have more time for teaching—and children will have more time to complete their work and to learn. How much time could I gain with Interactive Modeling? A little time spent on teaching students exactly how you want them to do things will pay big dividends throughout the year. Say you lose a few minutes every hour to repeating instructions and dealing with interruptions. That can add up to twenty or thirty minutes of lost instructional and work time each day—2½ hours each week. That’s 100 lost hours every school year! Think of what you could do with that time. Think of what your students could do. What does Interactive Modeling look and sound like in action? In a fifth grade classroom, Mrs. K wants the students to understand how to work productively during independent work time. She teaches this Interactive Modeling lesson: - Say what you will model and why. Mrs. K: “Our goal is for everyone to do high-quality work during independent work time. Watch how Carlos and I work hard on our assignment and let others do the same.” - Model the behavior. Mrs. K and Carlos (coached in advance) demonstrate how to work on a research assignment at the same table. They work quietly, but to show that it’s okay to talk, they each exchange one fact from their research, briefly and in low voices. Then they get right back to work. - Ask students what they noticed. Mrs. K: “What did you notice?” Her students note the key elements of the demonstration, such as how Carlos and Mrs. K stayed in their seats, worked quietly, and talked in low voices for only a short time. Mrs. K prompts students to name any key behaviors they missed. For example: “What did we do with our papers and other materials?” - Invite one or more students to model. Mrs. K chooses four more students, who demonstrate how to work independently at the same table just as she and Carlos did. - Again, ask students what they noticed. Mrs. K: “What did you notice this time?” Her students point out the key elements, just as they did in Step 3, helping to reinforce these behaviors for themselves. Again, she prompts them if they miss any key behaviors. - Have all students practice. Mrs. K gives all her students a short survey to work on so she can observe and coach them. - Provide feedback. Mrs. K: “I see everyone focused on the survey, working quietly. That kind of focus will help you and your classmates complete your assignments and learn a lot this year.”

http://www.responsiveclassroom.org/what-is-interactive-modeling/

There are two types of atomic bonds - ionic bonds and covalent bonds. They differ in their structure and properties. Covalent bonds consist of pairs of electrons shared by two atoms, and bind the atoms in a fixed orientation. Relatively high energies are required to break them (50 - 200 kcal/mol). Whether two atoms can form a covalent bond depends upon their electronegativity i.e. the power of an atom in a molecule to attract electrons to itself. If two atoms differ considerably in their electronegativity - as sodium and chloride do - then one of the atoms will lose its electron to the other atom. This results in a positively charged ion (cation) and negatively charged ion (anion). The bond between these two ions is called an ionic bond. |Covalent Bonds||Ionic Bonds| |Formation||A covalent bond is formed between two non-metals that have similar electronegativities. Neither atom is "strong" enough to attract electrons from the other. For stabilization, they share their electrons from outer molecular orbit with others.||An ionic bond is formed between a metal and a non-metal. Non-metals(-ve ion) are "stronger" than the metal(+ve ion) and can get electrons very easily from the metal. These two opposite ions attract each other and form the ionic bond.| |Shape||Definite shape||No definite shape| |What is it?||Covalent bonding is a form of chemical bonding between two non metallic atoms which is characterized by the sharing of pairs of electrons between atoms and other covalent bonds.||Ionic bond, also known as electrovalent bond is a type of bond formed from the electrostatic attraction between oppositely charged ions in a chemical compound. These kinds of bonds occur mainly between a metallic and a non metallic atom.| |Examples||Methane (CH4), Hydro Chloric acid (HCl)||Sodium chloride (NaCl), Sulphuric Acid (H2SO4 )| |Occurs between||Two non-metals||One metal and one non-metal| |State at room temperature||Liquid or gaseous||Solid| About Covalent and Ionic Bonds The covalent bond is formed when two atoms are able to share electrons whereas the ionic bond is formed when the "sharing" is so unequal that an electron from atom A is completely lost to atom B, resulting in a pair of ions. Each atom consists of protons, neutrons and electrons. At the centre of the atom, neutrons and protons stay together. But electrons revolve in orbit around the center. Each of these molecular orbits can have a certain number of electrons to form a stable atom. But apart from Inert gas, this configuration is not present with most of the atoms. So to stabilize the atom, each atom shares half of its electrons. Covalent bonding is a form of chemical bonding between two non metallic atoms which is characterized by the sharing of pairs of electrons between atoms and other covalent bonds. Ionic bond, also known as electrovalent bond, is a type of bond formed from the electrostatic attraction between oppositely charged ions in a chemical compound. This kind of bonds occurs mainly between a metallic and a non metallic atom. Formation and examples Covalent bonds are formed as a result of the sharing of one or more pairs of bonding electrons. The electro negativities (electron attracting ability) of the two bonded atoms are either equal or the difference is no greater than 1.7. As long as the electro-negativity difference is no greater than 1.7, the atoms can only share the bonding electrons. For example, let us consider a Methane molecule i.e.CH4. Carbon has 6 electrons and its electronic configuration is 1s22s22p2, i.e. it has 4 electrons in its outer orbit. According to the Octate rule ( It states that atoms tend to gain, lose, or share electrons so that each atom has full outermost energy level which is typically 8 electrons.), to be in a stable state, it needs 4 more electrons. So it forms covalent bond with Hydrogen (1s1), and by sharing electrons with hydrogen it forms Methane or CH4. If the electro-negativity difference is greater than 1.7 then the higher electronegative atom has an electron attracting ability which is large enough to force the transfer of electrons from the lesser electronegative atom. This causes the formation of ionic bonds. For example, in common table salt (NaCl) the individual atoms are sodium and chlorine. Chlorine has seven valence electrons in its outer orbit but to be in a stable condition, it needs eight electrons in outer orbit. On the other hand, Sodium has one valence electron and it also needs eight electrons. Since chlorine has a high electro-negativity, 3.16 compared to sodium’s 0.9, (so the difference between their electro-negativity is more than 1.7) chlorine can easily attract sodium's one valence electron. In this manner they form an Ionic bond, and share each other’s electrons and both will have 8 electrons in their outer shell. Characteristics of the bonds Covalent bonds have a definite and predictable shape and have low melting and boiling points. They can be easily broken into its primary structure as the atoms are close by to share the electrons. These are mostly gaseous and even a slight negative or positive charge at opposite ends of a covalent bond gives them molecular polarity. Ionic bonds normally form crystalline atoms and have higher melting points and boiling points compared to covalent compounds. These conduct electricity in molten or solution state and they are extremely polar bonds. Most of them are soluble in water but insoluble in non-polar solvents. They require much more energy than covalent bond to break the bond between them. The reason for the difference in the melting and boiling points for ionic and covalent bonds can be illustrated through an example of NaCl (ionic bond) and Cl2 (covalent bond). This example can be found at Cartage.org. - Wikipedia: Double bond - Covalent Bonds - The City University of New York - Chemical Bonding - Georgia State University - Covalent and Ionic Bonds - Access Excellence - Electron Sharing and Covalent Bonds - University of Oxford - Wikipedia: Molecular orbital diagram - Wikipedia: Electron configuration - Ionic Bond - Encyclopedia Britannica

http://www.diffen.com/difference/Covalent_Bonds_vs_Ionic_Bonds

Anna Mika, National Geographic Education Video length: 30 minutes.Learn more about Teaching Climate Literacy and Energy Awareness» See how this Activity supports the Next Generation Science Standards» Middle School: 1 Disciplinary Core Idea, 4 Cross Cutting Concepts, 2 Science and Engineering Practices About Teaching Climate Literacy Other materials addressing 2a Other materials addressing 4b Other materials addressing 4d Excellence in Environmental Education Guidelines Other materials addressing: A) Processes that shape the Earth. Notes From Our Reviewers The CLEAN collection is hand-picked and rigorously reviewed for scientific accuracy and classroom effectiveness. Read what our review team had to say about this resource below or learn more about how CLEAN reviews teaching materials Teaching Tips | Science | Pedagogy | - Educators are encouraged to start the activity by activating students' prior knowledge about extreme weather on Earth. - As an extension to this activity, educator could encourage students to investigate what constitutes extreme weather. In some areas, certain weather-related events may not be classified as extreme. About the Science - Students investigate types of extreme weather and their contributing factors, and then explore similarities and differences between weather and climate. - Comment from expert scientist: Good overview and use of imagery. Language is appropriate. About the Pedagogy - Video and brilliant photography archived on the National Geographic Society's Website make this lesson stimulating; otherwise, it is a standard paper-and-pencil activity. - Text is provided underneath photographs, explaining the images. Technical Details/Ease of Use - Educator is given clear instruction on how to introduce and guide this activity. - All resources needed for the activity are provided on the resource Webpage but cannot be downloaded at this time. - Video is challenging to find on their new site http://video.nationalgeographic.com/video/science/earth-sci/weather-101-sci/; no transcript on the video. Next Generation Science Standards See how this Activity supports: Disciplinary Core Ideas: 1 MS-ESS2.D1:Weather and climate are influenced by interactions involving sunlight, the ocean, the atmosphere, ice, landforms, and living things. These interactions vary with latitude, altitude, and local and regional geography, all of which can affect oceanic and atmospheric flow patterns. Cross Cutting Concepts: 4 MS-C5.2: Within a natural or designed system, the transfer of energy drives the motion and/or cycling of matter. MS-C7.1: Explanations of stability and change in natural or designed systems can be constructed by examining the changes over time and forces at different scales, including the atomic scale. MS-C1.4:Graphs, charts, and images can be used to identify patterns in data. MS-C2.2:Cause and effect relationships may be used to predict phenomena in natural or designed systems. Science and Engineering Practices: 2 MS-P4.3: Distinguish between causal and correlational relationships in data. MS-P6.1:Construct an explanation that includes qualitative or quantitative relationships between variables that predict(s) and/or describe(s) phenomena.

http://cleanet.org/resources/43389.html

Word problems in mathematics often pose a challenge because they require that students read and comprehend the text of the problem, identify the question that needs to be answered, and finally create and solve a numerical equation. Many ELLs may have difficulty reading and understanding the written content in a word problem. If a student is learning English as a second language, he might not yet know key terminology needed to solve the equation. In other words, ELLs who have had formal education in their home countries generally do not have mathematical difficulties; hence, their struggles begin when they encounter word problems in a second language that they have not yet mastered (Bernardo, 2005). For this reason it is recommended that students learn key terminology prior to attempting to solve mathematical word problems. Once English language learners know the key terminology used in mathematical word problems, it will be easier to learn how to write numerical equations. It is important for teachers to provide ELLs with opportunities to learn and practice key vocabulary words. While key words are very important, they are only part of the process. Understanding the language in word problems is critical for all students. They need to know the meaning of words. But because words are often used differently and problems are set up differently, there are some cautionary messages. Here is an example of problem that uses "fewer than" to set up a subtraction equation. Maria has 24 marbles which is 8 fewer than Paolo has. How many marbles does Paolo have? If we were to only focus on using key words, "fewer than" is a signal to pick out the numbers and subtract. The student may immediately make the conclusion that the answer is 16, but that is not what the problem is asking, and the child would be wrong. (The correct answer, by the way, is 32). What research has found is that if we ask students to only rely on knowing that certain key words signal specific operations, we can actually lead them away from trying to understand the problems. They will tend to look only for those words and whatever numbers are in the problem, even if they are not relevant to the answer. This will not help them be mathematically proficient later, even when they are proficient with English. Although the finding on key words was done with regular students, the consequences for ELL students of relying on them is the same. They would not be able to solve the problem above. However, if teachers follow the suggested process of reading a problem several times (at lower as well as upper grades) and discussing what it means, students will understand. Another good tool is to teach them to draw or model the problems. To illustrate the problem above, you could state: "Here's Maria's 24." Then, draw 24 units, figures, shapes, etc. to represent 24. "Here's Paolo's; he has more because Maria has fewer than he does". Draw 24 units, figures, shapes, etc. to represent 24 and add 8 more. "So Paolo's has to come to more than 24. How many more? 8. So what is Paolo's total?" The difference is between knowing the meaning of the words "fewer than" and using "fewer than" as a key to an operation. We want students to know the meaning of the words, but also to see them in the context of the whole problem. - How many students brought their homework today? - How many more children brought their homework yesterday? - We had 8 markers on the board, but now we only have 3. How many did we take away? - How many animals are there in this magazine? How many are mammals? How many are birds? (introduction to fractions and percentages) Continue to use key terminology daily and put it in context (e.g., less than, more than, difference, times, each, etc.). Show students how easy it might be to misunderstand the problem. - Read word problems slowly and carefully several times so that all students comprehend. - If possible, break up the problem into smaller segments. - Allow students to act out the word problems to better comprehend what they are being asked to solve. - Provide manipulatives to help students visualize the problem. - Take field or walking trips to figure out distances, speed, area covered, etc. - Ask students to do surveys, interviews, hands-on research in real-world situations to figure out percentages, differences, and higher-order math skills. - Allow students to make drawings or diagrams to help them understand problems. This article provides an overview of the challenges ELLs face in their content-area classes, such as math, science, and social studies. Understanding these challenges will help both ELL and content-area teachers adapt instruction for ELL students. The National Center for Research on the Educational Achievement and Teaching of English Language Learners (CREATE) is a research program designed to improve educational outcomes for ELLs by using a combination of strategies focused on readers in Grades 4-8 and teacher professional development. The CREATE Web site includes webcast seminars hosted by CREATE researchers, and listings of CREATE conferences and presentations around the country. TeacherVision offers lessons, printables, and quizzes to support math word problem instruction for grades K-12. This is a great site for teachers in the elementary levels, as it provides a list of keywords you can teach your ELLs to look for as they read word problems. Also included are useful ideas and tricks to better prepare students to understand written math problems.

http://www.colorincolorado.org/article/reading-and-understanding-written-math-problems

||It has been suggested that Process (science) be merged into this article. (Discuss) Proposed since December 2015.| The scientific method is a body of techniques for investigating phenomena, acquiring new knowledge, or correcting and integrating previous knowledge. To be termed scientific, a method of inquiry is commonly based on empirical or measurable evidence subject to specific principles of reasoning. The Oxford English Dictionary defines the scientific method as "a method or procedure that has characterized natural science since the 17th century, consisting in systematic observation, measurement, and experiment, and the formulation, testing, and modification of hypotheses." The scientific method is an ongoing process, which usually begins with observations about the natural world. Human beings are naturally inquisitive, so they often come up with questions about things they see or hear and often develop ideas (hypotheses) about why things are the way they are. The best hypotheses lead to predictions that can be tested in various ways, including making further observations about nature. In general, the strongest tests of hypotheses come from carefully controlled and replicated experiments that gather empirical data. Depending on how well the tests match the predictions, the original hypothesis may require refinement, alteration, expansion or even rejection. If a particular hypothesis becomes very well supported a general theory may be developed. Although procedures vary from one field of inquiry to another, identifiable features are frequently shared in common between them. The overall process of the scientific method involves making conjectures (hypotheses), deriving predictions from them as logical consequences, and then carrying out experiments based on those predictions. A hypothesis is a conjecture, based on knowledge obtained while formulating the question. The hypothesis might be very specific or it might be broad. Scientists then test hypotheses by conducting experiments. Under modern interpretations, a scientific hypothesis must be falsifiable, implying that it is possible to identify a possible outcome of an experiment that conflicts with predictions deduced from the hypothesis; otherwise, the hypothesis cannot be meaningfully tested. The purpose of an experiment is to determine whether observations agree with or conflict with the predictions derived from a hypothesis. Experiments can take place in a college lab, on a kitchen table, at CERN's Large Hadron Collider, at the bottom of an ocean, on Mars, and so on. There are difficulties in a formulaic statement of method, however. Though the scientific method is often presented as a fixed sequence of steps, it represents rather a set of general principles. Not all steps take place in every scientific inquiry (or to the same degree), and are not always in the same order. - 1 Overview - 2 Scientific inquiry - 3 Elements of the scientific method - 3.1 Characterizations - 3.2 Hypothesis development - 3.3 Predictions from the hypothesis - 3.4 Experiments - 3.5 Evaluation and improvement - 3.6 Confirmation - 4 Models of scientific inquiry - 5 Communication and community - 6 Philosophy and sociology of science - 7 History - 8 Relationship with mathematics - 9 Relationship with statistics - 10 See also - 11 Notes - 12 References - 13 Further reading - 14 External links - The DNA example below is a synopsis of this method The scientific method is the process by which science is carried out. As in other areas of inquiry, science (through the scientific method) can build on previous knowledge and develop a more sophisticated understanding of its topics of study over time. This model can be seen to underlay the scientific revolution. One thousand years ago, Alhazen argued the importance of forming questions and subsequently testing them, an approach which was advocated by Galileo in 1638 with the publication of Two New Sciences. The current method is based on a hypothetico-deductive model formulated in the 20th century, although it has undergone significant revision since first proposed (for a more formal discussion, see below). The overall process involves making conjectures (hypotheses), deriving predictions from them as logical consequences, and then carrying out experiments based on those predictions to determine whether the original conjecture was correct. There are difficulties in a formulaic statement of method, however. Though the scientific method is often presented as a fixed sequence of steps, they are better considered as general principles. Not all steps take place in every scientific inquiry (or to the same degree), and are not always in the same order. As noted by William Whewell (1794–1866), "invention, sagacity, [and] genius" are required at every step. Formulation of a question The question can refer to the explanation of a specific observation, as in "Why is the sky blue?", but can also be open-ended, as in "How can I design a drug to cure this particular disease?" This stage frequently involves looking up and evaluating evidence from previous experiments, personal scientific observations or assertions, and/or the work of other scientists. If the answer is already known, a different question that builds on the previous evidence can be posed. When applying the scientific method to scientific research, determining a good question can be very difficult and affects the final outcome of the investigation. A hypothesis is a conjecture, based on knowledge obtained while formulating the question, that may explain the observed behavior of a part of our universe. The hypothesis might be very specific, e.g., Einstein's equivalence principle or Francis Crick's "DNA makes RNA makes protein", or it might be broad, e.g., unknown species of life dwell in the unexplored depths of the oceans. A statistical hypothesis is a conjecture about some population. For example, the population might be people with a particular disease. The conjecture might be that a new drug will cure the disease in some of those people. Terms commonly associated with statistical hypotheses are null hypothesis and alternative hypothesis. A null hypothesis is the conjecture that the statistical hypothesis is false, e.g., that the new drug does nothing and that any cures are due to chance effects. Researchers normally want to show that the null hypothesis is false. The alternative hypothesis is the desired outcome, e.g., that the drug does better than chance. A final point: a scientific hypothesis must be falsifiable, meaning that one can identify a possible outcome of an experiment that conflicts with predictions deduced from the hypothesis; otherwise, it cannot be meaningfully tested. This step involves determining the logical consequences of the hypothesis. One or more predictions are then selected for further testing. The more unlikely that a prediction would be correct simply by coincidence, then the more convincing it would be if the prediction were fulfilled; evidence is also stronger if the answer to the prediction is not already known, due to the effects of hindsight bias (see also postdiction). Ideally, the prediction must also distinguish the hypothesis from likely alternatives; if two hypotheses make the same prediction, observing the prediction to be correct is not evidence for either one over the other. (These statements about the relative strength of evidence can be mathematically derived using Bayes' Theorem). This is an investigation of whether the real world behaves as predicted by the hypothesis. Scientists (and other people) test hypotheses by conducting experiments. The purpose of an experiment is to determine whether observations of the real world agree with or conflict with the predictions derived from a hypothesis. If they agree, confidence in the hypothesis increases; otherwise, it decreases. Agreement does not assure that the hypothesis is true; future experiments may reveal problems. Karl Popper advised scientists to try to falsify hypotheses, i.e., to search for and test those experiments that seem most doubtful. Large numbers of successful confirmations are not convincing if they arise from experiments that avoid risk. Experiments should be designed to minimize possible errors, especially through the use of appropriate scientific controls. For example, tests of medical treatments are commonly run as double-blind tests. Test personnel, who might unwittingly reveal to test subjects which samples are the desired test drugs and which are placebos, are kept ignorant of which are which. Such hints can bias the responses of the test subjects. Furthermore, failure of an experiment does not necessarily mean the hypothesis is false. Experiments always depend on several hypotheses, e.g., that the test equipment is working properly, and a failure may be a failure of one of the auxiliary hypotheses. (See the Duhem–Quine thesis.) Experiments can be conducted in a college lab, on a kitchen table, at CERN's Large Hadron Collider, at the bottom of an ocean, on Mars (using one of the working rovers), and so on. Astronomers do experiments, searching for planets around distant stars. Finally, most individual experiments address highly specific topics for reasons of practicality. As a result, evidence about broader topics is usually accumulated gradually. This involves determining what the results of the experiment show and deciding on the next actions to take. The predictions of the hypothesis are compared to those of the null hypothesis, to determine which is better able to explain the data. In cases where an experiment is repeated many times, a statistical analysis such as a chi-squared test may be required. If the evidence has falsified the hypothesis, a new hypothesis is required; if the experiment supports the hypothesis but the evidence is not strong enough for high confidence, other predictions from the hypothesis must be tested. Once a hypothesis is strongly supported by evidence, a new question can be asked to provide further insight on the same topic. Evidence from other scientists and experience are frequently incorporated at any stage in the process. Depending on the complexity of the experiment, many iterations may be required to gather sufficient evidence to answer a question with confidence, or to build up many answers to highly specific questions in order to answer a single broader question. |The basic elements of the scientific method are illustrated by the following example from the discovery of the structure of DNA: The discovery became the starting point for many further studies involving the genetic material, such as the field of molecular genetics, and it was awarded the Nobel Prize in 1962. Each step of the example is examined in more detail later in the article. The scientific method also includes other components required even when all the iterations of the steps above have been completed: If an experiment cannot be repeated to produce the same results, this implies that the original results might have been in error. As a result, it is common for a single experiment to be performed multiple times, especially when there are uncontrolled variables or other indications of experimental error. For significant or surprising results, other scientists may also attempt to replicate the results for themselves, especially if those results would be important to their own work. The process of peer review involves evaluation of the experiment by experts, who typically give their opinions anonymously. Some journals request that the experimenter provide lists of possible peer reviewers, especially if the field is highly specialized. Peer review does not certify correctness of the results, only that, in the opinion of the reviewer, the experiments themselves were sound (based on the description supplied by the experimenter). If the work passes peer review, which occasionally may require new experiments requested by the reviewers, it will be published in a peer-reviewed scientific journal. The specific journal that publishes the results indicates the perceived quality of the work. Data recording and sharing Scientists typically are careful in recording their data, a requirement promoted by Ludwik Fleck (1896–1961) and others. Though not typically required, they might be requested to supply this data to other scientists who wish to replicate their original results (or parts of their original results), extending to the sharing of any experimental samples that may be difficult to obtain. Scientific inquiry generally aims to obtain knowledge in the form of testable explanations that can be used to predict the results of future experiments. This allows scientists to gain a better understanding of the topic being studied, and later be able to use that understanding to intervene in its causal mechanisms (such as to cure disease). The better an explanation is at making predictions, the more useful it frequently can be, and the more likely it is to continue explaining a body of evidence better than its alternatives. The most successful explanations, which explain and make accurate predictions in a wide range of circumstances, are often called scientific theories. Most experimental results do not produce large changes in human understanding; improvements in theoretical scientific understanding is typically the result of a gradual process of development over time, sometimes across different domains of science. Scientific models vary in the extent to which they have been experimentally tested and for how long, and in their acceptance in the scientific community. In general, explanations become accepted over time as evidence accumulates on a given topic, and the explanation in question is more powerful than its alternatives at explaining the evidence. Often the explanations are altered over time, or explanations are combined to produce new explanations. Properties of scientific inquiry Scientific knowledge is closely tied to empirical findings, and can remain subject to falsification if new experimental observation incompatible with it is found. That is, no theory can ever be considered final, since new problematic evidence might be discovered. If such evidence is found, a new theory may be proposed, or (more commonly) it is found that modifications to the previous theory are sufficient to explain the new evidence. The strength of a theory can be argued to be related to how long it has persisted without major alteration to its core principles. Theories can also subject to subsumption by other theories. For example, thousands of years of scientific observations of the planets were explained almost perfectly by Newton's laws. However, these laws were then determined to be special cases of a more general theory (relativity), which explained both the (previously unexplained) exceptions to Newton's laws and predicting and explaining other observations such as the deflection of light by gravity. Thus, in certain cases independent, unconnected, scientific observations can be connected to each other, unified by principles of increasing explanatory power. Since new theories might be more comprehensive than what preceded them, and thus be able to explain more than previous ones, successor theories might be able to meet a higher standard by explaining a larger body of observations than their predecessors. For example, the theory of evolution explains the diversity of life on Earth, how species adapt to their environments, and many other patterns observed in the natural world; its most recent major modification was unification with genetics to form the modern evolutionary synthesis. In subsequent modifications, it has also subsumed aspects of many other fields such as biochemistry and molecular biology. Beliefs and biases Scientific methodology often directs that hypotheses be tested in controlled conditions wherever possible. This is frequently possible in certain areas, such as in the biological sciences, and more difficult in other areas, such as in astronomy. The practice of experimental control and reproducibility can have the effect of diminishing the potentially harmful effects of circumstance, and to a degree, personal bias. For example, pre-existing beliefs can alter the interpretation of results, as in confirmation bias; this is a heuristic that leads a person with a particular belief to see things as reinforcing their belief, even if another observer might disagree (in other words, people tend to observe what they expect to observe). A historical example is the belief that the legs of a galloping horse are splayed at the point when none of the horse's legs touches the ground, to the point of this image being included in paintings by its supporters. However, the first stop-action pictures of a horse's gallop by Eadweard Muybridge showed this to be false, and that the legs are instead gathered together. Another important human bias that plays a role is a preference for new, surprising statements (see appeal to novelty), which can result in a search for evidence that the new is true. In contrast to this standard in the scientific method, poorly attested beliefs can be believed and acted upon via a less rigorous heuristic, sometimes taking advantage of the narrative fallacy that when narrative is constructed its elements become easier to believe. Sometimes, these have their elements assumed a priori, or contain some other logical or methodological flaw in the process that ultimately produced them. Elements of the scientific method There are different ways of outlining the basic method used for scientific inquiry. The scientific community and philosophers of science generally agree on the following classification of method components. These methodological elements and organization of procedures tend to be more characteristic of natural sciences than social sciences. Nonetheless, the cycle of formulating hypotheses, testing and analyzing the results, and formulating new hypotheses, will resemble the cycle described below. - Four essential elements of the scientific method are iterations, recursions, interleavings, or orderings of the following: - Characterizations (observations, definitions, and measurements of the subject of inquiry) - Hypotheses (theoretical, hypothetical explanations of observations and measurements of the subject) - Predictions (reasoning including deductive reasoning from the hypothesis or theory) - Experiments (tests of all of the above) Each element of the scientific method is subject to peer review for possible mistakes. These activities do not describe all that scientists do (see below) but apply mostly to experimental sciences (e.g., physics, chemistry, and biology). The elements above are often taught in the educational system as "the scientific method". The scientific method is not a single recipe: it requires intelligence, imagination, and creativity. In this sense, it is not a mindless set of standards and procedures to follow, but is rather an ongoing cycle, constantly developing more useful, accurate and comprehensive models and methods. For example, when Einstein developed the Special and General Theories of Relativity, he did not in any way refute or discount Newton's Principia. On the contrary, if the astronomically large, the vanishingly small, and the extremely fast are removed from Einstein's theories – all phenomena Newton could not have observed – Newton's equations are what remain. Einstein's theories are expansions and refinements of Newton's theories and, thus, increase our confidence in Newton's work. A linearized, pragmatic scheme of the four points above is sometimes offered as a guideline for proceeding: - Define a question - Gather information and resources (observe) - Form an explanatory hypothesis - Test the hypothesis by performing an experiment and collecting data in a reproducible manner - Analyze the data - Interpret the data and draw conclusions that serve as a starting point for new hypothesis - Publish results - Retest (frequently done by other scientists) The iterative cycle inherent in this step-by-step method goes from point 3 to 6 back to 3 again. While this schema outlines a typical hypothesis/testing method, it should also be noted that a number of philosophers, historians and sociologists of science (perhaps most notably Paul Feyerabend) claim that such descriptions of scientific method have little relation to the ways that science is actually practiced. The scientific method depends upon increasingly sophisticated characterizations of the subjects of investigation. (The subjects can also be called unsolved problems or the unknowns.) For example, Benjamin Franklin conjectured, correctly, that St. Elmo's fire was electrical in nature, but it has taken a long series of experiments and theoretical changes to establish this. While seeking the pertinent properties of the subjects, careful thought may also entail some definitions and observations; the observations often demand careful measurements and/or counting. The systematic, careful collection of measurements or counts of relevant quantities is often the critical difference between pseudo-sciences, such as alchemy, and science, such as chemistry or biology. Scientific measurements are usually tabulated, graphed, or mapped, and statistical manipulations, such as correlation and regression, performed on them. The measurements might be made in a controlled setting, such as a laboratory, or made on more or less inaccessible or unmanipulatable objects such as stars or human populations. The measurements often require specialized scientific instruments such as thermometers, spectroscopes, particle accelerators, or voltmeters, and the progress of a scientific field is usually intimately tied to their invention and improvement. Measurements in scientific work are also usually accompanied by estimates of their uncertainty. The uncertainty is often estimated by making repeated measurements of the desired quantity. Uncertainties may also be calculated by consideration of the uncertainties of the individual underlying quantities used. Counts of things, such as the number of people in a nation at a particular time, may also have an uncertainty due to data collection limitations. Or counts may represent a sample of desired quantities, with an uncertainty that depends upon the sampling method used and the number of samples taken. Measurements demand the use of operational definitions of relevant quantities. That is, a scientific quantity is described or defined by how it is measured, as opposed to some more vague, inexact or "idealized" definition. For example, electric current, measured in amperes, may be operationally defined in terms of the mass of silver deposited in a certain time on an electrode in an electrochemical device that is described in some detail. The operational definition of a thing often relies on comparisons with standards: the operational definition of "mass" ultimately relies on the use of an artifact, such as a particular kilogram of platinum-iridium kept in a laboratory in France. The scientific definition of a term sometimes differs substantially from its natural language usage. For example, mass and weight overlap in meaning in common discourse, but have distinct meanings in mechanics. Scientific quantities are often characterized by their units of measure which can later be described in terms of conventional physical units when communicating the work. New theories are sometimes developed after realizing certain terms have not previously been sufficiently clearly defined. For example, Albert Einstein's first paper on relativity begins by defining simultaneity and the means for determining length. These ideas were skipped over by Isaac Newton with, "I do not define time, space, place and motion, as being well known to all." Einstein's paper then demonstrates that they (viz., absolute time and length independent of motion) were approximations. Francis Crick cautions us that when characterizing a subject, however, it can be premature to define something when it remains ill-understood. In Crick's study of consciousness, he actually found it easier to study awareness in the visual system, rather than to study free will, for example. His cautionary example was the gene; the gene was much more poorly understood before Watson and Crick's pioneering discovery of the structure of DNA; it would have been counterproductive to spend much time on the definition of the gene, before them. The history of the discovery of the structure of DNA is a classic example of the elements of the scientific method: in 1950 it was known that genetic inheritance had a mathematical description, starting with the studies of Gregor Mendel, and that DNA contained genetic information (Oswald Avery's transforming principle). But the mechanism of storing genetic information (i.e., genes) in DNA was unclear. Researchers in Bragg's laboratory at Cambridge University made X-ray diffraction pictures of various molecules, starting with crystals of salt, and proceeding to more complicated substances. Using clues painstakingly assembled over decades, beginning with its chemical composition, it was determined that it should be possible to characterize the physical structure of DNA, and the X-ray images would be the vehicle. ..2. DNA-hypotheses Another example: precession of Mercury The characterization element can require extended and extensive study, even centuries. It took thousands of years of measurements, from the Chaldean, Indian, Persian, Greek, Arabic and European astronomers, to fully record the motion of planet Earth. Newton was able to include those measurements into consequences of his laws of motion. But the perihelion of the planet Mercury's orbit exhibits a precession that cannot be fully explained by Newton's laws of motion (see diagram to the right), as Leverrier pointed out in 1859. The observed difference for Mercury's precession between Newtonian theory and observation was one of the things that occurred to Einstein as a possible early test of his theory of General Relativity. His relativistic calculations matched observation much more closely than did Newtonian theory. The difference is approximately 43 arc-seconds per century. A hypothesis is a suggested explanation of a phenomenon, or alternately a reasoned proposal suggesting a possible correlation between or among a set of phenomena. Normally hypotheses have the form of a mathematical model. Sometimes, but not always, they can also be formulated as existential statements, stating that some particular instance of the phenomenon being studied has some characteristic and causal explanations, which have the general form of universal statements, stating that every instance of the phenomenon has a particular characteristic. Scientists are free to use whatever resources they have – their own creativity, ideas from other fields, inductive reasoning, Bayesian inference, and so on – to imagine possible explanations for a phenomenon under study. Charles Sanders Peirce, borrowing a page from Aristotle (Prior Analytics, 2.25) described the incipient stages of inquiry, instigated by the "irritation of doubt" to venture a plausible guess, as abductive reasoning. The history of science is filled with stories of scientists claiming a "flash of inspiration", or a hunch, which then motivated them to look for evidence to support or refute their idea. Michael Polanyi made such creativity the centerpiece of his discussion of methodology. William Glen observes that - the success of a hypothesis, or its service to science, lies not simply in its perceived "truth", or power to displace, subsume or reduce a predecessor idea, but perhaps more in its ability to stimulate the research that will illuminate ... bald suppositions and areas of vagueness. In general scientists tend to look for theories that are "elegant" or "beautiful". In contrast to the usual English use of these terms, they here refer to a theory in accordance with the known facts, which is nevertheless relatively simple and easy to handle. Occam's Razor serves as a rule of thumb for choosing the most desirable amongst a group of equally explanatory hypotheses. Linus Pauling proposed that DNA might be a triple helix. This hypothesis was also considered by Francis Crick and James D. Watson but discarded. When Watson and Crick learned of Pauling's hypothesis, they understood from existing data that Pauling was wrong and that Pauling would soon admit his difficulties with that structure. So, the race was on to figure out the correct structure (except that Pauling did not realize at the time that he was in a race) ..3. DNA-predictions Predictions from the hypothesis Any useful hypothesis will enable predictions, by reasoning including deductive reasoning. It might predict the outcome of an experiment in a laboratory setting or the observation of a phenomenon in nature. The prediction can also be statistical and deal only with probabilities. It is essential that the outcome of testing such a prediction be currently unknown. Only in this case does a successful outcome increase the probability that the hypothesis is true. If the outcome is already known, it is called a consequence and should have already been considered while formulating the hypothesis. If the predictions are not accessible by observation or experience, the hypothesis is not yet testable and so will remain to that extent unscientific in a strict sense. A new technology or theory might make the necessary experiments feasible. Thus, much scientifically based speculation might convince one (or many) that the hypothesis that other intelligent species exist is true. But since there no experiment now known which can test this hypothesis, science itself can have little to say about the possibility. In future, some new technique might lead to an experimental test and the speculation would then become part of accepted science. James D. Watson, Francis Crick, and others hypothesized that DNA had a helical structure. This implied that DNA's X-ray diffraction pattern would be 'x shaped'. This prediction followed from the work of Cochran, Crick and Vand (and independently by Stokes). The Cochran-Crick-Vand-Stokes theorem provided a mathematical explanation for the empirical observation that diffraction from helical structures produces x shaped patterns. In their first paper, Watson and Crick also noted that the double helix structure they proposed provided a simple mechanism for DNA replication, writing, "It has not escaped our notice that the specific pairing we have postulated immediately suggests a possible copying mechanism for the genetic material". ..4. DNA-experiments Another example: general relativity Einstein's theory of General Relativity makes several specific predictions about the observable structure of space-time, such as that light bends in a gravitational field, and that the amount of bending depends in a precise way on the strength of that gravitational field. Arthur Eddington's observations made during a 1919 solar eclipse supported General Relativity rather than Newtonian gravitation. Once predictions are made, they can be sought by experiments. If the test results contradict the predictions, the hypotheses which entailed them are called into question and become less tenable. Sometimes the experiments are conducted incorrectly or are not very well designed, when compared to a crucial experiment. If the experimental results confirm the predictions, then the hypotheses are considered more likely to be correct, but might still be wrong and continue to be subject to further testing. The experimental control is a technique for dealing with observational error. This technique uses the contrast between multiple samples (or observations) under differing conditions to see what varies or what remains the same. We vary the conditions for each measurement, to help isolate what has changed. Mill's canons can then help us figure out what the important factor is. Factor analysis is one technique for discovering the important factor in an effect. Depending on the predictions, the experiments can have different shapes. It could be a classical experiment in a laboratory setting, a double-blind study or an archaeological excavation. Even taking a plane from New York to Paris is an experiment which tests the aerodynamical hypotheses used for constructing the plane. Scientists assume an attitude of openness and accountability on the part of those conducting an experiment. Detailed record keeping is essential, to aid in recording and reporting on the experimental results, and supports the effectiveness and integrity of the procedure. They will also assist in reproducing the experimental results, likely by others. Traces of this approach can be seen in the work of Hipparchus (190–120 BCE), when determining a value for the precession of the Earth, while controlled experiments can be seen in the works of Jābir ibn Hayyān (721–815 CE), al-Battani (853–929) and Alhazen (965–1039). Watson and Crick showed an initial (and incorrect) proposal for the structure of DNA to a team from Kings College – Rosalind Franklin, Maurice Wilkins, and Raymond Gosling. Franklin immediately spotted the flaws which concerned the water content. Later Watson saw Franklin's detailed X-ray diffraction images which showed an X-shape and was able to confirm the structure was helical. This rekindled Watson and Crick's model building and led to the correct structure. ..1. DNA-characterizations Evaluation and improvement The scientific method is iterative. At any stage it is possible to refine its accuracy and precision, so that some consideration will lead the scientist to repeat an earlier part of the process. Failure to develop an interesting hypothesis may lead a scientist to re-define the subject under consideration. Failure of a hypothesis to produce interesting and testable predictions may lead to reconsideration of the hypothesis or of the definition of the subject. Failure of an experiment to produce interesting results may lead a scientist to reconsider the experimental method, the hypothesis, or the definition of the subject. Other scientists may start their own research and enter the process at any stage. They might adopt the characterization and formulate their own hypothesis, or they might adopt the hypothesis and deduce their own predictions. Often the experiment is not done by the person who made the prediction, and the characterization is based on experiments done by someone else. Published results of experiments can also serve as a hypothesis predicting their own reproducibility. After considerable fruitless experimentation, being discouraged by their superior from continuing, and numerous false starts, Watson and Crick were able to infer the essential structure of DNA by concrete modeling of the physical shapes of the nucleotides which comprise it. They were guided by the bond lengths which had been deduced by Linus Pauling and by Rosalind Franklin's X-ray diffraction images. ..DNA Example Science is a social enterprise, and scientific work tends to be accepted by the scientific community when it has been confirmed. Crucially, experimental and theoretical results must be reproduced by others within the scientific community. Researchers have given their lives for this vision; Georg Wilhelm Richmann was killed by ball lightning (1753) when attempting to replicate the 1752 kite-flying experiment of Benjamin Franklin. To protect against bad science and fraudulent data, government research-granting agencies such as the National Science Foundation, and science journals, including Nature and Science, have a policy that researchers must archive their data and methods so that other researchers can test the data and methods and build on the research that has gone before. Scientific data archiving can be done at a number of national archives in the U.S. or in the World Data Center. Models of scientific inquiry The classical model of scientific inquiry derives from Aristotle, who distinguished the forms of approximate and exact reasoning, set out the threefold scheme of abductive, deductive, and inductive inference, and also treated the compound forms such as reasoning by analogy. In 1877, Charles Sanders Peirce (// like "purse"; 1839–1914) characterized inquiry in general not as the pursuit of truth per se but as the struggle to move from irritating, inhibitory doubts born of surprises, disagreements, and the like, and to reach a secure belief, belief being that on which one is prepared to act. He framed scientific inquiry as part of a broader spectrum and as spurred, like inquiry generally, by actual doubt, not mere verbal or hyperbolic doubt, which he held to be fruitless. He outlined four methods of settling opinion, ordered from least to most successful: - The method of tenacity (policy of sticking to initial belief) – which brings comforts and decisiveness but leads to trying to ignore contrary information and others' views as if truth were intrinsically private, not public. It goes against the social impulse and easily falters since one may well notice when another's opinion is as good as one's own initial opinion. Its successes can shine but tend to be transitory. - The method of authority – which overcomes disagreements but sometimes brutally. Its successes can be majestic and long-lived, but it cannot operate thoroughly enough to suppress doubts indefinitely, especially when people learn of other societies present and past. - The method of the a priori – which promotes conformity less brutally but fosters opinions as something like tastes, arising in conversation and comparisons of perspectives in terms of "what is agreeable to reason." Thereby it depends on fashion in paradigms and goes in circles over time. It is more intellectual and respectable but, like the first two methods, sustains accidental and capricious beliefs, destining some minds to doubt it. - The scientific method – the method wherein inquiry regards itself as fallible and purposely tests itself and criticizes, corrects, and improves itself. Peirce held that slow, stumbling ratiocination can be dangerously inferior to instinct and traditional sentiment in practical matters, and that the scientific method is best suited to theoretical research, which in turn should not be trammeled by the other methods and practical ends; reason's "first rule" is that, in order to learn, one must desire to learn and, as a corollary, must not block the way of inquiry. The scientific method excels the others by being deliberately designed to arrive – eventually – at the most secure beliefs, upon which the most successful practices can be based. Starting from the idea that people seek not truth per se but instead to subdue irritating, inhibitory doubt, Peirce showed how, through the struggle, some can come to submit to truth for the sake of belief's integrity, seek as truth the guidance of potential practice correctly to its given goal, and wed themselves to the scientific method. For Peirce, rational inquiry implies presuppositions about truth and the real; to reason is to presuppose (and at least to hope), as a principle of the reasoner's self-regulation, that the real is discoverable and independent of our vagaries of opinion. In that vein he defined truth as the correspondence of a sign (in particular, a proposition) to its object and, pragmatically, not as actual consensus of some definite, finite community (such that to inquire would be to poll the experts), but instead as that final opinion which all investigators would reach sooner or later but still inevitably, if they were to push investigation far enough, even when they start from different points. In tandem he defined the real as a true sign's object (be that object a possibility or quality, or an actuality or brute fact, or a necessity or norm or law), which is what it is independently of any finite community's opinion and, pragmatically, depends only on the final opinion destined in a sufficient investigation. That is a destination as far, or near, as the truth itself to you or me or the given finite community. Thus, his theory of inquiry boils down to "Do the science." Those conceptions of truth and the real involve the idea of a community both without definite limits (and thus potentially self-correcting as far as needed) and capable of definite increase of knowledge. As inference, "logic is rooted in the social principle" since it depends on a standpoint that is, in a sense, unlimited. Paying special attention to the generation of explanations, Peirce outlined the scientific method as a coordination of three kinds of inference in a purposeful cycle aimed at settling doubts, as follows (in §III–IV in "A Neglected Argument" except as otherwise noted): - Abduction (or retroduction). Guessing, inference to explanatory hypotheses for selection of those best worth trying. From abduction, Peirce distinguishes induction as inferring, on the basis of tests, the proportion of truth in the hypothesis. Every inquiry, whether into ideas, brute facts, or norms and laws, arises from surprising observations in one or more of those realms (and for example at any stage of an inquiry already underway). All explanatory content of theories comes from abduction, which guesses a new or outside idea so as to account in a simple, economical way for a surprising or complicative phenomenon. Oftenest, even a well-prepared mind guesses wrong. But the modicum of success of our guesses far exceeds that of sheer luck and seems born of attunement to nature by instincts developed or inherent, especially insofar as best guesses are optimally plausible and simple in the sense, said Peirce, of the "facile and natural", as by Galileo's natural light of reason and as distinct from "logical simplicity". Abduction is the most fertile but least secure mode of inference. Its general rationale is inductive: it succeeds often enough and, without it, there is no hope of sufficiently expediting inquiry (often multi-generational) toward new truths. Coordinative method leads from abducing a plausible hypothesis to judging it for its testability and for how its trial would economize inquiry itself. Peirce calls his pragmatism "the logic of abduction". His pragmatic maxim is: "Consider what effects that might conceivably have practical bearings you conceive the objects of your conception to have. Then, your conception of those effects is the whole of your conception of the object". His pragmatism is a method of reducing conceptual confusions fruitfully by equating the meaning of any conception with the conceivable practical implications of its object's conceived effects—a method of experimentational mental reflection hospitable to forming hypotheses and conducive to testing them. It favors efficiency. The hypothesis, being insecure, needs to have practical implications leading at least to mental tests and, in science, lending themselves to scientific tests. A simple but unlikely guess, if uncostly to test for falsity, may belong first in line for testing. A guess is intrinsically worth testing if it has instinctive plausibility or reasoned objective probability, while subjective likelihood, though reasoned, can be misleadingly seductive. Guesses can be chosen for trial strategically, for their caution (for which Peirce gave as example the game of Twenty Questions), breadth, and incomplexity. One can hope to discover only that which time would reveal through a learner's sufficient experience anyway, so the point is to expedite it; the economy of research is what demands the leap, so to speak, of abduction and governs its art. - Deduction. Two stages: - Explication. Unclearly premissed, but deductive, analysis of the hypothesis in order to render its parts as clear as possible. - Demonstration: Deductive Argumentation, Euclidean in procedure. Explicit deduction of hypothesis's consequences as predictions, for induction to test, about evidence to be found. Corollarial or, if needed, Theorematic. - Induction. The long-run validity of the rule of induction is deducible from the principle (presuppositional to reasoning in general) that the real is only the object of the final opinion to which adequate investigation would lead; anything to which no such process would ever lead would not be real. Induction involving ongoing tests or observations follows a method which, sufficiently persisted in, will diminish its error below any predesignate degree. Three stages: - Classification. Unclearly premissed, but inductive, classing of objects of experience under general ideas. - Probation: direct inductive argumentation. Crude (the enumeration of instances) or gradual (new estimate of proportion of truth in the hypothesis after each test). Gradual induction is qualitative or quantitative; if qualitative, then dependent on weightings of qualities or characters; if quantitative, then dependent on measurements, or on statistics, or on countings. - Sentential Induction. "...which, by inductive reasonings, appraises the different probations singly, then their combinations, then makes self-appraisal of these very appraisals themselves, and passes final judgment on the whole result". Communication and community Frequently the scientific method is employed not only by a single person, but also by several people cooperating directly or indirectly. Such cooperation can be regarded as an important element of a scientific community. Various standards of scientific methodology are used within such an environment. Peer review evaluation Scientific journals use a process of peer review, in which scientists' manuscripts are submitted by editors of scientific journals to (usually one to three) fellow (usually anonymous) scientists familiar with the field for evaluation. In certain journals, the journal itself selects the referees; while in others (especially journals that are extremely specialized), the manuscript author might recommend referees. The referees may or may not recommend publication, or they might recommend publication with suggested modifications, or sometimes, publication in another journal. This standard is practiced to various degrees by different journals, and can have the effect of keeping the literature free of obvious errors and to generally improve the quality of the material, especially in the journals who use the standard most rigorously. The peer review process can have limitations when considering research outside the conventional scientific paradigm: problems of "groupthink" can interfere with open and fair deliberation of some new research. Documentation and replication Sometimes experimenters may make systematic errors during their experiments, veer from standard methods and practices (Pathological science) for various reasons, or, in rare cases, deliberately report false results. Occasionally because of this then, other scientists might attempt to repeat the experiments in order to duplicate the results. Researchers sometimes practice scientific data archiving, such as in compliance with the policies of government funding agencies and scientific journals. In these cases, detailed records of their experimental procedures, raw data, statistical analyses and source code can be preserved in order to provide evidence of the methodology and practice of the procedure and assist in any potential future attempts to reproduce the result. These procedural records may also assist in the conception of new experiments to test the hypothesis, and may prove useful to engineers who might examine the potential practical applications of a discovery. When additional information is needed before a study can be reproduced, the author of the study might be asked to provide it. They might provide it, or if the author refuses to share data, appeals can be made to the journal editors who published the study or to the institution which funded the research. Since it is impossible for a scientist to record everything that took place in an experiment, facts selected for their apparent relevance are reported. This may lead, unavoidably, to problems later if some supposedly irrelevant feature is questioned. For example, Heinrich Hertz did not report the size of the room used to test Maxwell's equations, which later turned out to account for a small deviation in the results. The problem is that parts of the theory itself need to be assumed in order to select and report the experimental conditions. The observations are hence sometimes described as being 'theory-laden'. Dimensions of practice The primary constraints on contemporary science are: - Publication, i.e. Peer review - Resources (mostly funding) It has not always been like this: in the old days of the "gentleman scientist" funding (and to a lesser extent publication) were far weaker constraints. Both of these constraints indirectly require scientific method – work that violates the constraints will be difficult to publish and difficult to get funded. Journals require submitted papers to conform to "good scientific practice" and to a degree this can be enforced by peer review. Originality, importance and interest are more important – see for example the author guidelines for Nature. Philosophy and sociology of science Philosophy of science looks at the underpinning logic of the scientific method, at what separates science from non-science, and the ethic that is implicit in science. There are basic assumptions, derived from philosophy by at least one prominent scientist, that form the base of the scientific method – namely, that reality is objective and consistent, that humans have the capacity to perceive reality accurately, and that rational explanations exist for elements of the real world. These assumptions from methodological naturalism form a basis on which science may be grounded. Logical Positivist, empiricist, falsificationist, and other theories have criticized these assumptions and given alternative accounts of the logic of science, but each has also itself been criticized. Thomas Kuhn examined the history of science in his The Structure of Scientific Revolutions, and found that the actual method used by scientists differed dramatically from the then-espoused method. His observations of science practice are essentially sociological and do not speak to how science is or can be practiced in other times and other cultures. Norwood Russell Hanson, Imre Lakatos and Thomas Kuhn have done extensive work on the "theory laden" character of observation. Hanson (1958) first coined the term for the idea that all observation is dependent on the conceptual framework of the observer, using the concept of gestalt to show how preconceptions can affect both observation and description. He opens Chapter 1 with a discussion of the Golgi bodies and their initial rejection as an artefact of staining technique, and a discussion of Brahe and Kepler observing the dawn and seeing a "different" sun rise despite the same physiological phenomenon. Kuhn and Feyerabend acknowledge the pioneering significance of his work. Kuhn (1961) said the scientist generally has a theory in mind before designing and undertaking experiments so as to make empirical observations, and that the "route from theory to measurement can almost never be traveled backward". This implies that the way in which theory is tested is dictated by the nature of the theory itself, which led Kuhn (1961, p. 166) to argue that "once it has been adopted by a profession ... no theory is recognized to be testable by any quantitative tests that it has not already passed". Paul Feyerabend similarly examined the history of science, and was led to deny that science is genuinely a methodological process. In his book Against Method he argues that scientific progress is not the result of applying any particular method. In essence, he says that for any specific method or norm of science, one can find a historic episode where violating it has contributed to the progress of science. Thus, if believers in scientific method wish to express a single universally valid rule, Feyerabend jokingly suggests, it should be 'anything goes'. Criticisms such as his led to the strong programme, a radical approach to the sociology of science. The postmodernist critiques of science have themselves been the subject of intense controversy. This ongoing debate, known as the science wars, is the result of conflicting values and assumptions between the postmodernist and realist camps. Whereas postmodernists assert that scientific knowledge is simply another discourse (note that this term has special meaning in this context) and not representative of any form of fundamental truth, realists in the scientific community maintain that scientific knowledge does reveal real and fundamental truths about reality. Many books have been written by scientists which take on this problem and challenge the assertions of the postmodernists while defending science as a legitimate method of deriving truth. Role of chance in discovery Somewhere between 33% and 50% of all scientific discoveries are estimated to have been stumbled upon, rather than sought out. This may explain why scientists so often express that they were lucky. Louis Pasteur is credited with the famous saying that "Luck favours the prepared mind", but some psychologists have begun to study what it means to be 'prepared for luck' in the scientific context. Research is showing that scientists are taught various heuristics that tend to harness chance and the unexpected. This is what Nassim Nicholas Taleb calls "Anti-fragility"; while some systems of investigation are fragile in the face of human error, human bias, and randomness, the scientific method is more than resistant or tough – it actually benefits from such randomness in many ways (it is anti-fragile). Taleb believes that the more anti-fragile the system, the more it will flourish in the real world. Psychologist Kevin Dunbar says the process of discovery often starts with researchers finding bugs in their experiments. These unexpected results lead researchers to try to fix what they think is an error in their method. Eventually, the researcher decides the error is too persistent and systematic to be a coincidence. The highly controlled, cautious and curious aspects of the scientific method are thus what make it well suited for identifying such persistent systematic errors. At this point, the researcher will begin to think of theoretical explanations for the error, often seeking the help of colleagues across different domains of expertise. ||The following text needs to be harmonized with text in History of scientific method. (June 2015)| ||This section may contain an excessive amount of intricate detail that may only interest a specific audience. (June 2015)| The development of the scientific method emerges in the history of science itself. Ancient Egyptian documents describe empirical methods in astronomy, mathematics, and medicine. The Greeks made contributions to the scientific method, most notably through Aristotle in his six works of logic collected as the Organon. Aristotle's inductive-deductive method used inductions from observations to infer general principles, deductions from those principles to check against further observations, and more cycles of induction and deduction to continue the advance of knowledge According to Karl Popper, Parmenides (fl. 5th century BCE) had conceived an axiomatic-deductive method. According to David Lindberg, Aristotle (4th century BCE) wrote about the scientific method even if he and his followers did not actually follow what he said. Lindberg also notes that Ptolemy (2nd century CE) and Ibn al-Haytham (11th century CE) are among the early examples of people who carried out scientific experiments. Also, John Losee writes that "the Physics and the Metaphysics contain discussions of certain aspects of scientific method", of which, he says "Aristotle viewed scientific inquiry as a progression from observations to general principles and back to observations." Early Christian leaders such as Clement of Alexandria (150–215) and Basil of Caesarea (330–379) encouraged future generations to view the Greek wisdom as "handmaidens to theology" and science was considered a means to more accurate understanding of the Bible and of God.:pp.4–5 Augustine of Hippo (354–430) who contributed great philosophical wealth to the Latin Middle Ages, advocated the study of science and was wary of philosophies that disagreed with the Bible, such as astrology and the Greek belief that the world had no beginning.:p.5 This Christian accommodation with Greek science "laid a foundation for the later widespread, intensive study of natural philosophy during the Late Middle Ages.":pp.8,9 However, the division of Latin-speaking Western Europe from the Greek-speaking East,:p.18 followed by barbarian invasions, the Plague of Justinian, and the Islamic conquests, resulted in the West largely losing access to Greek wisdom. By the 8th century Islam had conquered the Christian lands of Syria, Iraq, Iran and Egypt. This swift conquest further severed Western Europe from many of the great works of Aristotle, Plato, Euclid and others, many of which were housed in the great library of Alexandria. Having come upon such a wealth of knowledge, the Arabs, who viewed non-Arab languages as inferior, even as a source of pollution, employed conquered Christians and Jews to translate these works from the native Greek and Syriac into Arabic. Thus equipped, Arab philosopher Alhazen (Ibn al-Haytham) performed optical and physiological experiments, reported in his manifold works, the most famous being Book of Optics (1021). He was thus a forerunner of scientific method, having understood that a controlled environment involving experimentation and measurement is required in order to draw educated conclusions. Other Arab polymaths of the same era produced copious works on mathematics, philosophy, astronomy and alchemy. Most stuck closely to Aristotle, being hesitant to admit that some of Aristotle's thinking was errant, while others strongly criticized him. During these years, occasionally a paraphrased translation from the Arabic, which itself had been translated from Greek and Syriac, might make its way to the West for scholarly study. It was not until 1204, during which the Latins conquered and took Constantinople from the Byzantines in the name of the fourth Crusade, that a renewed scholarly interest in the original Greek manuscripts began to grow. Due to the new easier access to the libraries of Constantinople by Western scholars, a certain revival in the study and analysis of the original Greek texts by Western scholars began. From that point a functional scientific method that would launch modern science was on the horizon. Grosseteste (1175–1253), an English statesman, scientist and Christian theologian, was "the principal figure" in bringing about "a more adequate method of scientific inquiry" by which "medieval scientists were able eventually to outstrip their ancient European and Muslim teachers" (Dales 1973, p. 62). ... His thinking influenced Roger Bacon, who spread Grosseteste's ideas from Oxford to the University of Paris during a visit there in the 1240s. From the prestigious universities in Oxford and Paris, the new experimental science spread rapidly throughout the medieval universities: "And so it went to Galileo, William Gilbert, Francis Bacon, William Harvey, Descartes, Robert Hooke, Newton, Leibniz, and the world of the seventeenth century" (Crombie 1953, p. 15). "So it went to us as well " (Gauch 2003, pp. 52–53). Roger Bacon (c. 1214 – c. 1292), an English thinker and experimenter, is recognized by many to be the father of modern scientific method. His view that mathematics was essential to a correct understanding of natural philosophy was considered to be 400 years ahead of its time.:2 He was viewed as "a lone genius proclaiming the truth about time," having correctly calculated the calendar:3 His work in optics provided the platform on which Newton, Descartes, Huygens and others later transformed the science of light. Bacon's groundbreaking advances were due largely to his discovery that experimental science must be based on mathematics. (186–187) His works Opus Majus and De Speculis Comburentibus contain many "carefully drawn diagrams showing Bacon's meticulous investigations into the behavior of light.":66 He gives detailed descriptions of systematic studies using prisms and measurements by which he shows how a rainbow functions.:200 Others who advanced scientific method during this era included Albertus Magnus (c. 1193 – 1280), Theodoric of Freiberg, (c. 1250 – c. 1310), William of Ockham (c. 1285 – c. 1350), and Jean Buridan (c. 1300 – c. 1358). These were not only scientists but leaders of the church – Christian archbishops, friars and priests. By the late 15th century, the physician-scholar Niccolò Leoniceno was finding errors in Pliny's Natural History. As a physician, Leoniceno was concerned about these botanical errors propagating to the materia medica on which medicines were based. To counter this, a botanical garden was established at Orto botanico di Padova, University of Padua (in use for teaching by 1546), in order that medical students might have empirical access to the plants of a pharmacopia. The philosopher and physician Francisco Sanches was led by his medical training at Rome, 1571–73, and by the philosophical skepticism recently placed in the European mainstream by the publication of Sextus Empiricus' "Outlines of Pyrrhonism", to search for a true method of knowing (modus sciendi), as nothing clear can be known by the methods of Aristotle and his followers – for example, syllogism fails upon circular reasoning. Following the physician Galen's method of medicine, Sanches lists the methods of judgement and experience, which are faulty in the wrong hands, and we are left with the bleak statement That Nothing is Known (1581). This challenge was taken up by René Descartes in the next generation (1637), but at the least, Sanches warns us that we ought to refrain from the methods, summaries, and commentaries on Aristotle, if we seek scientific knowledge. In this, he is echoed by Francis Bacon, also influenced by skepticism; Sanches cites the humanist Juan Luis Vives who sought a better educational system, as well as a statement of human rights as a pathway for improvement of the lot of the poor. The modern scientific method crystallized no later than in the 17th and 18th centuries. In his work Novum Organum (1620) – a reference to Aristotle's Organon – Francis Bacon outlined a new system of logic to improve upon the old philosophical process of syllogism. Then, in 1637, René Descartes established the framework for scientific method's guiding principles in his treatise, Discourse on Method. The writings of Alhazen, Bacon and Descartes are considered critical in the historical development of the modern scientific method, as are those of John Stuart Mill. In the late 19th century, Charles Sanders Peirce proposed a schema that would turn out to have considerable influence in the development of current scientific methodology generally. Peirce accelerated the progress on several fronts. Firstly, speaking in broader context in "How to Make Our Ideas Clear" (1878), Peirce outlined an objectively verifiable method to test the truth of putative knowledge on a way that goes beyond mere foundational alternatives, focusing upon both deduction and induction. He thus placed induction and deduction in a complementary rather than competitive context (the latter of which had been the primary trend at least since David Hume, who wrote in the mid-to-late 18th century). Secondly, and of more direct importance to modern method, Peirce put forth the basic schema for hypothesis/testing that continues to prevail today. Extracting the theory of inquiry from its raw materials in classical logic, he refined it in parallel with the early development of symbolic logic to address the then-current problems in scientific reasoning. Peirce examined and articulated the three fundamental modes of reasoning that, as discussed above in this article, play a role in inquiry today, the processes that are currently known as abductive, deductive, and inductive inference. Thirdly, he played a major role in the progress of symbolic logic itself – indeed this was his primary specialty. Beginning in the 1930s, Karl Popper argued that there is no such thing as inductive reasoning. All inferences ever made, including in science, are purely deductive according to this view. Accordingly, he claimed that the empirical character of science has nothing to do with induction – but with the deductive property of falsifiability that scientific hypotheses have. Contrasting his views with inductivism and positivism, he even denied the existence of the scientific method: "(1) There is no method of discovering a scientific theory (2) There is no method for ascertaining the truth of a scientific hypothesis, i.e., no method of verification; (3) There is no method for ascertaining whether a hypothesis is 'probable', or probably true". Instead, he held that there is only one universal method, a method not particular to science: The negative method of criticism, or colloquially termed trial and error. It covers not only all products of the human mind, including science, mathematics, philosophy, art and so on, but also the evolution of life. Following Peirce and others, Popper argued that science is fallible and has no authority. In contrast to empiricist-inductivist views, he welcomed metaphysics and philosophical discussion and even gave qualified support to myths and pseudosciences. Popper's view has become known as critical rationalism. Although science in a broad sense existed before the modern era, and in many historical civilizations (as described above), modern science is so distinct in its approach and successful in its results that it now defines what science is in the strictest sense of the term. Relationship with mathematics Science is the process of gathering, comparing, and evaluating proposed models against observables. A model can be a simulation, mathematical or chemical formula, or set of proposed steps. Science is like mathematics in that researchers in both disciplines can clearly distinguish what is known from what is unknown at each stage of discovery. Models, in both science and mathematics, need to be internally consistent and also ought to be falsifiable (capable of disproof). In mathematics, a statement need not yet be proven; at such a stage, that statement would be called a conjecture. But when a statement has attained mathematical proof, that statement gains a kind of immortality which is highly prized by mathematicians, and for which some mathematicians devote their lives. Mathematical work and scientific work can inspire each other. For example, the technical concept of time arose in science, and timelessness was a hallmark of a mathematical topic. But today, the Poincaré conjecture has been proven using time as a mathematical concept in which objects can flow (see Ricci flow). Nevertheless, the connection between mathematics and reality (and so science to the extent it describes reality) remains obscure. Eugene Wigner's paper, The Unreasonable Effectiveness of Mathematics in the Natural Sciences, is a very well known account of the issue from a Nobel Prize-winning physicist. In fact, some observers (including some well known mathematicians such as Gregory Chaitin, and others such as Lakoff and Núñez) have suggested that mathematics is the result of practitioner bias and human limitation (including cultural ones), somewhat like the post-modernist view of science. George Pólya's work on problem solving, the construction of mathematical proofs, and heuristic show that the mathematical method and the scientific method differ in detail, while nevertheless resembling each other in using iterative or recursive steps. |Mathematical method||Scientific method| |1||Understanding||Characterization from experience and observation| |2||Analysis||Hypothesis: a proposed explanation| |3||Synthesis||Deduction: prediction from the hypothesis| |4||Review/Extend||Test and experiment| In Pólya's view, understanding involves restating unfamiliar definitions in your own words, resorting to geometrical figures, and questioning what we know and do not know already; analysis, which Pólya takes from Pappus, involves free and heuristic construction of plausible arguments, working backward from the goal, and devising a plan for constructing the proof; synthesis is the strict Euclidean exposition of step-by-step details of the proof; review involves reconsidering and re-examining the result and the path taken to it. Imre Lakatos argued that mathematicians actually use contradiction, criticism and revision as principles for improving their work. In like manner to science, where truth is sought, but certainty is not found, in Proofs and refutations (1976), what Lakatos tried to establish was that no theorem of informal mathematics is final or perfect. This means that we should not think that a theorem is ultimately true, only that no counterexample has yet been found. Once a counterexample, i.e. an entity contradicting/not explained by the theorem is found, we adjust the theorem, possibly extending the domain of its validity. This is a continuous way our knowledge accumulates, through the logic and process of proofs and refutations. (If axioms are given for a branch of mathematics, however, Lakatos claimed that proofs from those axioms were tautological, i.e. logically true, by rewriting them, as did Poincaré (Proofs and Refutations, 1976).) Lakatos proposed an account of mathematical knowledge based on Polya's idea of heuristics. In Proofs and Refutations, Lakatos gave several basic rules for finding proofs and counterexamples to conjectures. He thought that mathematical 'thought experiments' are a valid way to discover mathematical conjectures and proofs. Relationship with statistics The scientific method has been extremely successful in bringing the world out of medieval times, especially once it was combined with industrial processes. However, when the scientific method employs statistics as part of its arsenal, there are a number of both mathematical and practical issues that can have a deleterious effect on the reliability of the output of the scientific methods. This is outlined in detail in the most downloaded 2005 scientific paper "Why Most Published Research Findings Are False" ever by John Ioannidis. The particular points raised are statistical ("The smaller the studies conducted in a scientific field, the less likely the research findings are to be true" and "The greater the flexibility in designs, definitions, outcomes, and analytical modes in a scientific field, the less likely the research findings are to be true.") and economical ("The greater the financial and other interests and prejudices in a scientific field, the less likely the research findings are to be true" and "The hotter a scientific field (with more scientific teams involved), the less likely the research findings are to be true.") Hence: "Most research findings are false for most research designs and for most fields" and "As shown, the majority of modern biomedical research is operating in areas with very low pre- and poststudy probability for true findings." However: "Nevertheless, most new discoveries will continue to stem from hypothesis-generating research with low or very low pre-study odds," which means that *new* discoveries will come from research that, when that research started, had low or very low odds (a low or very low chance) of succeeding. Hence, if the scientific method is used to expand the frontiers of knowledge, research into areas that are outside the mainstream will yield most new discoveries. - Armchair theorizing - Empirical limits in science - Fuzzy logic - Information theory - Quantitative research - Replication crisis - Social research - Statistical hypothesis testing - Strong inference Problems and issues History, philosophy, sociology - Goldhaber & Nieto 2010, p. 940 - " Rules for the study of natural philosophy", Newton transl 1999, pp. 794–6, after Book 3, The System of the World. - From the Oxford English Dictionary definition for "scientific". - Peirce (1908), "A Neglected Argument for the Reality of God", Hibbert Journal v. 7, pp. 90–112. s:A Neglected Argument for the Reality of God with added notes. Reprinted with previously unpublished part, Collected Papers v. 6, paragraphs 452–85, The Essential Peirce v. 2, pp. 434–50, and elsewhere. - See, for example, Galileo 1638. His thought experiments disprove Aristotle's physics of falling bodies, in Two New Sciences. - Popper 1959:p273 - Karl R. Popper, Conjectures and Refutations: The Growth of Scientific Knowledge, Routledge, 2003 ISBN 0-415-28594-1 - Gauch, Hugh G. (2003). Scientific Method in Practice (Reprint ed.). Cambridge University Press. p. 3. ISBN 9780521017084. Retrieved 2015-01-26. The scientific method 'is often misrepresented as a fixed sequence of steps,' rather than being seen for what it truly is, 'a highly variable and creative process' (AAAS 2000:18). The claim here is that science has general principles that must be mastered to increase productivity and enhance perspective, not that these principles provide a simple and automated sequence of steps to follow. - History of Inductive Science (1837), and in Philosophy of Inductive Science (1840) - Jim Al-Khalili (4 January 2009). "The 'first true scientist'". BBC News. - Tracey Tokuhama-Espinosa (2010). Mind, Brain, and Education Science: A Comprehensive Guide to the New Brain-Based Teaching. W. W. Norton & Company. p. 39. ISBN 9780393706079. Alhazen (or Al-Haytham; 965–1039 C.E.) was perhaps one of the greatest physicists of all times and a product of the Islamic Golden Age or Islamic Renaissance (7th–13th centuries). He made significant contributions to anatomy, astronomy, engineering, mathematics, medicine, ophthalmology, philosophy, physics, psychology, and visual perception and is primarily attributed as the inventor of the scientific method, for which author Bradley Steffens (2006) describes him as the "first scientist". - Peirce, C. S., Collected Papers v. 1, paragraph 74. - Morris Kline (1985) Mathematics for the nonmathematician. Courier Dover Publications. p. 284. ISBN 0-486-24823-2 - Shapere, Dudley (1974). Galileo: A Philosophical Study. University of Chicago Press. ISBN 0-226-75007-8. - " The thesis of this book, as set forth in Chapter One, is that there are general principles applicable to all the sciences." __ Gauch 2003, p. xv - Peirce (1877), "The Fixation of Belief", Popular Science Monthly, v. 12, pp. 1–15. Reprinted often, including (Collected Papers of Charles Sanders Peirce v. 5, paragraphs 358–87), (The Essential Peirce, v. 1, pp. 109–23). Peirce.org Eprint. Wikisource Eprint. - Gauch 2003, p. 1 The scientific method can function in the same way; This is the principle of noncontradiction. - Francis Bacon(1629) New Organon, lists 4 types of error: Idols of the tribe (error due to the entire human race), the cave (errors due to an individual's own intellect), the marketplace (errors due to false words), and the theater (errors due to incredulous acceptance). - Peirce, C. S., Collected Papers v. 5, in paragraph 582, from 1898: ... [rational] inquiry of every type, fully carried out, has the vital power of self-correction and of growth. This is a property so deeply saturating its inmost nature that it may truly be said that there is but one thing needful for learning the truth, and that is a hearty and active desire to learn what is true. - Taleb contributes a brief description of anti-fragility, http://www.edge.org/q2011/q11_3.html - For example, the concept of falsification (first proposed in 1934) formalizes the attempt to disprove hypotheses rather than prove them. Karl R. Popper (1963), 'The Logic of Scientific Discovery'. The Logic of Scientific Discovery pp. 17–20, 249–252, 437–438, and elsewhere. - Leon Lederman, for teaching physics first, illustrates how to avoid confirmation bias: Ian Shelton, in Chile, was initially skeptical that supernova 1987a was real, but possibly an artifact of instrumentation (null hypothesis), so he went outside and disproved his null hypothesis by observing SN 1987a with the naked eye. The Kamiokande experiment, in Japan, independently observed neutrinos from SN 1987a at the same time. - Lindberg 2007, pp. 2–3: "There is a danger that must be avoided. ... If we wish to do justice to the historical enterprise, we must take the past for what it was. And that means we must resist the temptation to scour the past for examples or precursors of modern science. ...My concern will be with the beginnings of scientific theories, the methods by which they were formulated, and the uses to which they were put; ... " - "How does light travel through transparent bodies? Light travels through transparent bodies in straight lines only.... We have explained this exhaustively in our Book of Optics. But let us now mention something to prove this convincingly: the fact that light travels in straight lines is clearly observed in the lights which enter into dark rooms through holes.... [T]he entering light will be clearly observable in the dust which fills the air. – Alhazen, translated into English from German by M. Schwarz, from "Abhandlung über das Licht", J. Baarmann (ed. 1882) Zeitschrift der Deutschen Morgenländischen Gesellschaft Vol 36 as quoted in Sambursky 1974, p. 136. - He demonstrated his conjecture that "light travels through transparent bodies in straight lines only" by placing a straight stick or a taut thread next to the light beam, as quoted in Sambursky 1974, p. 136 to prove that light travels in a straight line. - David Hockney, (2001, 2006) in Secret Knowledge: rediscovering the lost techniques of the old masters ISBN 0-14-200512-6 (expanded edition) cites Alhazen several times as the likely source for the portraiture technique using the camera obscura, which Hockney rediscovered with the aid of an optical suggestion from Charles M. Falco. Kitab al-Manazir, which is Alhazen's Book of Optics, at that time denoted Opticae Thesaurus, Alhazen Arabis, was translated from Arabic into Latin for European use as early as 1270. Hockney cites Friedrich Risner's 1572 Basle edition of Opticae Thesaurus. Hockney quotes Alhazen as the first clear description of the camera obscura in Hockney, p. 240. - Galilei, Galileo (1638), Discorsi e Dimonstrazioni Matematiche, intorno a due nuoue scienze, Leida: Apresso gli Elsevirri, ISBN 0-486-60099-8, Dover reprint of the 1914 Macmillan translation by Henry Crew and Alfonso de Salvio of Two New Sciences, Galileo Galilei Linceo (1638). Additional publication information is from the collection of first editions of the Library of Congress surveyed by Bruno 1989, pp. 261–264. - Godfrey-Smith 2003 p. 236. - Gauch 2003, p. 3 - Schuster and Powers (2005), Translational and Experimental Clinical Research, Ch. 1. Link. This chapter also discusses the different types of research questions and how they are produced. - This phrasing is attributed to Marshall Nirenberg. - Note: for a discussion of multiple hypotheses, see Bayesian inference#Informal - October 1951, as noted in McElheny 2004, p. 40:"That's what a helix should look like!" Crick exclaimed in delight (This is the Cochran-Crick-Vand-Stokes theory of the transform of a helix). - June 1952, as noted in McElheny 2004, p. 43: Watson had succeeded in getting X-ray pictures of TMV showing a diffraction pattern consistent with the transform of a helix. - Watson did enough work on Tobacco mosaic virus to produce the diffraction pattern for a helix, per Crick's work on the transform of a helix. pp. 137–138, Horace Freeland Judson (1979) The Eighth Day of Creation ISBN 0-671-22540-5 - – Cochran W, Crick FHC and Vand V. (1952) "The Structure of Synthetic Polypeptides. I. The Transform of Atoms on a Helix", Acta Cryst., 5, 581–586. - Friday, January 30, 1953. Tea time, as noted in McElheny 2004, p. 52: Franklin confronts Watson and his paper – "Of course it [Pauling's pre-print] is wrong. DNA is not a helix." However, Watson then visits Wilkins' office, sees photo 51, and immediately recognizes the diffraction pattern of a helical structure. But additional questions remained, requiring additional iterations of their research. For example, the number of strands in the backbone of the helix (Crick suspected 2 strands, but cautioned Watson to examine that more critically), the location of the base pairs (inside the backbone or outside the backbone), etc. One key point was that they realized that the quickest way to reach a result was not to continue a mathematical analysis, but to build a physical model. - "The instant I saw the picture my mouth fell open and my pulse began to race." – Watson 1968, p. 167 Page 168 shows the X-shaped pattern of the B-form of DNA, clearly indicating crucial details of its helical structure to Watson and Crick. - McElheny 2004 p.52 dates the Franklin-Watson confrontation as Friday, January 30, 1953. Later that evening, Watson urges Wilkins to begin model-building immediately. But Wilkins agrees to do so only after Franklin's departure. - Saturday, February 28, 1953, as noted in McElheny 2004, pp. 57–59: Watson found the base pairing mechanism which explained Chargaff's rules using his cardboard models. - Galileo Galilei (1638) Two new sciences - Reconstruction of Galileo Galilei's experiment – the inclined plane - In Two new sciences, there are three 'reviewers': Simplicio, Sagredo, and Salviati, who serve as foil, antagonist, and protagonist. Galileo speaks for himself only briefly. But note that Einstein's 1905 papers were not peer reviewed before their publication. - Fleck 1979, pp. xxvii–xxviii - "NIH Data Sharing Policy." - Stanovich, Keith E. (2007). How to Think Straight About Psychology. Boston: Pearson Education. pg 123 - Brody 1993, pp. 44–45 - Hall, B. K.; Hallgrímsson, B., eds. (2008). Strickberger's Evolution (4th ed.). Jones & Bartlett. p. 762. ISBN 0-7637-0066-5. - Cracraft, J.; Donoghue, M. J., eds. (2005). Assembling the tree of life. Oxford University Press. p. 592. ISBN 0-19-517234-5. - Needham & Wang 1954 p.166 shows how the 'flying gallop' image propagated from China to the West. - "A myth is a belief given uncritical acceptance by members of a group ..." – Weiss, Business Ethics p. 15, as cited by Ronald R. Sims (2003) Ethics and corporate social responsibility: why giants fall p.21 - Imre Lakatos (1976), Proofs and Refutations. Taleb 2007, p. 72 lists ways to avoid narrative fallacy and confirmation bias. - For more on the narrative fallacy, see also Fleck 1979, p. 27: "Words and ideas are originally phonetic and mental equivalences of the experiences coinciding with them. ... Such proto-ideas are at first always too broad and insufficiently specialized. ... Once a structurally complete and closed system of opinions consisting of many details and relations has been formed, it offers enduring resistance to anything that contradicts it." - The scientific method requires testing and validation a posteriori before ideas are accepted. "Invariably one came up against fundamental physical limits to the accuracy of measurement. ... The art of physical measurement seemed to be a matter of compromise, of choosing between reciprocally related uncertainties. ... Multiplying together the conjugate pairs of uncertainty limits mentioned, however, I found that they formed invariant products of not one but two distinct kinds. ... The first group of limits were calculable a priori from a specification of the instrument. The second group could be calculated only a posteriori from a specification of what was done with the instrument. ... In the first case each unit [of information] would add one additional dimension (conceptual category), whereas in the second each unit would add one additional atomic fact.", – pp. 1–4: MacKay, Donald M. (1969), Information, Mechanism, and Meaning, Cambridge, MA: MIT Press, ISBN 0-262-63-032-X - See the hypothethico-deductive method, for example, Godfrey-Smith 2003, p. 236. - Jevons 1874, pp. 265–6. - pp. 65,73, 92, 398 – Andrew J. Galambos, Sic Itur ad Astra ISBN 0-88078-004-5(AJG learned scientific method from Felix Ehrenhaft - Galileo 1638, pp. v–xii,1–300 - Brody 1993, pp. 10–24 calls this the "epistemic cycle": "The epistemic cycle starts from an initial model; iterations of the cycle then improve the model until an adequate fit is achieved." - Iteration example: Chaldean astronomers such as Kidinnu compiled astronomical data. Hipparchus was to use this data to calculate the precession of the Earth's axis. Fifteen hundred years after Kidinnu, Al-Batani, born in what is now Turkey, would use the collected data and improve Hipparchus' value for the precession of the Earth's axis. Al-Batani's value, 54.5 arc-seconds per year, compares well to the current value of 49.8 arc-seconds per year (26,000 years for Earth's axis to round the circle of nutation). - Recursion example: the Earth is itself a magnet, with its own North and South Poles William Gilbert (in Latin 1600) De Magnete, or On Magnetism and Magnetic Bodies. Translated from Latin to English, selection by Moulton & Schifferes 1960, pp. 113–117. Gilbert created a terrella, a lodestone ground into a spherical shape, which served as Gilbert's model for the Earth itself, as noted in Bruno 1989, p. 277. - "The foundation of general physics ... is experience. These ... everyday experiences we do not discover without deliberately directing our attention to them. Collecting information about these is observation." – Hans Christian Ørsted("First Introduction to General Physics" ¶13, part of a series of public lectures at the University of Copenhagen. Copenhagen 1811, in Danish, printed by Johan Frederik Schulz. In Kirstine Meyer's 1920 edition of Ørsted's works, vol.III pp. 151–190. ) "First Introduction to Physics: the Spirit, Meaning, and Goal of Natural Science". Reprinted in German in 1822, Schweigger's Journal für Chemie und Physik 36, pp. 458–488, as translated in Ørsted 1997, p. 292 - "When it is not clear under which law of nature an effect or class of effect belongs, we try to fill this gap by means of a guess. Such guesses have been given the name conjectures or hypotheses." – Hans Christian Ørsted(1811) "First Introduction to General Physics" as translated in Ørsted 1997, p. 297. - "In general we look for a new law by the following process. First we guess it. ...", – Feynman 1965, p. 156 - "... the statement of a law – A depends on B – always transcends experience." – Born 1949, p. 6 - "The student of nature ... regards as his property the experiences which the mathematician can only borrow. This is why he deduces theorems directly from the nature of an effect while the mathematician only arrives at them circuitously." – Hans Christian Ørsted(1811) "First Introduction to General Physics" ¶17. as translated in Ørsted 1997, p. 297. - Salviati speaks: "I greatly doubt that Aristotle ever tested by experiment whether it be true that two stones, one weighing ten times as much as the other, if allowed to fall, at the same instant, from a height of, say, 100 cubits, would so differ in speed that when the heavier had reached the ground, the other would not have fallen more than 10 cubits." Two New Sciences (1638) – Galileo 1638, pp. 61–62. A more extended quotation is referenced by Moulton & Schifferes 1960, pp. 80–81. - In the inquiry-based education paradigm, the stage of "characterization, observation, definition, ..." is more briefly summed up under the rubric of a Question - "To raise new questions, new possibilities, to regard old problems from a new angle, requires creative imagination and marks real advance in science." – Einstein & Infeld 1938, p. 92. - Crawford S, Stucki L (1990), "Peer review and the changing research record", "J Am Soc Info Science", vol. 41, pp. 223–228 - See, e.g., Gauch 2003, esp. chapters 5–8 - Andreas Vesalius, Epistola, Rationem, Modumque Propinandi Radicis Chynae Decocti (1546), 141. Quoted and translated in C.D. O'Malley, Andreas Vesalius of Brussels, (1964), 116. As quoted by Bynum & Porter 2005, p. 597: Andreas Vesalius,597#1. - Crick, Francis (1994), The Astonishing Hypothesis ISBN 0-684-19431-7 p.20 - McElheny 2004 p.34 - Glen 1994, pp. 37–38. - "The structure that we propose is a three-chain structure, each chain being a helix" – Linus Pauling, as quoted on p. 157 by Horace Freeland Judson (1979), The Eighth Day of Creation ISBN 0-671-22540-5 - McElheny 2004, pp. 49–50: January 28, 1953 – Watson read Pauling's pre-print, and realized that in Pauling's model, DNA's phosphate groups had to be un-ionized. But DNA is an acid, which contradicts Pauling's model. - June 1952. as noted in McElheny 2004, p. 43: Watson had succeeded in getting X-ray pictures of TMV showing a diffraction pattern consistent with the transform of a helix. - McElheny 2004 p.68: Nature April 25, 1953. - In March 1917, the Royal Astronomical Society announced that on May 29, 1919, the occasion of a total eclipse of the sun would afford favorable conditions for testing Einstein's General theory of relativity. One expedition, to Sobral, Ceará, Brazil, and Eddington's expedition to the island of Principe yielded a set of photographs, which, when compared to photographs taken at Sobral and at Greenwich Observatory showed that the deviation of light was measured to be 1.69 arc-seconds, as compared to Einstein's desk prediction of 1.75 arc-seconds. – Antonina Vallentin (1954), Einstein, as quoted by Samuel Rapport and Helen Wright (1965), Physics, New York: Washington Square Press, pp 294–295. - Mill, John Stuart, "A System of Logic", University Press of the Pacific, Honolulu, 2002, ISBN 1-4102-0252-6. - al-Battani, De Motu Stellarum translation from Arabic to Latin in 1116, as cited by "Battani, al-" (c. 858 – 929) Encyclopædia Britannica, 15th. ed. Al-Battani is known for his accurate observations at al-Raqqah in Syria, beginning in 877. His work includes measurement of the annual precession of the equinoxes. - McElheny 2004 p.53: The weekend (January 31 – February 1) after seeing photo 51, Watson informed Bragg of the X-ray diffraction image of DNA in B form. Bragg gave them permission to restart their research on DNA (that is, model building). - McElheny 2004 p.54: On Sunday February 8, 1953, Maurice Wilkes gave Watson and Crick permission to work on models, as Wilkes would not be building models until Franklin left DNA research. - McElheny 2004 p.56: Jerry Donohue, on sabbatical from Pauling's lab and visiting Cambridge, advises Watson that textbook form of the base pairs was incorrect for DNA base pairs; rather, the keto form of the base pairs should be used instead. This form allowed the bases' hydrogen bonds to pair 'unlike' with 'unlike', rather than to pair 'like' with 'like', as Watson was inclined to model, on the basis of the textbook statements. On February 27, 1953, Watson was convinced enough to make cardboard models of the nucleotides in their keto form. - "Suddenly I became aware that an adenine-thymine pair held together by two hydrogen bonds was identical in shape to a guanine-cytosine pair held together by at least two hydrogen bonds. ..." – Watson 1968, pp. 194–197. - McElheny 2004 p.57 Saturday, February 28, 1953, Watson tried 'like with like' and admitted these base pairs didn't have hydrogen bonds that line up. But after trying 'unlike with unlike', and getting Jerry Donohue's approval, the base pairs turned out to be identical in shape (as Watson stated above in his 1968 Double Helix memoir quoted above). Watson now felt confident enough to inform Crick. (Of course, 'unlike with unlike' increases the number of possible codons, if this scheme were a genetic code.) - See, e.g., Physics Today, 59(1), p42. Richmann electrocuted in St. Petersburg (1753) - Aristotle, "Prior Analytics", Hugh Tredennick (trans.), pp. 181–531 in Aristotle, Volume 1, Loeb Classical Library, William Heinemann, London, UK, 1938. - "What one does not in the least doubt one should not pretend to doubt; but a man should train himself to doubt," said Peirce in a brief intellectual autobiography; see Ketner, Kenneth Laine (2009) "Charles Sanders Peirce: Interdisciplinary Scientist" in The Logic of Interdisciplinarity). Peirce held that actual, genuine doubt originates externally, usually in surprise, but also that it is to be sought and cultivated, "provided only that it be the weighty and noble metal itself, and no counterfeit nor paper substitute"; in "Issues of Pragmaticism", The Monist, v. XV, n. 4, pp. 481–99, see p. 484, and p. 491. (Reprinted in Collected Papers v. 5, paragraphs 438–63, see 443 and 451). - But see Scientific method and religion. - Peirce (1898), "Philosophy and the Conduct of Life", Lecture 1 of the Cambridge (MA) Conferences Lectures, published in Collected Papers v. 1, paragraphs 616–48 in part and in Reasoning and the Logic of Things, Ketner (ed., intro.) and Putnam (intro., comm.), pp. 105–22, reprinted in Essential Peirce v. 2, pp. 27–41. - " ... in order to learn, one must desire to learn ..." – Peirce (1899), "F.R.L." [First Rule of Logic], Collected Papers v. 1, paragraphs 135–40, Eprint at the Wayback Machine (archived January 6, 2012) - Peirce (1877), "How to Make Our Ideas Clear", Popular Science Monthly, v. 12, pp. 286–302. Reprinted often, including Collected Papers v. 5, paragraphs 388–410, Essential Peirce v. 1, pp. 124–41. ArisbeEprint. Wikisource Eprint. - Peirce (1868), "Some Consequences of Four Incapacities", Journal of Speculative Philosophy v. 2, n. 3, pp. 140–57. Reprinted Collected Papers v. 5, paragraphs 264–317, The Essential Peirce v. 1, pp. 28–55, and elsewhere. Arisbe Eprint - Peirce (1878), "The Doctrine of Chances", Popular Science Monthly v. 12, pp. 604–15, see pp. 610-11 via Internet Archive. Reprinted Collected Papers v. 2, paragraphs 645–68, Essential Peirce v. 1, pp. 142–54. "...death makes the number of our risks, the number of our inferences, finite, and so makes their mean result uncertain. The very idea of probability and of reasoning rests on the assumption that this number is indefinitely great. .... ...logicality inexorably requires that our interests shall not be limited. .... Logic is rooted in the social principle." - Peirce (c. 1906), "PAP (Prolegomena for an Apology to Pragmatism)" (Manuscript 293, not the like-named article), The New Elements of Mathematics (NEM) 4:319–320, see first quote under "Abduction" at Commens Dictionary of Peirce's Terms. - Peirce, Carnegie application (L75, 1902), New Elements of Mathematics v. 4, pp. 37–38: For it is not sufficient that a hypothesis should be a justifiable one. Any hypothesis which explains the facts is justified critically. But among justifiable hypotheses we have to select that one which is suitable for being tested by experiment. - Peirce (1902), Carnegie application, see MS L75.329–330, from Draft D of Memoir 27: Consequently, to discover is simply to expedite an event that would occur sooner or later, if we had not troubled ourselves to make the discovery. Consequently, the art of discovery is purely a question of economics. The economics of research is, so far as logic is concerned, the leading doctrine with reference to the art of discovery. Consequently, the conduct of abduction, which is chiefly a question of heuretic and is the first question of heuretic, is to be governed by economical considerations. - Peirce (1903), "Pragmatism – The Logic of Abduction", Collected Papers v. 5, paragraphs 195–205, especially 196. Eprint. - Peirce, "On the Logic of Drawing Ancient History from Documents", Essential Peirce v. 2, see pp. 107–9. On Twenty Questions, p. 109: Thus, twenty skillful hypotheses will ascertain what 200,000 stupid ones might fail to do. - Peirce (1878), "The Probability of Induction", Popular Science Monthly, v. 12, pp. 705–18, see 718 Google Books; 718 via Internet Archive. Reprinted often, including (Collected Papers v. 2, paragraphs 669–93), (The Essential Peirce v. 1, pp. 155–69). - Peirce (1905 draft "G" of "A Neglected Argument"), "Crude, Quantitative, and Qualitative Induction", Collected Papers v. 2, paragraphs 755–760, see 759. Find under "Induction" at Commens Dictionary of Peirce's Terms. - . Brown, C. (2005) Overcoming Barriers to Use of Promising Research Among Elite Middle East Policy Groups, Journal of Social Behaviour and Personality, Select Press. - Einstein, Albert (1936, 1956) One may say "the eternal mystery of the world is its comprehensibility." From the article "Physics and Reality" (1936), reprinted in Out of My Later Years (1956). 'It is one of the great realizations of Immanuel Kant that the setting up of a real external world would be senseless without this comprehensibility.' - Hanson, Norwood (1958), Patterns of Discovery, Cambridge University Press, ISBN 0-521-05197-5 - Kuhn 1962, p. 113 ISBN 978-1-4432-5544-8 - Feyerabend, Paul K (1960) "Patterns of Discovery" The Philosophical Review (1960) vol. 69 (2) pp. 247–252 - Kuhn, Thomas S., "The Function of Measurement in Modern Physical Science", ISIS 52(2), 161–193, 1961. - Feyerabend, Paul K., Against Method, Outline of an Anarchistic Theory of Knowledge, 1st published, 1975. Reprinted, Verso, London, UK, 1978. - Higher Superstition: The Academic Left and Its Quarrels with Science, The Johns Hopkins University Press, 1997 - Fashionable Nonsense: Postmodern Intellectuals' Abuse of Science, Picador; 1st Picador USA Pbk. Ed edition, 1999 - The Sokal Hoax: The Sham That Shook the Academy, University of Nebraska Press, 2000 ISBN 0-8032-7995-7 - A House Built on Sand: Exposing Postmodernist Myths About Science, Oxford University Press, 2000 - Intellectual Impostures, Economist Books, 2003 - Dunbar, K., & Fugelsang, J. (2005). Causal thinking in science: How scientists and students interpret the unexpected. In M. E. Gorman, R. D. Tweney, D. Gooding & A. Kincannon (Eds.), Scientific and Technical Thinking (pp. 57–79). Mahwah, NJ: Lawrence Erlbaum Associates. - Oliver, J.E. (1991) Ch2. of The incomplete guide to the art of discovery. New York:NY, Columbia University Press. - Riccardo Pozzo (2004) The impact of Aristotelianism on modern philosophy. CUA Press. p.41. ISBN 0-8132-1347-9 - The ancient Egyptians observed that heliacal rising of a certain star, Sothis (Greek for Sopdet (Egyptian), known to the West as Sirius), marked the annual flooding of the Nile river. See Neugebauer, Otto (1969) , The Exact Sciences in Antiquity (2 ed.), Dover Publications, ISBN 978-0-486-22332-2, p.82, and also the 1911 Britannica, "Egypt". - The Rhind papyrus lists practical examples in arithmetic and geometry – 1911 Britannica, "Egypt". - The Ebers papyrus lists some of the 'mysteries of the physician', as cited in the 1911 Britannica, "Egypt" - Gauch, Hugh G. (2003). Scientific Method in Practice. Cambridge University Press. p. 45. ISBN 9780521017084. Retrieved 10 February 2015. - Popper, Karl (1998). The world of Parmenides: essays on the Presocratic enlightenment. Routledge. p. 91. ISBN 0415173019. So what was really new in Parmenides was his axiomatic-deductive method, which Leucippus and Democritus turned into a hypothetical-deductive method, and thus made part of scientific methodology. - Lindberg, David (2007). The beginnings of western science: the European scientific tradition in philosophical, religious, and institutional context, Prehistory to A.D. 1450. The University of Chicago Press. p. 362. ISBN 0226482057. - Losee, John (2001). A Historical Introduction to the Philosophy of Science. Oxford University Press. pp. 4–5. ISBN 0198700555. - Grant, Edward. “The Foundations of Modern Science in the Middle Ages”. Cambridge University Press, UK. - Hodges, Richard and David Whitehouse. Mohammed, Charlemagne and the Origins of Europe. Cornell University Press, Ithaca, NY, 1983, p. 76. - Lewis, Bernard. Muslim Discovery of Europe. W. W. Norton and Company Ltd., New York, NY, 2001, p. 74. - Dowley, Tim, Ed. "The Baker Atlas of Christian History." 2001, p. 89 - Lewis, Muslim Discovery, p. 72. - Meri, Josef W. and Jere L. Bacharach. “Medieval Islamic Civilization”. Vol. 1 Index A – K. 2006, p. 304. - R. L. Verma (1969). Al-Hazen: father of modern optics. - Jaki, Stanley. Science and Creation: From Eternal Cycles to an Oscillating Universe. Science History Publications, NY, 1974, p. 206. - Mathpages. “Translating Aristotle”. Accessed July 2013. - George Sampson (1970). The concise Cambridge history of English literature. Cambridge University Press. p.174. ISBN 0-521-09581-6 - Clegg, Brian. “The First Scientist: A Life of Roger Bacon”. Carroll and Graf Publishers, NY, 2003 - Niccolò Leoniceno (1509), De Plinii et aliorum erroribus liber apud Ferrara, as cited by Sanches, Limbrick & Thomson 1988, p. 13 - 'I have sometimes seen a verbose quibbler attempting to persuade some ignorant person that white was black; to which the latter replied, "I do not understand your reasoning, since I have not studied as much as you have; yet I honestly believe that white differs from black. But pray go on refuting me for just as long as you like." ' – Sanches, Limbrick & Thomson 1988, p. 276 - Sanches, Limbrick & Thomson 1988, p. 278. - Bacon, Francis Novum Organum (The New Organon), 1620. Bacon's work described many of the accepted principles, underscoring the importance of empirical results, data gathering and experiment. Encyclopædia Britannica (1911), "Bacon, Francis" states: [In Novum Organum, we ] "proceed to apply what is perhaps the most valuable part of the Baconian method, the process of exclusion or rejection. This elimination of the non-essential, ..., is the most important of Bacon's contributions to the logic of induction, and that in which, as he repeatedly says, his method differs from all previous philosophies." - "John Stuart Mill (Stanford Encyclopedia of Philosophy)". plato.stanford.edu. Retrieved 2009-07-31. - Logik der Forschung, new appendices *XVII–*XIX (not yet available in the English edition Logic of scientific discovery) - Logic of Scientific discovery, p. 20 - Karl Popper: On the non-existence of scientific method. Realism and the Aim of Science (1983) - Karl Popper: Science: Conjectures and Refutations. Conjectures and Refuations, section VII - Karl Popper: On knowledge. In search of a better world, section II - "The historian ... requires a very broad definition of "science" – one that ... will help us to understand the modern scientific enterprise. We need to be broad and inclusive, rather than narrow and exclusive ... and we should expect that the farther back we go [in time] the broader we will need to be." – David Pingree (1992), "Hellenophilia versus the History of Science" Isis 83 554–63, as cited on p.3, David C. Lindberg (2007), The beginnings of Western science: the European Scientific tradition in philosophical, religious, and institutional context, Second ed. Chicago: Univ. of Chicago Press ISBN 978-0-226-48205-7 - "When we are working intensively, we feel keenly the progress of our work; we are elated when our progress is rapid, we are depressed when it is slow." – the mathematician Pólya 1957, p. 131 in the section on 'Modern heuristic'. - "Philosophy [i.e., physics] is written in this grand book – I mean the universe – which stands continually open to our gaze, but it cannot be understood unless one first learns to comprehend the language and interpret the characters in which it is written. It is written in the language of mathematics, and its characters are triangles, circles, and other geometrical figures, without which it is humanly impossible to understand a single word of it; without these, one is wandering around in a dark labyrinth." – Galileo Galilei, Il Saggiatore (The Assayer, 1623), as translated by Stillman Drake (1957), Discoveries and Opinions of Galileo pp. 237–8, as quoted by di Francia 1981, p. 10. - Pólya 1957 2nd ed. - George Pólya (1954), Mathematics and Plausible Reasoning Volume I: Induction and Analogy in Mathematics, - George Pólya (1954), Mathematics and Plausible Reasoning Volume II: Patterns of Plausible Reasoning. - Pólya 1957, p. 142 - Pólya 1957, p. 144 - Mackay 1991 p.100 - See the development, by generations of mathematicians, of Euler's formula for polyhedra as documented by Lakatos, Imre (1976), Proofs and refutations, Cambridge: Cambridge University Press, ISBN 0-521-29038-4 - Lakatos, Imre (Worrall & Zahar, eds. 1976) Proofs and Refutations, p.55 - Rosenberg, Nathan, Luther Earle Birdzell, and Glenn William Mitchell. How the West grew rich. Popular Prakashan, 1986. - Ioannidis, John P. A. (2005-08-01). "Why Most Published Research Findings Are False". PLoS Medicine 2 (8). doi:10.1371/journal.pmed.0020124. ISSN 1549-1277. PMC 1182327. PMID 16060722. - Born, Max (1949), Natural Philosophy of Cause and Chance, Peter Smith, also published by Dover, 1964. From the Waynflete Lectures, 1948. On the web. N.B.: the web version does not have the 3 addenda by Born, 1950, 1964, in which he notes that all knowledge is subjective. Born then proposes a solution in Appendix 3 (1964) - Brody, Thomas A. (1993), The Philosophy Behind Physics, Springer Verlag, ISBN 0-387-55914-0. (Luis De La Peña and Peter E. Hodgson, eds.) - Bruno, Leonard C. (1989), The Landmarks of Science, ISBN 0-8160-2137-6 - Bynum, W.F.; Porter, Roy (2005), Oxford Dictionary of Scientific Quotations, Oxford, ISBN 0-19-858409-1. - Dales, Richard C. (1973), The Scientific Achievement of the Middle Ages (The Middle Ages Series), University of Pennsylvania Press, ISBN 9780812210576 - di Francia, G. Toraldo (1981), The Investigation of the Physical World, Cambridge University Press, ISBN 0-521-29925-X. - Einstein, Albert; Infeld, Leopold (1938), The Evolution of Physics: from early concepts to relativity and quanta, New York: Simon and Schuster, ISBN 0-671-20156-5 - Feynman, Richard (1965), The Character of Physical Law, Cambridge: M.I.T. Press, ISBN 0-262-56003-8. - Fleck, Ludwik (1979), Genesis and Development of a Scientific Fact, Univ. of Chicago, ISBN 0-226-25325-2. (written in German, 1935, Entstehung und Entwickelung einer wissenschaftlichen Tatsache: Einführung in die Lehre vom Denkstil und Denkkollectiv) English translation, 1979 - Galileo (1638), Two New Sciences, Leiden: Lodewijk Elzevir, ISBN 0-486-60099-8 Translated from Italian to English in 1914 by Henry Crew and Alfonso de Salvio. Introduction by Antonio Favaro. xxv+300 pages, index. New York: Macmillan, with later reprintings by Dover. - Gauch, Hugh G., Jr. (2003), Scientific Method in Practice, Cambridge University Press, ISBN 0-521-01708-4 435 pages - Glen, William (ed.) (1994), The Mass-Extinction Debates: How Science Works in a Crisis, Stanford, CA: Stanford University Press, ISBN 0-8047-2285-4. - Godfrey-Smith, Peter (2003), Theory and Reality: An introduction to the philosophy of science, University of Chicago Press, ISBN 0-226-30063-3. - Goldhaber, Alfred Scharff; Nieto, Michael Martin (January–March 2010), "Photon and graviton mass limits", Rev. Mod. Phys. (American Physical Society) 82: 939, doi:10.1103/RevModPhys.82.939. pages 939–979. - Jevons, William Stanley (1874), The Principles of Science: A Treatise on Logic and Scientific Method, Dover Publications, ISBN 1-4304-8775-5. 1877, 1879. Reprinted with a foreword by Ernst Nagel, New York, NY, 1958. - Kuhn, Thomas S. (1962), The Structure of Scientific Revolutions, Chicago, IL: University of Chicago Press. 2nd edition 1970. 3rd edition 1996. - Lindberg, David C. (2007), The Beginnings of Western Science, University of Chicago Press 2nd edition 2007. - Mackay, Alan L. (ed.) (1991), Dictionary of Scientific Quotations, London: IOP Publishing Ltd, ISBN 0-7503-0106-6 - McElheny, Victor K. (2004), Watson & DNA: Making a scientific revolution, Basic Books, ISBN 0-7382-0866-3. - Moulton, Forest Ray; Schifferes, Justus J. (eds., Second Edition) (1960), The Autobiography of Science, Doubleday. - Needham, Joseph; Wang, Ling (王玲) (1954), Science and Civilisation in China, 1 Introductory Orientations, Cambridge University Press - Newton, Isaac (1999) [1687, 1713, 1726], Philosophiae Naturalis Principia Mathematica, University of California Press, ISBN 0-520-08817-4, Third edition. From I. Bernard Cohen and Anne Whitman's 1999 translation, 974 pages. - Ørsted, Hans Christian (1997), Selected Scientific Works of Hans Christian Ørsted, Princeton, ISBN 0-691-04334-5. Translated to English by Karen Jelved, Andrew D. Jackson, and Ole Knudsen, (translators 1997). - Peirce, C. S. – see Charles Sanders Peirce bibliography. - Poincaré, Henri (1905), Science and Hypothesis Eprint - Pólya, George (1957), How to Solve It, Princeton University Press, ISBN 978-4871878302, OCLC 706968824 (reprinted 2009) - Popper, Karl R. (1959), The Logic of Scientific Discovery 1934, 1959. - Sambursky, Shmuel (ed.) (1974), Physical Thought from the Presocratics to the Quantum Physicists, Pica Press, ISBN 0-87663-712-8. - Sanches, Francisco; Limbrick, Elaine. Introduction, Notes, and Bibliography; Thomson, Douglas F.S. Latin text established, annotated, and translated. (1988), That Nothing is Known, Cambridge: Cambridge University Press, ISBN 0-521-35077-8 Critical edition. - Taleb, Nassim Nicholas (2007), The Black Swan, Random House, ISBN 978-1-4000-6351-2 - Bauer, Henry H., Scientific Literacy and the Myth of the Scientific Method, University of Illinois Press, Champaign, IL, 1992 - Beveridge, William I. B., The Art of Scientific Investigation, Heinemann, Melbourne, Australia, 1950. - Bernstein, Richard J., Beyond Objectivism and Relativism: Science, Hermeneutics, and Praxis, University of Pennsylvania Press, Philadelphia, PA, 1983. - Brody, Baruch A. and Capaldi, Nicholas, Science: Men, Methods, Goals: A Reader: Methods of Physical Science, W. A. Benjamin, 1968 - Brody, Baruch A., and Grandy, Richard E., Readings in the Philosophy of Science, 2nd edition, Prentice Hall, Englewood Cliffs, NJ, 1989. - Burks, Arthur W., Chance, Cause, Reason – An Inquiry into the Nature of Scientific Evidence, University of Chicago Press, Chicago, IL, 1977. - Alan Chalmers. What is this thing called science?. Queensland University Press and Open University Press, 1976. - Crick, Francis (1988), What Mad Pursuit: A Personal View of Scientific Discovery, New York: Basic Books, ISBN 0-465-09137-7. - Crombie, A. C. (1953), Robert Grosseteste and the Origins of Experimental Science 1100–1700, Oxford - Dewey, John, How We Think, D.C. Heath, Lexington, MA, 1910. Reprinted, Prometheus Books, Buffalo, NY, 1991. - Earman, John (ed.), Inference, Explanation, and Other Frustrations: Essays in the Philosophy of Science, University of California Press, Berkeley & Los Angeles, CA, 1992. - Fraassen, Bas C. van, The Scientific Image, Oxford University Press, Oxford, UK, 1980. - Franklin, James (2009), What Science Knows: And How It Knows It, New York: Encounter Books, ISBN 1-59403-207-6. - Gadamer, Hans-Georg, Reason in the Age of Science, Frederick G. Lawrence (trans.), MIT Press, Cambridge, MA, 1981. - Giere, Ronald N. (ed.), Cognitive Models of Science, vol. 15 in 'Minnesota Studies in the Philosophy of Science', University of Minnesota Press, Minneapolis, MN, 1992. - Hacking, Ian, Representing and Intervening, Introductory Topics in the Philosophy of Natural Science, Cambridge University Press, Cambridge, UK, 1983. - Heisenberg, Werner, Physics and Beyond, Encounters and Conversations, A.J. Pomerans (trans.), Harper and Row, New York, NY 1971, pp. 63–64. - Holton, Gerald, Thematic Origins of Scientific Thought, Kepler to Einstein, 1st edition 1973, revised edition, Harvard University Press, Cambridge, MA, 1988. - Kuhn, Thomas S., The Essential Tension, Selected Studies in Scientific Tradition and Change, University of Chicago Press, Chicago, IL, 1977. - Latour, Bruno, Science in Action, How to Follow Scientists and Engineers through Society, Harvard University Press, Cambridge, MA, 1987. - Losee, John, A Historical Introduction to the Philosophy of Science, Oxford University Press, Oxford, UK, 1972. 2nd edition, 1980. - Maxwell, Nicholas, The Comprehensibility of the Universe: A New Conception of Science, Oxford University Press, Oxford, 1998. Paperback 2003. - McCarty, Maclyn (1985), The Transforming Principle: Discovering that genes are made of DNA, New York: W. W. Norton, pp. 252 , ISBN 0-393-30450-7. Memoir of a researcher in the Avery–MacLeod–McCarty experiment. - McComas, William F., ed. PDF (189 KB), from The Nature of Science in Science Education, pp53–70, Kluwer Academic Publishers, Netherlands 1998. - Misak, Cheryl J., Truth and the End of Inquiry, A Peircean Account of Truth, Oxford University Press, Oxford, UK, 1991. - Piattelli-Palmarini, Massimo (ed.), Language and Learning, The Debate between Jean Piaget and Noam Chomsky, Harvard University Press, Cambridge, MA, 1980. - Popper, Karl R., Unended Quest, An Intellectual Autobiography, Open Court, La Salle, IL, 1982. - Putnam, Hilary, Renewing Philosophy, Harvard University Press, Cambridge, MA, 1992. - Rorty, Richard, Philosophy and the Mirror of Nature, Princeton University Press, Princeton, NJ, 1979. - Salmon, Wesley C., Four Decades of Scientific Explanation, University of Minnesota Press, Minneapolis, MN, 1990. - Shimony, Abner, Search for a Naturalistic World View: Vol. 1, Scientific Method and Epistemology, Vol. 2, Natural Science and Metaphysics, Cambridge University Press, Cambridge, UK, 1993. - Thagard, Paul, Conceptual Revolutions, Princeton University Press, Princeton, NJ, 1992. - Ziman, John (2000). Real Science: what it is, and what it means. Cambridge, UK: Cambridge University Press. |Wikibooks has a book on the topic of: The Scientific Method| - Scientific Method entry by Anne Andersen and Brian Hepburn in the Stanford Encyclopedia of Philosophy - Confirmation and Induction entry in the Internet Encyclopedia of Philosophy - Scientific method at PhilPapers - Scientific method at the Indiana Philosophy Ontology Project - An Introduction to Science: Scientific Thinking and a scientific method by Steven D. Schafersman. - Introduction to the scientific method at the University of Rochester - Theory-ladenness by Paul Newall at The Galilean Library - Lecture on Scientific Method by Greg Anderson - Using the scientific method for designing science fair projects - SCIENTIFIC METHODS an online book by Richard D. Jarrard - Richard Feynman on the Key to Science (one minute, three seconds), from the Cornell Lectures. - Lectures on the Scientific Method by Nick Josh Karean, Kevin Padian, Michael Shermer and Richard Dawkins |Library resources about

https://en.wikipedia.org/wiki/Scientific_pluralism

On a ship, the fire room, or FR or boiler room or stokehold, referred to the space of a vessel where water was brought to a boil. The steam was then transmitted to a separate engine room, located immediately aft, where it was utilized to power the vessel. To increase the safety and damage survivability of a vessel, the machinery necessary for operations may be segregated into various spaces, the fire room was one of these spaces, and was among the largest physical compartment of the machinery space. On some ships, the space comprised more than one fire room, such as forward and aft, or port or starboard fire rooms, or may be simply numbered. Each room was connected to a stack ventilating smoke. By their nature, fire rooms were less complex than their allied engine room and were normally supervised by less senior personnel. On a large percentage of vessels, ships and boats, the fire room was located near the bottom, and at the rear, or aft, end of the vessel, and usually comprised few compartments. This design maximized the cargo carrying capacity of the vessel. The fire room on some ships was situated amid-ships, especially on vessels built from the 1880s to the 1990s. Vessels typically contained several engines for different purposes. Main, or propulsion engines are used to turn the ship's propeller and move the ship through the water. Their allied fire room typically burned heavy fuel oil, replacing the earlier use of coal. There was a mechanism for removing ash from combustion that did not rise out of the stack. On a steamship, power for both electricity and propulsion is provided by one or more large boilers giving rise to the alternate name boiler room. The latter name was preferred in the British Navy, among others. High pressure steam from the boiler is piped to the engine room to drive reciprocating engines or turbines for propulsion, and turbo generators for electricity. Naval ships typically were able to generate a large volume of smoke by changing the fuel mix. Prior to the heavy use of radar, a smoke screen could be used to mask the movement of ships. Damage control was enhanced by the separation of the fire and engine rooms. In the event of damage to its associated engine room, steam could be transmitted to another engine room. In turn, an engine room could still operate though its associated fire room had become inoperative. Two engineering advances resulted in the disappearance of the fire room in the early 1990s. The first was the movement by naval shipbuilding to nuclear-powered vessels. If a room containing nuclear material was subjected to damage, it was assumed that the event would likely result in abandonment of the ship regardless of the separation of rooms. The second was the adoption of gas turbines in place of oil-fired boilers for all other navy ships. These powered engines directly and needed no boilers. Fire rooms were hot, sometimes dirty, and potentially dangerous. The presence of flammable fuel meant that a fire hazard existed in the fire room, which was monitored continuously by the ship's engineering staff and various monitoring systems. Fire rooms employed some means of providing air for the operation of the flame to ignite the oil and associated ventilation. Only spot ventilation was practical to keep personnel cool. This would require an unrestricted hull opening of the same size as the intake area of the engine itself assuming the hull opening is in the fire room itself. Forced draft fire rooms were used until World War II. These required that personnel enter through an air lock to maintain the pressure. These were abandoned when the forced draft occasionally failed and blowback occurred killing fire room personnel. Commonly, screens were placed over openings reducing airflow by approximately 50% so the opening area was increased appropriately. The requirement for general ventilation and the requirement for sufficient combustion air are quite different. A typical arrangement might be to make the opening large enough to provide intake air plus 1,000 cubic feet (28 m3) per Minute (CFM) for additional ventilation. Engines pull sufficient air into the fire room for their own operation. However, additional airflow for ventilation usually requires intake and exhaust blowers. When fired up, there were personnel assignments specified underway, as well as in port. For example, for an Iowa-class battleship, in normal steaming four boilers were operated. This was sufficient to power the ships at speeds up to 27 knots (50 km/h). For higher speeds, all eight boilers were lit. Each operating boiler required a minimum of four trained operators on watch: a boiler supervisor (BTOW), a superheater burnerman and saturated burnerman to control the steam temperature and pressure and a checkman, who monitored and controlled the water level in the steam drum. In addition, there was a fireroom messenger and a lower level pumpman on duty whenever the fireroom was steaming. - Kearny incident - Engineering department - Marine propulsion - Marine fuel management - Mechanical room - Electrical room - International Marine Engineering 22. Simmons-Boardman Publishing Company. 1917. p. 298. ISSN 0272-2879. Retrieved 2015-02-22. - Personnel Qualification Standard for BB-61 Class Engineering (NAVEDTRA 43404-7A). Chief of Naval Education and Training. 1986.

https://en.wikipedia.org/wiki/Fire_room

|Target disease||Haemophilus influenzae type b| |(what is this?)| Haemophilus influenzae type B vaccine is a vaccine used to prevent Haemophilus influenzae type b (Hib) infection. In countries that include it as a routine vaccine, rates of severe Hib infections have decreased more than 90%. It has therefore resulted in a decrease in the rate of meningitis, pneumonia, and epiglottitis. It is recommended by both the World Health Organization and Centers for Disease Control and Prevention. Two or three doses should be given before six months of age. The first dose is recommended around six weeks of age with four weeks between doses. If only two doses are used, another dose later in life is recommended. It is given by injection into a muscle. Severe side effects are uncommon. About 20 to 25% of people develop pain at the site of injection while about 2% develop a fever. There is no clear association with severe allergic reactions. The Hib vaccine is available by itself, in combination with the diptheria/tetanus/pertussis vaccine, and in combination with the hepatitis B vaccine, among others. All Hib vaccines that are currently used are conjugate vaccine. An initial Hib vaccine was developed in 1977 which was replaced by a more effective formulation in the 1990s. As of 2013, 184 countries include it in their routine vaccinations. It is on the World Health Organization's List of Essential Medicines, the most important medications needed in a basic health system. The wholesale cost of a pentavalent vaccine which includes Hib was 15.40 USD per dose as of 2014. In the United States it costs about 25 to 50 USD per dose. Hib conjugate vaccines have been shown to be universally effective against all manifestations of Hib disease, with a clinical efficacy among fully vaccinated children estimated to be between 95–100%. The vaccine has also been shown to be immunogenic in patients at high risk of invasive disease. Hib vaccine is not effective against non-type B Haemophilus influenzae. However, non-type B disease is rare in comparison to pre-vaccine Haemophilus influenzae type B disease. Prior to introduction of the conjugate vaccine, Hib was a leading cause of childhood meningitis, pneumonia, and epiglottitis in the United States, causing an estimated 20,000 cases a year in the early 1980s, mostly in children under 5 years old. Since routine vaccination began, the incidence of Hib disease has declined by greater than 99%, effectively eliminating Hib as a public health problem. Similar reductions in disease occurred after introduction of the vaccine in Western Europe and developing countries. Although Hib vaccine is given to children, Hib infections have also decreased in adults. This decrease occurred because of herd immunity; children infected with Hib carry the bacteria in their nasal passages while clearing the infection. These Hib-carrying children would regularly infect adults. The practice of vaccinating children eliminated the source of the bacteria, reducing the rate of Hib in adults. The CDC and WHO currently recommend that all infants be vaccinated using a polysaccharide-protein conjugate Hib vaccine, starting after the age of 6 weeks. The vaccination is also indicated in asplenic patients. Clinical trials and ongoing surveillance have shown Hib vaccine to be safe. In general, adverse reactions to the vaccine are mild. The most common reactions are mild fever, loss of appetite, transient redness, swelling, or pain at the site of injection, occurring in 5–30% of vaccine recipients. More severe reactions are extremely rare. Introduction of Hib vaccine in developing countries lagged behind that in developed countries for several reasons. The expense of the vaccine was large in comparison to the standard EPI vaccines. Poor disease surveillance systems and inadequate hospital laboratories failed to detect the disease, leading many experts to believe that Hib did not exist in their countries. And health systems in many countries were struggling with the current vaccines they were trying to deliver. GAVI and the Hib Initiative In order to remedy these issues, the GAVI Alliance took active interest in the vaccine. GAVI offers substantial subsidization of Hib vaccine for countries interested in using the vaccine, as well as financial support for vaccine systems and safe injections. In addition, GAVI created the Hib Initiative to catalyze uptake of the vaccine. The Hib Initiative uses a combination of collecting and disseminating existing data, research, and advocacy to assist countries in the making a decision about using the Hib vaccine. Currently[update], 61 out of 72 low-income countries are planning on introducing the vaccine by the end of 2009. The first Hib vaccine licensed was a pure polysaccharide vaccine, first marketed in the US in 1985. Similar to other polysaccharide vaccines, immune response to the vaccine was highly age-dependent. Children under 18 months of age did not produce a positive response for this vaccine. As a result, the age group with the highest incidence of Hib disease was unprotected, limiting the usefulness of the vaccine. The vaccine was withdrawn from the market in 1988. The shortcomings of the polysaccharide vaccine led to the production of the Hib polysaccharide-protein conjugate vaccine. Attaching Hib polysaccharide to a protein carrier greatly increased the ability of the immune system of young children to recognize the polysaccharide and develop immunity. There are currently three types of conjugate vaccine, utilizing different carrier proteins for the conjugation process, all of which are highly effective and safe: inactivated tetanospasmin (also called tetanus toxoid), mutant diphtheria protein, and meningococcal group B outer membrane protein. Multiple combinations of Hib and other vaccines have been licensed in the United States, reducing the number of shots necessary to vaccinate a child. Hib vaccine combined with diphtheria-tetanus-pertussis–polio vaccines and Hepatitis B vaccines are available in the US. The World Health Organization (WHO) has certified several Hib vaccine combinations, including a pentavalent diphtheria-pertussis-tetanus-hepatitis B-Hib, for use in developing countries. There is not yet sufficient evidence on how effective this combined pentavalent vaccine is in relation to the individual vaccines. - "Haemophilus influenzae type b (Hib) Vaccination Position Paper – July 2013." (PDF). Wkly Epidemiol Rec 88 (39): 413–26. Sep 27, 2013. PMID 24143842. - "Haemophilus b conjugate vaccines for prevention of Haemophilus influenzae type b disease among infants and children two months of age and older. Recommendations of the immunization practices advisory committee (ACIP).". MMWR Recomm Rep 40 (RR-1): 1–7. Jan 11, 1991. PMID 1899280. - "WHO Model List of EssentialMedicines" (PDF). World Health Organization. October 2013. Retrieved 22 April 2014. - "Vaccine, Pentavalent". International Drug Price Indicator Guide. Retrieved 7 December 2015. - Hamilton, Richart (2015). Tarascon Pocket Pharmacopoeia 2015 Deluxe Lab-Coat Edition. Jones & Bartlett Learning. p. 313. ISBN 9781284057560. - "Recommendation of the Immunization Practices Advisory Committee (ACIP) Polysaccharide Vaccine for Prevention of Haemophilus influenzae Type b Disease". MMWR Weekly 34 (15): 201–5. 1985-04-19. ISSN 0149-2195. Retrieved 2008-10-03. - "Haemophilus influenzae Disease (Including Hib)". Disease Listing. Centers for Disease Control and Prevention. 2012-09-25. Retrieved 2014-01-31. - "Hib Initiative". Retrieved 2008-10-03. 61 of 72 GAVI countries have introduced or will introduce Hib vaccine into their routine immunization program [sic] by 2009 - Centers for Disease Control and Prevention (2006). Atkinson W, Hamborsky J, McIntyre L, Wolfe S, ed. Epidemiology and Prevention of Vaccine-Preventable Diseases (9th ed.). Washington, D.C.: Public Health Foundation. - Bar-On, ES; Goldberg, E; Hellmann, S; Leibovici, L (18 April 2012). "Combined DTP-HBV-HIB vaccine versus separately administered DTP-HBV and HIB vaccines for primary prevention of diphtheria, tetanus, pertussis, hepatitis B and Haemophilus influenzae B (HIB).". The Cochrane database of systematic reviews 4: CD005530. doi:10.1002/14651858.CD005530.pub3. PMID 22513932.

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A new solar laser could be instrumental in the quest to use magnesium as a source of energy. A new kind of efficient, solar-powered laser has been developed by researchers at the Tokyo Institute of Technology, in Japan. They hope to use the laser to help them realize their goal of developing a magnesium combustion engine. The researchers described the new laser in a recent issue of Applied Physics Letters. The idea, says Takashi Yabe, a professor of mechanical engineering and science at the Tokyo Institute, is to make a powerful laser capable of combusting the magnesium content of seawater. In the process, large amounts of heat and hydrogen are given off. Magnesium has great potential as an energy source because it has an energy storage density about 10 times higher than that of hydrogen, says Yabe. It is also highly abundant, with about 1.3 grams found in every liter of seawater, or about 1,800 trillion metric tons in our oceans, he says. Moreover, the magnesium oxide resulting from the reaction can be converted back into magnesium, says Yabe. The catch? Recycling the magnesium oxide back into magnesium requires temperatures of 4,000 kelvins (3,726 ºC)–hence the need for a laser to generate such temperatures on a small spot. But for a magnesium combustion engine to function as a practical source of energy, the lasers need to be powered by a renewable energy source, such as solar power. Solar-pumped lasers already exist: they work by concentrating sunlight onto crystalline materials such as neodymium-doped yttrium aluminium garnet, causing them to emit laser light. Until now, however, most solar-pumped lasers have relied on extremely large mirrors to focus the sunlight on the crystal. Yabe and his colleagues have developed a compact laser that offers a threefold improvement in efficiency over previous designs, in terms of how much power it can deliver compared with the available sunlight. This is partly due to the use of Nd:YAG crystals that are additionally doped with chromium, enabling them to absorb a broader range of light. Adding the chromium makes a greater proportion of the spectrum available, says Yabe: “Thus the efficiency from sunlight to laser is greatly enhanced.” The other innovation of Yabe’s laser is the use of a small Fresnel lens instead of large mirror lenses. Fresnel lenses reduce the size and amount of material needed to build a lens by breaking it into concentric rings of lenses. Typically, 10 percent of incident light is focused on the crystal, whereas with the Fresnel, it’s around 80 percent. “In our case, we used only 1.3 meter squared and achieved 25 watts,” says Yabe. Although this is only a threefold increase, the laser output exponentially increases with the increasing area. “So we are expecting 300 to 400 watts with the four-meter-squared Fresnel lens,” he says. It’s an unusual approach, says Sunita Satyapal, head of the Department of Energy’s hydrogen-storage team, in Washington, DC. But it’s not the first time that metals, such as magnesium, and water have been explored as a means of hydrogen production, she says. What is needed now is a total-efficiency budget for the entire system, says Satyapal: “The key issue is cost and total efficiency.” There are much simpler ways of generating hydrogen using sunlight, such as by employing solar cells to split water using electrolysis, she adds.

https://www.technologyreview.com/s/408698/solar-powered-laser/

Yellow fever is a viral infection spread by mosquitoes. Yellow fever is caused by a virus carried by mosquitoes. You can catch this disease if you are bitten by a mosquito infected with this virus. This disease is common in South America and in sub-Saharan Africa. Anyone can get yellow fever, but the elderly have a higher risk of severe infection. If a person is bitten by an infected mosquito, symptoms usually develop 3 - 6 days later. Yellow fever has three stages: - Stage 1 (infection): Headache, muscle and joint aches, fever, flushing, loss of appetite, vomiting, and jaundice are common. Symptoms often go away briefly after about 3 - 4 days. - Stage 2 (remission): Fever and other symptoms go away. Most people will recover at this stage, but others may get worse within 24 hours. - Stage 3 (intoxication): Problems with many organs may occur, including the heart, liver, and kidney. Bleeding disorders, seizures, coma, and delirium may also occur. Symptoms may include: Exams and Tests The health care provider will perform a physical examination and request selected blood tests. These blood tests may show liver and kidney failure and shock. It is important to tell your doctor if you have traveled to areas where the disease is known to thrive. Blood tests can confirm the diagnosis. There is no specific treatment for yellow fever. Treatment for symptoms can include: - Blood products for severe bleeding - Dialysis for kidney failure - Fluids through a vein (intravenous fluids) Yellow fever can cause severe problems, including internal bleeding. Death is possible. When to Contact a Medical Professional Get medical attention at least 10 - 14 days before traveling to an endemic area for yellow fever to find out whether you should be vaccinated against the disease. Tell your health care provider right away if you or your child develop fever, headache, muscle aches, vomiting, or jaundice, especially if you have traveled to an area where yellow fever is known to occur. Some countries require proof of vaccination to gain entry. If you will be traveling to an area where yellow fever is common: - Sleep in screened housing - Use mosquito repellents - Wear clothing that fully covers your body There is an effective vaccine against yellow fever. Ask your doctor at least 10 - 14 days before traveling if you should be vaccinated against yellow fever. Bausch DG. Viral hemorrhagic fevers. In: Goldman L, Schafer AI, eds. Cecil Medicine. 24th ed. Philadelphia, Pa: Saunders Elsevier; 2011:chap 389. Reviewed By: Jatin M. Vyas, MD, PhD, Assistant Professor in Medicine, Harvard Medical School; Assistant in Medicine, Division of Infectious Disease, Department of Medicine, Massachusetts General Hospital. Also reviewed by David Zieve, MD, MHA, Isla Ogilvie, PhD, and the A.D.A.M. Editorial team.

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Dinosaur Appearance and Behavior Dinosaurs with almost complete fossil skeletons give us clues about what they were like. We get a good idea of the dinosaurs' size, weight, and appearance in life. The cell structure of fossilized dinosaur bones can tell us about the biology of a dinosaur. We then have information about how rapidly they grew, and perhaps about whether they were warm-blooded or cold-blooded. Fossilized bones sometimes leave evidence of bone diseases and the tooth marks of predators. Muscle tissue is almost never preserved (only two good "dinosaur mummies," with soft parts intact, have been discovered). But we can still tell how dinosaurs moved from the traces of ligaments and muscle scars on the bones. Dinosaurs' teeth can tell us what kinds of foods they ate. Occasionally the actual stomach contents are preserved, so scientists can study what a dinosaur had for its last meal. The smooth "stomach stones" with which some dinosaurs ground up their food are sometimes preserved, and even fossilized dinosaur droppings, which are known as coprolites, have been found. All this gives us more information on dinosaur diets. Fossilized skin impressions can be seen on very well preserved skeletons. The fine-grained sandy rocks of Dinosaur Provincial Park in Alberta, Canada, have preserved the skin impressions of duckbilled and horned dinosaurs. These provide an idea of what dinosaur skin looked like. No dinosaur has ever been found with feather traces, so there is no evidence that any dinosaur had feathers. Scientists study locations and distribution of dinosaur bones for information about the dinosaurs' environments. The direction in which bones are pointed and the way they are arranged in dinosaur bone beds are clues about the size and strength of rivers dinosaurs had to wade through. Scientists compare the number of different types of dinosaurs from one location to find out about dinosaur habits and lifestyles. For example, if there were few predators and many animals that would have been their prey, it would show that the predatory dinosaurs had a quick metabolism and were probably warm-blooded.

http://animals.howstuffworks.com/dinosaurs/dinosaur-bones3.htm

Scientists from across the world came together in London on 12-13 January to review the scientific and technical status of the LISA mission, the world’s first gravitational wave observatory, at a meeting organised by the Royal Astronomical Society (RAS) and the Institute of Physics. Scheduled for launch in 2016, LISA will be the largest scientific instrument ever constructed, consisting of three spacecraft, each separated by 5 million kilometres (3 million miles). Its task will be to detect the elusive gravitational waves which were predicted by Einstein’s Theory of General Relativity, published in 1916. To date, although astronomers have indirect evidence of their existence, none have yet been detected directly. LISA will be one of the most challenging space science missions ever flown. In order to detect the passage of a gravitational wave, the distance between the spacecraft must be measured by laser beams to an accuracy of ten picometres, about one millionth of the diameter of a human hair! Gravitational waves are emitted when very massive objects such as black holes spiral violently together or when neutron stars collide at high speed. These invisible waves squeeze and stretch spacetime as they travel to us from distant parts of the universe, The waves travel from the source without absorption and this allows scientists to study objects at very great distances and the events that took place immediately after the birth of the Universe. Various models of the early universe predict gravitational wave emission during the first tiny fractions of a second, and if these can be detected by LISA scientists will learn a great deal about the processes active at that time. The technology needed for gravitational wave detection in space is being developed in Europe and the US, with a major role being played by the UK. Groups at the Universities of Glasgow, Birmingham, Imperial College London and the Rutherford Appleton Laboratory have been working for over ten years to perfect the necessary instrumentation and a flight test of this hardware is planned for 2009 on a space mission called LISA Pathfinder. Over a period of at least 2 years, LISA will detect gravitational waves from a variety of compact objects, ranging from massive black holes at great distances from the Earth to sub-solar mass white dwarfs – extremely dense, glowing remnants of dead stars - in our Galaxy. The mission will consist of three spacecraft flying 5 million kilometres (3 million miles) apart in an equilateral triangle formation. Laser beams traveling between the spacecraft will be reflected from two test masses in each satellite. By obtaining extremely accurate measurements of the distance between the spacecraft, it will be possible to determine whether the fabric of spacetime in which they are traveling is being distorted by passing gravitational waves. The formation of three spacecraft will face the Sun and lie in a plane that is tilted at 60 degrees to the Earth’s orbit. The trio will orbit the Sun, following 20 degrees behind the Earth, and will rotate once per year. This orbital motion will help to detect the direction of each source of gravitational radiation. Although LISA will not be affected by vibrations that influence ground-based observatories, the test masses must be cocooned within active shields to protect them from the constant buffeting by charged particles pouring out of the Sun. Sensors will detect the relative motion of the spacecraft and the delicate test mass mirrors, and will command thrusters to minimise the relative motion. Source: Royal Astronomical Society (RAS) Explore further: LISA pathfinder thrusters operated successfully

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Short division is similar to long division, but it involves less written work and more mental arithmetic. The general method for both short and long division is the same, but in short division you write down less of your work, doing the simple subtraction and multiplication mentally. To understand short division, you must have mastered the basic skills of subtraction and multiplication. Short division is ideal when the divisor, the number that you're dividing into another number, is less than 10. Doing Short Division 1Write the problem. To write the problem correctly, place the divisor, the number that you're dividing into another number, outside the long division bar. Place the dividend, the number that you'll be dividing by the divisor, inside the long division bar. The quotient, or your result, will go on top of the division bar. Remember that for short division to work, your divisor has to be less than 10. - For example: In 847/5, 5 is the divisor, so write it outside the division bar and 847 is the dividend, so place it inside the division bar. - The quotient is blank because you haven't started dividing yet. 2Divide the first number of the dividend by the divisor. In this case, 5 goes into 8 one time with 3 left over. Write the number 1, the first number of the quotient, on top of the division bar. This leftover number is called the remainder. - If you were using long division, you would write out 8-5 equals 3 and then bring down the 4. Short division simplifies this written process. - At the beginning of a problem the divisor may not go into first number of the dividend. For example, 567/7. In this case 7 doesn’t go into 5, but it does go into 56 eight times. When solving this problem, write the first number of the quotient over the 4 instead of the 5 and continue solving. The final answer is 81. - If you encounter a problem where the divisor does not go into the dividend, simply write a zero in the quotient, trying again with the next number in the dividend until the number can be divided. For example, 3208/8, 8 goes into 32 four times, but does not go into 0. You would add a 0 and then divide into the next number. 8 goes into 8 one time, therefore, the solution would be 401. 3Write the remainder next to the first number of the dividend. Write a small 3 to the top right of the number 8. This will remind you that there was a remainder of 3 when you divided 8 by 5. The next number you will divide into is the combination of the remainder and the second number. - In our example, the next number is 34. 4Divide the number formed by the first remainder and the second number in the dividend by the divisor. The remainder is 3 and the second number of the dividend is 4, so the new number you'll be working with is 34. - Now, divide 34 by 5. 5 goes into 34 six times (5 x 6 =30) with a remainder of 4. - Write your quotient, 6, on the division bar to the right of the 1. - Again, keep in mind you are doing most of the math mentally. 5Write the second remainder above the second number in the dividend and divide. Just as you did the first time, simply write a small 4 above and to the right of the number 4. The next number you will be dividing by is 47. - Now, divide 47 by 5. 5 goes into 47 9 times (5 x 9 = 45) with a remainder of 2. - Write your quotient, 9, on the division bar to the right of the 6. 6Write the final remainder on the division bar. Write "r 2" to the right of the quotient on the division bar. The final answer of 847/5 is 169 with a remainder 2. Questions and Answers Give us 3 minutes of knowledge! |A video on how to do short division. | In other languages: Português: Do Short Divisão, Italiano: Fare le Divisioni in Linea, Français: faire une division raccourcie, Deutsch: Kurz Division ausführen, 中文: 做短除法, Русский: выполнить ускоренное деление, Español: hacer una división corta, Bahasa Indonesia: Melakukan Pembagian Pendek, Nederlands: Snel delen, العربية: استخدام القسمة المختصرة, Tiếng Việt: Làm Phép Chia Ngắn Thanks to all authors for creating a page that has been read 257,261 times.

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The retina is a thin nerve membrane that detects light entering the eye. Nerve cells in the retina send signals of what the eye sees along the optic nerve to the brain. The retina lines the back two-thirds of the eye and is made up of two layers: the sensory retina and the retinal pigment epithelium (RPE). The macula, near the center of the retina at the back of the eyeball, provides the sharp, detailed, central vision a person uses for focusing on what is directly in the line of sight. The rest of the retina provides side (peripheral) vision, which lets a person see shapes but not fine details. eMedicineHealth Medical Reference from Healthwise To learn more visit Healthwise.org - Early Care for Your Premature Baby - What to Eat When You Have Cancer - When to Take More Pain Medication

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How to multiply and divide rational expressions. How to divide polynomials using synthetic division. How to rationalize the denominator when dealing with an imaginary number. How to use long division to divide polynomials. Dividing fractions by splitting into equal parts or multiplying by the reciprocal Multiplying or dividing with mixed numbers How to divide two complex numbers in trigonometric form. Looking at how multiplication represents repeated addition, as well as special cases of multiplying and dividing whole numbers How to divide polynomials using long division. How to rationalize denominators. How to simplify a complex fraction. How to identify the centroid and the way it divides each of the medians. Learn how to multiply and divide positive and negative numbers. Multiplying and dividing decimals by ignoring and later replacing the decimal point How to define the apothem and center of a polygon; how to divide a regular polygon into congruent triangles.

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Roman Numerals are numbers in a system. Surely everyone knows that I is 1, II is 2, III is 3, and IV is 4, but do they know that V is 5, or that X is 10? If you don't know that and the numbers that follow, and would like to know, read on! It will definitely help you when looking at copyright dates of books published in MCMXXXIV (1934) or to understand that the book on the Statue of Liberty bears a label MDCCLXXVI (1776)! 1Understand the concept of Roman Numerals. In Roman Numerals, the first number I is followed by II, then there is III. If there were not four parts in a quarter, then IV would not be known, and one might assume the next number is IIII. This is incorrect. The next number is IV, and following IV is V. After V is VI, VII, and VIII. After VIII comes IX, then X. The concept repeats itself continuously. This may sound confusing, though it is easy to understand if you look at Roman Numerals from 1 to 20: I II III IV V VI VII VIII IX X XI XII XIII XIV XV XVI XVII XVIII XIX XX 2Learn the values. As you may know, the Roman numbering system is based on letters rather than a whole new alphabet for numbers. - I 1 - V 5 - X 10 - L 50 - C 100 - D 500 - M 1000 3Put them together. The concept in step 1 is actually quite basic. That pattern repeats continuously throughout the system, no matter how great a number may be. For example: - From 40 to 50: XL XLI XLII XLIII XLIV XLV XLVI XLVIII XLIX L - From 100 to 110: C CI CII CIII CIV CV CVI CVII CVIII CIX CX - By hundreds, from 100 to 1000: C CC CCC CD D DC DCC DCCC CM M - From 40 to 50: 4Learn the special rules for 4000 and above. There is no character representing 5000, so there are 3 ways to do this. - The Romans themselves just wrote MMMM for 4000. - To be faster, write the Roman Numeral value for 4 in parentheses like so: (IV) - Finally, a bar could be drawn above the number rather than putting the number in parentheses. 5Understand how to break up a number to read it. After understanding how the individual numbers are written, it is easy to understand how MCMXXXIV (used as an example in the introduction) means 1934. To get MCMXXXIV, simply add: MCM (1900) + XXX (30) + IV (4). To get MDCCLXXVI (1776), simply add: MDCC (1700) + LXX (70) + VI (6). Questions and Answers Give us 3 minutes of knowledge! - Since it is likely that numbers over 100 will probably be used to mark years beyond the year 1000, try memorizing the numbers for 1000, 1100, 1200, 1300, and so on, up to 2000, so that you know which century each one marks, especially for 1700, 1800, and 1900.

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|This article does not cite any sources. (December 2009)| In thermodynamics, motive power is a natural agent, such as water or steam, wind or electricity, used to impart motion to machinery such as an engine. Motive power may also be a locomotive or a motor, which provides motive power to a system. Motive power may be thought of as a synonym for either "work", i.e. force times distance, or "power". In 1679, physicist Denis Papin conceived the idea of using steam to power a piston and cylinder engine, by watching a steam release valve of a bone-digester rhythmically move up and down. In 1698, based on Papin’s designs, mechanical designer Thomas Savery built the first engine. The first scientific treatise on the energetics of engines was the 1824 book: Reflections on the Motive Power of Fire written by French physicist Sadi Carnot. As an example, the Newcomen engine of 1711 was able to replace a team of 500 horses that had “powered” a wheel to pump water out of a mine, i.e. to “move” buckets of water vertically out of a mine. Hence, we have the precursory model to the term motive power. Based on this model, in 1832, Carnot defined work as “weight lifted through a height”, being the very same definition used to this day. Carnot states, in the footnotes to his famous 1824 publication, “We use here the expression motive power to express the useful effect that a motor is capable of producing. This effect can always be likened to the elevation of a weight to a certain height. It has, as we know, as a measure, the product of the weight multiplied by the height to which it is raised.” In this manner, Carnot is actually referring to "motive power" in the same manner we currently define "work". If we were to include a unit of time in Carnot's definition, we would then have the modern-day definition for power: Thus Carnot's definition of motive power is not consistent with the modern physics definition of "power", nor the modern usage of the term. In 1834, the French mining engineer Émile Clapeyron refers to Carnot’s motive power as “mechanical action”. As an example, during the expansion stroke of a piston engine he states that: “the gas will have developed a quantity of mechanical action during its expansion given by the integral of the product of the pressure times the differential of the volume.” Clapeyron then goes on to use graphical methods to show how this "mechanical action", i.e. work in modern terms, could be calculated.

https://en.wikipedia.org/wiki/Motive_power

Arab-Israeli History 101History Q & A During World War I, Britain made three different promises regarding historic Palestine. Arab leaders were assured that the land would become independent; in the Balfour declaration, Britain indicated its support for a Jewish national home in Palestine; and secretly Britain arranged with its allies to divide up Ottoman territory, with Palestine becoming part of the British empire. Historians have engaged in detailed exegesis of the relevant texts and maps, but the fundamental point is that Britain had no moral right to assign Palestine to anyone: by right Palestine belonged to its inhabitants. In the late years of the 19th century, anti-Semitism became especially virulent in Russia and re-emerged in France. Some Jews concluded that only in a Jewish state would Jews be safe and thus founded Zionism. Most Jews at the time rejected Zionism, preferring instead to address the problem of anti-Semitism through revolutionary or reformist politics or assimilation. And for many orthodox Jews, especially the small Jewish community in Palestine, a Jewish state could only be established by God, not by humans. At first Zionists were willing to consider other sites for their Jewish state, but they eventually focused on Palestine for its biblical connections. The problem, however, was that although a Zionist slogan called Palestine"a land without people for a people without land," the land was not at all empty. Following World War I, Britain arranged for the League of Nations to make Palestine a British"mandate," which is to say a colony to be administered by Britain and prepared for independence. To help justify its rule over Arab land, Britain arranged that one of its duties as the mandatory power would be to promote a Jewish national home. Who were the Jews who came to Palestine? The early Zionist settlers were idealistic, often socialist, individuals, fleeing oppression. In this respect they were like the early American colonists. But also like the American colonists, many Zionists had racist attitudes toward the indigenous people and little regard for their well-being.1 Some Zionists thought in terms of Arab-Jewish cooperation and a bi-national state, but many were determined to set up an exclusively Jewish state (though to avoid antagonizing the Palestinians, they decided to use the term Jewish"national home" rather than"state" until they were able to bring enough Jews to Palestine). Jewish immigration to Palestine was relatively limited until the 1930s,.when Hitler came to power. The U.S. and Europe closed their doors to immigration by desperate jews, making Palestine one of the few options. Who were the indigenous people of Palestine? Pro-Israel propaganda has argued that most Palestinians actually entered Palestine after 1917, drawn to the economic dynamism of the growing Jewish community, and thus have no rights to Palestine. This argument has been elaborated in Joan Peters' widely promoted book, From Time Immemorial. However, the book's claim is false.2 The indigenous population was mostly Muslim, with a Christian and a smaller Jewish minority. As Zionists arrived from Europe, the Muslims and Christians began to adopt a distinctly Palestinian national identity. How did the Zionists acquire land in Palestine? Some was acquired illegally and some was purchased from Arab landlords with funds provided by wealthy Jews in Europe. Even the legal purchases, however, were often morally questionable as they sometimes involved buying land from absentee landlords and then throwing the poor Arab peasants off the land. Land thus purchased became part of the Jewish National Fund which specified that the land could never be sold or leased to Arabs. Even with these purchases, Jews owned only about 6% of the land by 1947. Was Palestinian opposition to Zionism a result of anti-Semitism? Anti-Semitism in the Arab world was generally far less severe than in Europe. Before the beginning of Zionist immigration, relations among the different religious groups in Palestine were relatively harmonious. There was Palestinian anti-Semitism, but no people will look favorably on another who enter one's territory with the intention of setting up their own sovereign state. The expulsion of peasants from their land and the frequent Zionist refusal to employ Arabs exacerbated relations. What was the impact of World War II on the Palestine question? As World War II approached, Britain shrewdly calculated that they could afford to alienate Jews -- who weren't going to switch to Hitler's side -- but not Arabs, so they greatly restricted Jewish immigration into Palestine. But, of course, this was precisely when the need for sanctuary for Europe's Jews was at its height. Many Jews smuggled their way into Palestine as the United States and other nations kept their borders closed to frantic refugees. At the end of the war, as the enormity of the Holocaust became evident, for the first time Zionism became a majority sentiment among world Jewry. Many U.S. Christians also supported Zionism as a way to absolve their guilt for what had happened, without having to allow Jews into the United States. U.S. Zionists, who during the war had subordinated rescue efforts to their goal of establishing a Jewish state,3 argued that the Holocaust proved more than ever the need for a Jewish state: Had Israel existed in 1939, millions of Jews might have been saved. Actually, Palestine just narrowly avoided being overrun by the Nazis, so Jews would have been far safer in the United States than in a Jewish Palestine. During the war many Jews in Palestine had joined the British army. By war's end, the Jewish community in Palestine was well armed, well-organized, and determined to fight. The Palestinians were poorly armed, with feudal leaders. The Mufti of Jerusalem had been exiled by the British for supporting an Arab revolt in 1936-39 and had made his way to Berlin during the war where he aided Nazi propaganda. From the Zionist point of view, it was considered a plus to have the extremist Mufti as the Palestinians' leader; as David Ben Gurion, the leader of the Jewish community in Palestine and Israel's first prime minister, advised in 1938,"rely on the Mufti."4 What were the various positions in 1947? Both the Palestinians and the Zionists wanted the British out so they could establish an independent state. The Zionists, particularly a right-wing faction led by Menachim Begin, launched a terror campaign against Britain. London, impoverished by the war, announced that it was washing its hands of the problem and turning it over to the United Nations (though Britain had various covert plans for remaining in the region). The Zionists declared that having gone through one of the great catastrophes of modern history, the Jewish people were entitled to a state of their own, one into which they could gather Jewish refugees, still languishing in the displaced persons camps of Europe. The Zionist bottom line was a sovereign state with full control over immigration. The Palestinians argued that the calamity that befell European Jews was hardly their fault. If Jews were entitled to a state, why not carve it out of Germany? As it was, Palestine had more Jewish refugees than any other place on Earth. Why should they bear the full burden of atoning for Europe's sins? They were willing to give full civil rights (though not national rights) to the Jewish minority in an independent Palestine, but they were not willing to give this minority the right to control immigration, and bring in more of their co-religionists until they were a majority to take over the whole of Palestine. A small left-wing minority among the Zionists called for a binational state in Palestine, where both peoples might live together, each with their national rights respected. This view had little support among Jews or Palestinians. What did the UN do and why? In November 1947, the UN General Assembly voted to partition Palestine into two independent states, a Jewish state and an Arab state, joined by an economic union, with Jerusalem internationalized. In 1947 the UN had many fewer members than it does today. Most Third World nations were still colonies and thus not members. Nevertheless, the partition resolution passed only because the Soviet Union and its allies voted in favor and because many small states were subject to improper pressure. For example, members of the U.S. Congress told the Philippines that it would not get U.S. economic aid unless it voted for partition. Moscow favored partition as a way to reduce British influence in the region; Israel was viewed as potentially less pro-Western than the dominant feudal monarchies. Didn't Palestinians have a chance for a state of their own in 1947, but they rejected it by going to war with Israel? In 1947 Jews were only one third of the population of Palestine and owned only 6% of the land. Yet the partition plan granted the Jewish state 55% of the total land area. The Arab state was to have an overwhelmingly Arab population, while the Jewish state would have almost as many Arabs as Jews. If it was unjust to force Jews to be a 1/3 minority in an Arab state, it was no more just to force Arabs to be an almost 50% minority in a Jewish state. The Palestinians rejected partition. The Zionists accepted it, but in private Zionist leaders had more expansive goals. In 1938, during earlier partition proposals, Ben Gurion stated,"when we become a strong power after the establishment of the state, we will abolish partition and spread throughout all of Palestine."5 The Mufti called Palestinians to war against partition, but in fact very few Palestinians responded. The"decisive majority" of Palestinians, confided Ben Gurion,"do not want to fight us." The majority"accept the partition as a fait accompli," reported a Zionist Arab affairs expert. The 1936-39 Arab revolt against the British had mass popular support, but the 1947-48 fighting between the Mufti's followers and the Zionist military forces had no such popular backing.6 But even if Palestinians were fully united in going to war against the partition plan, this can provide no moral justification for denying them their basic right of self- determination for more than half a century. This right is not a function of this or that agreement, but a basic right to which every person is entitled. (Israelis don't lose their right to self-determination because their government violated countless UN cease-fire resolutions.) Didn't Israel achieve larger borders in 1948 as a result of a defensive war of independence? Arab armies crossed the border on May 15, 1948, after Israel declared its independence. But this declaration came three and a half months before the date specified in the partition resolution. The U.S. had proposed a three month truce on the condition that Israel postpone its declaration of independence. The Arab states accepted and Israel rejected, in part because it had worked out a secret deal with Jordan's King Abdullah, whereby his Arab Legion would invade the Palestinian territory assigned to the Palestinian state and not interfere with the Jewish state. (Since Jordan was closely allied to Britain, the scheme also provided a way for London to maintain its position in the region.) The other Arab states invaded as much to thwart Abdullah's designs as to defeat Israel.7 Most of the fighting that ensued took place on territory that was to be part of the Palestinian state or the internationalized Jerusalem. Thus, Israel was primarily fighting not for its survival, but to expand its borders at the expense of the Palestinians. For most of the war, the Israelis actually held both a quantitative and qualitative military edge, even apart from the fact that the Arab armies were uncoordinated and operating at cross purposes.8 When the armistice agreements were signed in 1949, the Palestinian state had disappeared, its territory taken over by Israel and Jordan, with Egypt in control of the Gaza Strip. Jerusalem, which was to have been internationalized, was divided between Israeli and Jordanian control. Israel now held 78% of Palestine. Some 700,000 Palestinians had become refugees. Why did Palestinians become refugees in 1948? The Israeli government claim is that Palestinians chose to leave Palestine voluntarily, instructed to do so via radio broadcasts from Arab leaders who wanted to clear a path for their armies. But radio broadcasts from the area were monitored by the British and American governments and no evidence of general orders to flee has ever been found. On the contrary, there are numerous instances of Arab leaders telling Palestinians to stay put, to keep their claim to the territory.9 People flee during wartime for a variety of reasons and that was certainly the case here. Some left because war zones are dangerous environments. Some because of Zionist atrocities -- most dramatically at Deir Yassin where in April 1948 254 defenseless civilians were slaughtered. Some left in panic, aided by Zionist psychological warfare which warned that Deir Yassin's fate awaited others. And some were driven out at gunpoint, with killings to speed them on their way, as in the towns of Ramle and Lydda.10 There is no longer any serious doubt that many Palestinians were forcibly expelled. The exact numbers driven out versus those who panicked or simply sought safety is still contested, but what permits us to say that all were victims of ethnic cleansing is that Israeli officials refused to allow any of them to return. (In Kosovo, any ethnic Albanian refugee, whether he or she was forced out at gunpoint, panicked, or even left to make it easier for NATO to bomb, was entitled to return.) In Israel, Arab villages were bulldozed over, citrus groves, lands, and property seized, and their owners and inhabitants prohibited from returning. Indeed, not only was the property of"absentee" Palestinians expropriated, but any Palestinians who moved from one place within Israel to another during the war were declared"present absentees" and their property expropriated as well. Of the 860,000 Arabs who had lived in areas of Palestine that became Israel, only 133,000 remained. Some 470,000 moved into refugee camps on the West Bank (controlled by Jordan) or the Gaza Strip (administered by Egypt). The rest dispersed to Lebanon, Syria, and other countries. Why did Israel expel the Palestinians? In part to remove a potential fifth column. In part to obtain their property. In part to make room for more Jewish immigrants. But mostly because the notion of a Jewish state with a large non-Jewish minority was extremely awkward for Israeli leaders. Indeed, because Israel took over some territory intended for the Palestinian state, there had actually been an Arab majority living within the borders of Israel. Nor was the idea of expelling Palestinians something that just emerged in the 1948 war. In 1937, Ben Gurion had written to his son,"We will expel the Arabs and take their places ... with the force at our disposal."11 How did the international community react to the problem of the Palestinian refugees? In December 1948, the General Assembly passed Resolution 194, which declared that"refugees wishing to return to their homes and live in peace with their neighbors should be permitted to do so" and that" compensation should be paid for the property of those choosing not to return." This same resolution was overwhelmingly adopted year after year. Israel repeatedly refused to carry out the terms of the resolution. Did the Arab countries take steps to resettle the Palestinian refugees? Only in Jordan were Palestinians eligible for citizenship. In Lebanon, the government feared that allowing Palestinians to become citizens would disturb the country's delicate Christian-Muslim balance; in Egypt, the shortage of arable land led the government to confine the Palestinians to the Gaza Strip. It must be noted, however, that the Palestinians were reluctant to leave the camps if that would mean acquiescing in the loss of homes and property or giving up their right to return. It is sometimes implied that the lack of assistance to Palestinians from Arab nations justifies Israel's refusal to acknowledge and address the claims of the refugees. But if you harm someone, you are responsible for redressing that harm, regardless of whether the victim's relatives are supportive. Hasn't there been a population exchange, with Jews from Arab lands coming to Israel and replacing the Palestinians? This argument makes individual Palestinians responsible for the wrong-doing of Arab governments. Jews left Arab countries under various circumstances: some were forced out, some came voluntarily, some were recruited by Zionist officials. In Iraq, Jews feared that they might be harmed, a fear possibly helped along by some covert bombs placed by Zionist agents.12 But whatever the case, there are no moral grounds for punishing Palestinians (or denying them their due) because of how Jews were treated in the Arab world. If Italy were to abuse American citizens, this would not justify the United States harming or expelling Italian-Americans. How were the Palestinians who remained within Israel treated? Most Arabs lived in the border areas of Israel and, until 1966, these areas were all declared military security zones, which essentially meant that Palestinians were living under martial law conditions for nearly 20 years. After 1966, Arab citizens of Israel continued to be the victims of harsh discrimination: most of the country's land is owned by the Jewish National Fund which prohibits its sale or lease to non-Jews; schools for Palestinians in Israel are, in the words of Human Rights Watch,"separate and unequal"; and government spending has been funneled so as to keep Arab villages underdeveloped. Thousands of Israeli Arabs live in villages declared"unrecognized" and hence ineligible for electricity or any other government services.13 Following 1948, didn't the Arab states continually try to destroy Israel? After Israel's victory in the 1948-49 war, there were several opportunities for peace. There was blame on all sides, but Israeli intransigence was surely a prime factor. In 1951, a UN peace plan was accepted by Egypt, Syria, Lebanon and Jordan, but rejected by Israel. When Nasser came to power in Egypt, he made overtures to Israel that were rebuffed. When Nasser negotiated an end to British control of the Suez Canal zone, Israeli intelligence covertly arranged a bombing campaign of western targets in Egypt as a way to discourage British withdrawal. The plot was foiled, Egypt executed some of the plotters, and Israel responded with a major military attack on Gaza.14 In 1956, Israel joined with Britain and France in invading Egypt, drawing condemnation from the United States and the UN. How were the Occupied Territories occupied? In June 1967, Israel launched a war in which it seized all of Palestine (the West Bank including East Jerusalem from Jordan and the Gaza Strip from Egypt), along with the Sinai from Egypt and the Golan Heights from Syria. Large numbers of Palestinians, some living in cities, towns, and villages, and some in refugee camps, came under Israeli control. (In 2001, half the Palestinian population of the Occupied Territories lived in refugee camps.15 The Israeli conquest also sent a new wave of refugees from Palestine to surrounding countries.) Israel's supporters argue that although Israel fired the first shots in this war, it was a justified preventive war, given that Arab armies were mobilizing on Israel's borders, with murderous rhetoric. The rhetoric was indeed blood-curdling, and many people around the world worried for Israel's safety. But those who understood the military situation -- in Tel Aviv and the Pentagon -- knew quite well that even if the Arabs struck first, Israel would prevail in any war. Nasser was looking for a way out and agreed to send his vice-president to Washington for negotiations. Israel attacked when it did in part because it rejected negotiations and the prospect of any face-saving compromise for Nasser. Menachem Begin, who was an enthusiastic supporter of this (and other) Israeli wars was quite clear about the necessity of launching an attack: In June 1967, he said, Israel"had a choice." Egyptian Army concentrations did not prove that Nasser was about to attack."We must be honest with ourselves. We decided to attack him."16 However, even if it were the case that the 1967 war was wholly defensive on Israel's part, this cannot justify the continued rule over Palestinians. A people do not lose their right to self-determination because the government of a neighboring state goes to war. Sure, punish Egypt and Jordan -- don't give them back Gaza and the West Bank (which they had no right to in the first place, having joined with Israel in carving up the stillborn Palestinian state envisioned in the UN's 1947 partition plan). But there is no basis for punishing the Palestinian population by forcing them to submit to foreign military occupation. Israel immediately incorporated occupied East Jerusalem into Israel proper, announcing that Jerusalem was its united and eternal capital. It then began to establish settlements in the Occupied Territories in violation of the Geneva Conventions which prohibit a conquering power from settling its population on occupied territory. These settlements, placed in strategic locations throughout the West Bank and Gaza were intended to" create facts" on the ground to make the occupation irreversible. How did the international community respond to the Israeli occupation? In November 1967, the UN Security Council unanimously passed resolution 242. The resolution emphasized"the inadmissibility of the acquisition of territory by war" and called for the"withdrawal of Israeli armed forces from territory occupied in the recent conflict." It also called for all countries in the region to end their state of war and to respect the right of each country"to live in peace within secure and recognized boundaries." Israel argued that because resolution 242 called for Israeli withdrawal from"territories," rather than"the territories," occupied in the recent conflict, it meant that Israel could keep some of them as a way to attain"secure" borders. The official French and Russian texts of the resolution include the definite article, but in any event U.S. officials told Arab delegates that it expected"virtually complete withdrawal" by Israel, and this was the view as well of Britain, France, and the Soviet Union.17 Palestinians objected to the resolution because it referred to them only in calling for"a just settlement to the refugee problem" rather than acknowledging their right to self- determination. By the mid-1970s, however, the international consensus -- rejected by Israel and the United States -- was expanded to include support for a Palestinian state in the West Bank and Gaza, perhaps with insignificant border adjustments. How did the United States respond to the Israeli occupation? Prior to the 1967 war, France, not the United States, was Israel's chief weapons supplier. But now U.S. officials determined that Israel would be an extremely valuable ally to have in the Middle East and Washington became Israel's principal military and diplomatic backer. Why, given the U.S. concern for Middle Eastern oil, was Washington supporting Israel? This assumes that the main conflict was Israel vs. the Arabs, rather than Israel and conservative, pro-Western Arab regimes vs. radical Arab nationalism. Egypt and Syria had been champions of the latter, armed by the Soviet Union, and threatening U.S. interests in the region. (On the eve of the 1967, for example, Egypt and Saudi Arabia were militarily backing opposite sides in a civil war in Yemen. Israel had plotted with Jordan against Palestinian nationalism in 1948, and in 1970 Israel was prepared to take Jordan's side in a war against Palestinians and Syria.) Diplomatically, the U.S. soon backed off the generally accepted interpretation of resolution 242, deciding that given Israel's military dominance no negotiations were necessary except on Israel's terms. So when Secretary of State Rogers put forward a reasonable peace plan, President Nixon privately sent word to Israel that the U.S. wouldn't press the proposal.18 When Anwar Sadat, Nasser's successor, proposed a peace plan that included cutting his ties with Moscow, Washington decided he hadn't groveled enough and ignored it. But after Egypt and Syria unsuccessfully went to war with Israel for the limited aim of regaining their lost territory, and Arab oil states called a limited oil embargo, Washington rethought its position. This led in 1979 to the Israeli-Egyptian Camp David Agreement under which Israel returned the Sinai to Egypt in return for peace and diplomatic relations. Egypt then joined Israel as a pillar of U.S. policy in the region and the two became the leading recipients of U.S. aid in the world. What progress was made toward justice for Palestinians during the first two decades of the occupation? The Palestine Liberation Organization was formed in 1964, but it was controlled by the Arab states until 1969, when Yasser Arafat became its leader. The PLO had many factions, advocating different tactics (some carried out hijackings) and different politics. At first the PLO took the position that Israel had no right to exist and that only Palestinians were entitled to national rights in Palestine. This was the mirror image of the official Israeli view -- of both the right-wing Likud party and the Labor party -- that there could be no recognition of the PLO under any circumstances, even if it renounced terrorism and recognized Israel, let alone acceptance of a Palestinian state on any part of the Occupied Territories. By 1976, however, the PLO view had come to accept the international consensus favoring a two-state solution. In January 1976 a resolution backed by the PLO, Egypt, Syria, Jordan, and the Soviet Union was introduced in the Security Council incorporating this consensus. Washington vetoed the resolution.19 The 1979 Camp David agreement established peace along the Egyptian-Israeli border, but it worsened the situation for Palestinians. With its southern border neutralized, Israel had a freer hand to invade Lebanon in 1982 (where the PLO was based) and to tighten its grip on the Occupied Territories. What was the first Intifada? Anger and frustration were growing in the Occupied Territories, fueled by iron-fisted Israeli repression, daily humiliations, and the establishment of sharply increasing numbers of Israeli settlements. In December 1987, Palestinians in Gaza launched an uprising, the Intifada, that quickly spread to the West Bank as well. The Intifada was locally organized, and enjoyed mass support among the Palestinian population. Guns and knives were banned and the main political demand was for an independent Palestinian state coexisting with Israel.20 Israel responded with great brutality, with hundreds of Palestinians killed. The Labor Party Defense Minister, Yitzhak Rabin, urged Israeli soldiers to break the bones of Palestinian demonstrators. PLO leader Khalil al-Wazir, who from Tunis had advised the rejection of arms, was assassinated (with the approval of Rabin); Israel was especially eager to repress Palestinian leaders who advocated a Palestinian state that would coexist with Israel.21 By 1989, the initial discipline of the uprising had faded, as a considerable number of individual acts of violence by Palestinians took place. Hamas, an organization initially promoted by the Israelis as a counterweight to the PLO,22 also gained strength; it called for armed attacks to achieve an Islamic state in all of Palestine. What were the Oslo Accords? Arafat had severely weakened his credibility by his flirtation with Saddam Hussein following the Iraqi invasion of Kuwait. (The Iraqi leader had opportunistically tried to link his withdrawal from Kuwait to an Israeli withdrawal from the Occupied Territories.) Israel saw Arafat's weakness as an opportunity. Better to deal with Arafat while he was weak, before Hamas gained too much influence. Let Arafat police the unruly Palestinians, while Israel would maintain its settlements and control over resources. The Oslo agreement consisted of"Letters of Mutual Recognition" and a Declaration of Principles. In Arafat's letter he recognized Israel's right to exist, accepted various UN resolutions, renounced terrorism and armed struggle. Israeli Prime Minister Rabin in his letter agreed to recognize the PLO as the representative of the Palestine people and commence negotiations with it, but there was no Israeli recognition of the Palestinian right to a state. The Declaration of Principles was signed on the White House lawn on September 13, 1993. In it, Israel agreed to redeploy its troops from the Gaza Strip and from the West Bank city of Jericho. These would be given self-governing status, except for the Israeli settlements in Gaza. A Palestinian Authority (PA) would be established, with a police force that would maintain internal order in areas from which Israeli forces withdrew. Left for future resolution in"permanent status" talks were all the critical and vexatious issues: Jerusalem, refugees, settlements, and borders. These talks were to commence by year three of the agreement. In September 1995 an interim agreement -- commonly called Oslo II -- was signed. This divided the Occupied Territories into three zones, Area A, Area B, and Area C. (No mention was made of a fourth area: Israeli-occupied East Jerusalem.) In area A, the PA was given civil and security control but not sovereignty; in area B the PA would have civil control and the Israelis security control; and area C was wholly under Israeli control (these included the settlements, the network of connecting roads, and most of the valuable land and water resources of the West Bank). In March 2000, 17% of the West Bank was designated area A -- where the vast majority of Palestinians lived -- 24% area B, and 59% area C. In the Gaza Strip, with a population of over a million Palestinians, 6,500 Israeli settlers lived in the 20% of the territory that made up area C. Palestinians thus were given limited autonomy -- not sovereignty -- over areas of dense population in the Gaza Strip and small, non-contiguous portions of the West Bank (there were 227 separate and disconnected enclaves),23 which meant that the PA was responsible chiefly for maintaining order over poor and angry Palestinians. How did Israel respond to the Oslo Accords? Whatever hopes Oslo may have inspired among the Palestinian population, most Israeli officials had an extremely restricted vision of where it would lead. In a speech in October 1995, Rabin declared that there would not be a return to the pre-1967 borders, Jerusalem would remain united and under exclusive Israeli sovereignty, and most of the settlements would remain under Israeli sovereignty. Rabin said he wanted the"entity" that Palestinians would get to be"less than a state."24 Under Rabin, settlements were expanded and he began a massive program of road-building, meant to link the settlements and carve up the West Bank. (These by-pass roads, built on confiscated Palestinian land and U.S.- funded, were for Israelis only.) In 1995, Rabin was assassinated by a right-wing Israeli and he was succeeded as prime minister by Shimon Peres. But Peres, noted his adviser Yossi Beilin, had an even more limited view than Rabin, wanting any future Palestinian state to be located only in Gaza.25 Yossi Sarid, head of the moderate left Israeli party Meretz, said that Peres's plan for the West Bank was"little different" from that of Ariel Sharon.26 Settlements and by-pass roads expanded further. In May 1996, Likud's Benjamin Netanyahu who was openly opposed to the Oslo accords was elected prime minister. Netanyahu reneged on most of the already agreed on Israeli troop withdrawals from occupied territory, continued building settlements and roads, stepped up the policy of sealing off the Palestinian enclaves, and refused to begin the final status talks required by Oslo.27 In 1999, Labor's Ehud Barak won election as prime minister. Barak had been a hardliner, but he had also confessed that if he had been born a Palestinian he probably would have joined a terrorist organization28 -- so his intentions were unclear. His policies, however, in his first year in office were more of the same: settlements grew at a more rapid pace than under Netanyahu, agreed-upon troops withdrawals were not carried out, and land confiscations and economic closures continued. His proposed 2001 government budget increased the subsidies supporting settlements in the Occupied Territories.29 What was the impact of the Oslo accords? The number of Israeli settlers since Oslo (1993) grew from 110,000 to 195,000 in the West Bank and Gaza; in annexed East Jerusalem, the Jewish population rose from 22,000 to 170,000.30 Thirty new settlements were established and more than 18,000 new housing units for settlers were constructed.31 From 1994-2000, Israeli authorities confiscated 35,000 acres of Arab land for roads and settlements.32 Poverty increased, so that in mid-2000, more than one out of five Palestinians had consumption levels below $2.10 a day.33 According to CIA figures, at the end of 2000, unemployment stood at 40%.34 Israeli closure policies meant that Palestinians had less freedom of movement -- from Gaza to the West Bank, to East Jerusalem, or from one Palestinian enclave to another -- than they had before Oslo.35 What was U.S. policy during this period? The United States has been the major international backer of Israel for more than three decades. Since 1976 Israel has been the leading annual recipient of U.S. foreign aid and is the largest cumulative recipient since World War II. And this doesn't include all sorts of special financial and military benefits, such as the use of U.S. military assistance for research and development in the United States. Israel's economy is not self-sufficient, and relies on foreign assistance and borrowing. During the Oslo years, Washington gave Israel more than $3 billion per year in aid, and $4 billion in FY 2000, the highest of any year except 1979. Of this aid, grant military aid was $1.8 billion a year since Oslo, and more than $3 billion in FY 2000, two thirds higher than ever before.36 Diplomatically, the U.S. retreated from various positions it had held for years. Since 1949, the U.S. had voted with the overwhelming majority of the General Assembly in calling for the right of return of Palestinian refugees. In 1994, the Clinton administration declared that because the refugee question was something to be resolved in the permanent status talks, the U.S. would no longer support the resolution. Likewise, although the U.S. had previously agreed with the rest of the world (and common sense) in considering East Jerusalem occupied territory, it now declared that Jerusalem's status too was to be decided in the permanent status talks. On three occasions in 1995 and 1997, the Security Council considered draft resolutions critical of Israeli expropriations and settlements in East Jerusalem; Washington vetoed all three.37 What happened at Camp David? Permanent status talks between Israel and the Palestinians as called for by the Oslo agreement finally took place in July 2000 at Camp David, in the United States, with U.S. mediators. The standard view is that Barak made an exceedingly generous offer to Arafat, but Arafat rejected it, choosing violence instead. A U.S. participant in the talks, Robert Malley, has challenged this view.38 Barak offered -- but never in writing and never in detail; in fact, says, Malley,"strictly speaking, there never was an Israeli offer" -- to give the Palestinians Israeli land equivalent to 1% of the West Bank (unspecified, but to be chosen by Israel) in return for 9% of the West Bank which housed settlements, highways, and military bases effectively dividing the West Bank into separate regions. Thus, there would have been no meaningfully independent Palestinian state, but a series of Bantustans, while all the best land and water aquifers would be in Israeli hands. Israel would also"temporarily" hold an additional 10 percent of West Bank land. (Given that Barak had not carried out the previous withdrawals to which Israel had committed, Palestinian skepticism regarding"temporary" Israeli occupation is not surprising.) It's a myth, Malley wrote,39 that"Israel's offer met most if not all of the Palestinians' legitimate aspirations" and a myth as well that the"Palestinians made no concession of their own." Some Israeli analysts made a similar assessment. For example, influential commentator Ze'ev Schiff wrote that, to Palestinians,"the prospect of being able to establish a viable state was fading right before their eyes. They were confronted with an intolerable set of options: to agree to the spreading occupation ... or to set up wretched Bantustans, or to launch an uprising."40 What caused the second Intifada? On September 28, 2000 Ariel Sharon, then a member of Parliament, accompanied by a thousand-strong security force, paid a provocative visit approved by Barak to the site of the Al Aqsa mosque. The next day Barak sent another large force of police and soldiers to the area and, when the anticipated rock throwing by some Palestinians occurred, the heavily-augmented police responded with lethal fire, killing four and wounding hundreds. Thus began the second Intifada. The underlying cause was the tremendous anger and frustration among the population of the Occupied Territories, who saw things getting worse, not better, under Oslo, whose hopes had been shattered, and whose patience after 33 years of occupation had reached the boiling point. Who is Ariel Sharon? Sharon was the commander of an Israeli force that massacred some seventy civilians in the Jordanian village of Qibya in 1953. He was Defense Minister in 1982, when Israel invaded Lebanon, causing the deaths of 17,000 civilians. In September 1982, Lebanese forces allied to Israel slaughtered hundreds of Palestinian non- combatants in the Sabra and Shitila refugee camps, a crime for which an Israeli commission found Sharon to bear indirect responsibility. As Housing Minister in various Israeli governments, Sharon vigorously promoted the settlements in the Occupied Territories. In January 2001, he took office as Prime Minister. How did Israel respond to this second Intifada? Israeli security forces responded to Palestinian demonstrations with lethal force even though, as a UN investigation reported, at these demonstrations the Israeli Defense Forces,"endured not a single serious casualty."41 Some Palestinians proceeded to arm themselves, and the killing escalated, with deaths on both sides, though the victims were disproportionately Palestinians. In November 2001, there was a week-long lull in the fighting. Sharon then ordered the assassination of Hamas leader Mahmoud Abu Hanoud, which, as everyone predicted, led to a rash of terror bombings, which in turn Sharon used as justification for further assaults on the PA.42 By March 2002, Amnesty International reported that more than 1000 Palestinians had been killed."Israeli security services have killed Palestinians, including more than 200 children, unlawfully, by shelling and bombing residential areas, random or targeted shooting, especially near checkpoints and borders, by extrajudicial executions and during demonstrations."43 Palestinian suicide bombings have targeted civilians. Amnesty International commented:"These actions are shocking. Yet they can never justify the human rights violations and grave breaches of the Geneva Conventions which, over the past 18 months, have been committed daily, hourly, even every minute, by the Israeli authorities against Palestinians. Israeli forces have consistently carried out killings when no lives were in danger." Medical personnel have been attacked and ambulances, including those of the Red Cross,"have been consistently shot at."44 Wounded people have been denied medical treatment. Israel has carried out targeted assassinations (sometimes the targets were probably connected to terrorism, sometimes not,45 but all of these extrajudicial executions have been condemned by human rights groups). The Israeli government criticized Arafat for not cracking down harder on terrorists and then responded by attacking his security forces, who might have allowed him to crack down, and restricting him to his compound in Ramallah. Israeli opinion became sharply polarized. At the same time that hundreds of military reservists have declared their refusal to serve in the West Bank and Gaza (www.couragetorefuse.org), polls show 46% of Israelis favor forcibly expelling all Palestinians from the Occupied Territories.46 What has U.S. policy been? U.S. military, economic, and diplomatic support has made possible the Israeli repression of the previous year and a half. Much of the weaponry Israel has been using in its attacks on Palestinians either was made in the United States (F-16s, attack helicopters, rockets, grenade launchers, Caterpillar bulldozers, airburst shells, M-40 ground launchers) or made in Israel with U.S. Department of Defense research and development funding (the Merkava tank). On March 26, 2001, the Security Council considered a resolution to establish an international presence in the Occupied Territories as a way to prevent human rights violations. The United States vetoed the resolution. Because Israel did not want the U.S. to get involved diplomatically, Washington did not name a special envoy to the region, General Zinni, until November 2001, more than a year after the Intifada began. Bush met four times with Sharon during the Intifada, never with Arafat. In February 2002, Vice President Cheney declared that Israel could"hang" Arafat.47 What caused the current crisis? As the Arab League was meeting to endorse a Saudi peace proposal -- recognition of Israel in return for full Israeli withdrawal to the 1967 borders -- a Hamas suicide bomber struck. Sharon, no doubt fearing a groundswell of support for the Arab League position, responded with massive force, breaking into Arafat's compound, confining him to several rooms. Then there were major invasions of all the Palestinian cities in the West Bank. There are many Palestinian casualties, though because Israel has kept reporters out, their extent is not known. In the early days of Sharon's offensive, Bush pointedly refused to criticize the Israeli action, reserving all his condemnation for Arafat, who, surrounded in a few rooms, was said to not be doing enough to stop terrorism. As demonstrations in the Arab world, especially in pro-U.S. Jordan and Egypt, threatened to destabilize the entire region, Bush finally called on Israel to withdraw from the cities. Sharon, recognizing that the U.S."demand" was uncoupled from any threat of consequences, kept up his onslaught. Is there a way out? A solution along the lines of the international consensus -- Israeli withdrawal from territories occupied in 1967, the establishment of a truly independent and viable Palestinian state in the West Bank and Gaza with its capital in East Jerusalem -- remains feasible. It needs only the backing of the United States and Israel. Don't the Arabs already have 22 states? Why do they need another one? Not all Arabs are the same. That other Arabs may already have their right of self- determination does not take away from Palestinians' basic rights. The fact that many Palestinians live in Jordan and have considerable influence and rights there, doesn't mean that the millions of Palestinians living under Israeli occupation or who were expelled from their homes and are now in refugee camps aren't entitled to their rights -- any more than the fact that there are a lot of Jews in the U.S., where they have considerable influence and rights, means that Israeli Jews should be packed off across the Atlantic. How can terrorists be given a state? If people whose independence movements use terrorism are not entitled to a state, then many current-day states would be illegitimate, not the least of them being Israel, whose independence struggle involved frequent terrorism against civilians. Won't an independent Palestinian state threaten Israeli security? Conquerors frequently justify their conquests by claiming security needs. This was the argument Israel gave for years why it couldn't return the Sinai to Egypt or pull out of Lebanon. Both of these were done, however, and Israel's security was enhanced rather than harmed. True, the Oslo Accords, which turned over disconnected swatches of territory to Palestinian administration, may not have improved Israeli security. But as Shimon Peres, one of the architects of the Oslo agreement and Sharon's current Foreign Minister acknowledged, Oslo was flawed from the start."Today we discover that autonomy puts the Palestinians in a worse situation." The second Intifada could have been avoided, Peres said, if the Palestinians had had a state from the outset."We cannot keep three and a half million Palestinians under siege without income, oppressed, poor, densely populated, near starvation."48 Israel is the region's only nuclear power. Beyond that, it is the strongest military power in the Middle East. Surely it cannot need to occupy neighboring territory in order to achieve security. Nothing would better guarantee the Israeli people peace and security than pulling out of the Occupied Territories. Isn't the Palestinian demand for the right of return just a ploy to destroy Israel? Allowing people who have been expelled from their homes the right to return is hardly an extreme demand. Obviously this can't mean throwing out people who have been living in these homes for many years now, and would need to be carefully worked out. Both Palestinian officials and the Arab League have indicated that in their view the right of return should be implemented in a way that would not create a demographic problem for Israel.49 Of course, one could reasonably argue that an officially Jewish state is problematic on basic democratic grounds. (Why should a Jew born in Brooklyn have a right to"return" to Israel while a Palestinian born in Haifa does not?) In any event, however, neither the Arab League nor Arafat have raised this objection.50 Don't Palestinians just view their own state as the first step in eliminating Israel entirely? Hamas and a few other, smaller Palestinian groups object not just to the occupation but to the very existence of Israel. But the Hamas et al. position is a distinctly minority sentiment among Palestinians, who are a largely secular community that has endorsed a two-state settlement. To be sure, Hamas has been growing in strength as a result of the inability of the Palestinian Authority to deliver a better life for Palestinians. If there were a truly independent Palestinian state, one can assume that Hamas would find far fewer volunteers for its suicide squads. It must be acknowledged, though, that the longer the mutual terror continues, the harder it will be to achieve long term peace. Is a two-state solution just? There is a broad international consensus on a two-state solution, along the lines of the Saudi peace proposal. Such a solution is by no means ideal. Palestine is a small territory to be divided into two states; it forms a natural economic unit. An Israeli state that discriminates in favor of Jews and a Palestinian state that will probably be equally discriminatory will depart substantially from a just outcome. What's needed is a single secular state that allows substantial autonomy to both national communities, something along the lines of the bi-national state proposed before 1948. This outcome, however, does not seem imminent. A two-state solution may be the temporary measure that will provide a modicum of justice and allow Jews and Palestinians to move peacefully forward to a more just future. - As Zionist writer Ahad Ha'am put it, his fellow Jews"treat the Arabs with hostility and cruelty, deprive them of their rights, offend them without cause, and even boast of these deeds." Quoted in Jews For Justice in The Middle East, The Origin of the Palestine- Israeli conflict, 3rd ed., P.O. Box 14561, Berkeley, CA, 94712, available at http://www.cactus48.com/truth.html. return - . Norman G. Finkelstein,"A Land Without a People: Joan Peters's 'Wilderness' Myth," in Image and Reality of the Israel Palestine Conflict, New York: Verso, 1995, pp. 21-50. return - See the sources cited by Noam Chomsky, Fateful Triangle: The United States, Israel and the Palestinians, updated edition, Cambridge: South End Press, 1999, p. 169n10. return - Simha Flapan, The Birth of Israel: Myths and Realities, New York: Pantheon, 1987, pp. 66-67. return - Quoted in Jerome Slater,"What Went Wrong? The Collapse of the Israeli-Palestinian Peace Process," Political Science Quarterly, vol. 116, no. 2, 2001, p. 174. return - Flapan, pp. 55, 73-77. return - Flapan, pp. 153-86. return - Flapan, pp. 187-199. return - Christopher Hitchens,"Broadcasts," in Blaming the Victims: Spurious Scholarship and the Palestinian Question, ed. Edward W. Said and Christopher Hitchens, New York: Verso, 1988, pp. 73-83. return - Benny Morris, The Birth of the Palestinian Refugee Problem, 1947-1949. New York: Cambridge University Press, 1987; Norman G. Finkelstein,"'Born of War, Not By Design," in Finkelstein, Image and Reality..., pp. 51-87. return - Slater, pp. 173-74. return - See Mark Tessler, A History of the Israeli-Palestinian Conflict, Bloomington: Indiana University Press, 1994, pp. 308-11; and sources in Noam Chomsky, Towards a New Cold War, New York: Pantheon, 1982, p. 462n33. return - Ian Lustick, Arabs in the Jewish State: Israel's Contorl of a National Minority, University of Texas, 1980; Human Rights Watch, Second Class: Discrimination Against Palestinian Arab Children in Israel's Schools, Sept. 2001, http://www.hrw.org/reports/2001/israel2/. On Israeli-Arab"unrecognized" villages, where some 100,000 people are forced to live without basic government services, including electricity and water, see http://www.assoc40.org/index_main.html. return - Charles D. Smith, Palestine and the Arab-Israeli Conflict, 4th ed., Boston: Bedford/St. Martin's, 2001, pp. 237-38. return - John Dugard, Kamal Hossain, and Richard Falk,"Question of The Violation of Human Rights in The Occupied Arab Territories, Including Palestine," Report of the human rights inquiry commission established pursuant to Commission resolution S-5/1 of 19 October 2000, E/CN.4/2001/121, 16 March 2001, para 29. return - Quoted in Chomsky, Fateful Triangle, p. 100. return - Smith, pp. 306, 334n10. return - Henry Kissinger, White House Years, Boston: Little, Brown, 1979, p. 376. return - Chomsky, Fateful Triangle, chap 3, esp. p. 67. return - Smith, pp. 418-21. return - Smith, pp. 422-24. return - Richard Sale,"Israel gave major aid to Hamas," UPI, Feb. 24, 2001. return - Geoffrey Aronson,"Recapitulating the Redeployments: The Israel-PLO 'Interim Agreements'," Information Brief No. 32, Center for Policy Analysis, 27 April 2000. return - Slater, p. 177, citing speech to Knesset of 5 October 1995, printed in Report on Israeli Settlement in the Occupied Territories 5 (November 1995). return - Slater, p. 178n9, quoting Ha'aretz, 7 March 1997. return - Slater, p. 178n9, quoting Report of the American Academy of Arts and Sciences, Israeli-Palestinian Security,1995. return - Slater, p. 179. return - Smith, p. 490. return - Slater, pp. 180-81. return - Edward Said,"Palestinians under Siege," in The New Intifada: Resisting Israel's Apartheid, ed. Roane Carey, New York: Verso, 2001, p. 29; Allegra Pacheco,"Flouting Convention: The Oslo Agreements," in Carey, p. 189. return - Sara Roy,"Decline and Disfigurement: The Palestinian Economy After Oslo," in Carey, p. 95; Pacheco, p. 187. return - Roy, p. 95. return - Roy, p. 101. return - CIA World Factbook 2001. return - Roy, pp. 98-100. return - Clyde R. Mark, Israel: U.S. Foreign Assistance, Updated March 15, 2002, CRS Issue Brief for Congress, Congressional Research Service, The Library of Congress, Order Code IB85066. Available at http:///www.fpc.gov/CRS_REPS/Crs_abs.htm. return - See the list of vetoed Security Council resolutions on Palestine at http://www.un.org/Depts/dpa/qpal/index.html. return - Robert Malley and Hussein Agha,"Camp David: The Tragedy of Errors," New York Review of Books, August 9, 2001. See also Deborah Sontag,"Quest for Mideast Peace: How and Why It Failed," New York Times, 26 July 2001, p. A1; and the critique of the Barak offer on the website of the"Peace Bloc," Gush Shalom, http:// comments powered by Disqus R Chand - 1/14/2004 Stephen Shalom's analysis is right on the mark N/A - 10/13/2003 I think the website is sooooooooooooo BORING!!! Ya'll are a bunch of LOSERS!!! Kenneth Stow - 6/17/2002 No, it is not propaganda, it is just bad history. I feel sorry for the students in this supposed course. The essay may overstate, for example, to say that most Arabs arrived in Palestine after 1917, but it is simply ignorant to imply that the land was thoroughly populated by Arabs from time immemorial, and to say that a Palestinian culture developed from 1917, which is at best a product of the later 1930s, and possibly later. In short, the Shalom essay is anything but irenic. It is shallow and as a graduate paper, it would get an F. From a professional historian, it is unacceptable. Shalom knows, or ought to know, the rules of evidence. Footnote: anti-Semitic. The term has nothing to do with origins. It refers to an attitude toward Jews whose basis is more than religious and which often stigmatizes Jews precisely for being Jews. A good source is John Gager's book on anti-Semistism's origins to learn the distinction between anti-Jewish and anti-Semitic as the terms are used in scholarly discourse. Sol Shalit - 6/14/2002 I substantially agree with Vought. Perhaps not propaganda, but certainly not history. This is a very biased and quite distorted -- in favor of Arabs -- reading of the Arab-Israeli conflict. All the footnotes will not save it; their purpose must have been to give the article the looks of objective scholarship, when in fact it is a not-so-subtle attempt to push a biased point of view, if not an agenda. I disagree with Vought that this entire web site is tainted becasue of this bad apple. In fact, I find the entire web site useful and quite balanced. I have seen good articles favoring one side or the other, and I have no objections as long as it was carefully researched and properly presented without obvious bias or agenda. It is still true that I found this bad apple offensive and believe that Vought has gone a bit too far in damning the entire site. He must have been angry, and rightly so. But the anger is against the author, not the web site. Finally, a note to the web site. You must, however, bear a bit of responsibility also. You could have recognized the obvious bias and agenda, by the tone of the first paragraph ("Israeli propaganda"). I'm not advocating censorship; just high standards. So far, at least in my book, you've done well. Keep it up and be on guard! Professor Emeritus of Economics Jim Williams - 6/13/2002 Since Semitic is a linguistic term for a related family of languages including Hebrew, Arabic, Aramaic, etc, both Palestinians and Hebrew-speaking Jews are Semites since they speak Semitic languages. Thus we should not use the term anti-Semitic but anti-Jewish of Palestinians who oppose Jews. To say Palestinians are anti-Semitic means that they hate themselves as well as Jews. Hans P. Vought - 6/12/2002 Real historians attempt to present an impartial, balanced account which fairly incorporates all available evidence. Like most of the rest of the pieces on this website, this is not history. but rather propaganda. One would never know from reading this article that Palestinians had ever killed anyone. Especially egregious is the author's attempt to make the intifada sound like nonviolent civil disobediance. Gandhi and King never threw rocks! I propose that the name of this website be changed to the Competing Propaganda Network. History has nothing to do with it. Pierre S. Troublion - 6/11/2002 Haven't had time to wade through all the footnotes yet, but it's nice to see some real history, instead of the usual Likud propaganda masquerading as Mideast history, on this website (for a change). - Rubio Surges Into Second In New Hampshire - Branstad Says Cruz Ran ‘Unethical’ Campaign - Christie Highlights Santorum’s Endorsement of Rubio - Portman Comes Out Against Trade Deal - Megyn Kelly Gets a Book Deal - A Big List of the Bad Things Clinton Has Done - An Unambiguous Sign Sanders Won Last Night’s Debate - Still Friends at the End - Quote of the Day - Trump Still Leads as Clinton Slips - Clinton Can’t Shake Image as Wall Street’s Friend - Maddow Doesn’t See Sanders Winning - Why Does the Media Still Shield Chelsea Clinton? - Bush Jokes His Mother May Have Abused Him - Rubio Closes the Gap in New Hampshire - Newly released interactive map shows images of destroyed monuments of Mosul - How the Rise of the Post Office Explains American Innovation - These Americans are reliving history and don’t mind repeating it - Britain largest home is saved for the nation - Shelter and the slums: capturing bleak Britain 50 years ago - WSJ features an article by a conservative calling for the abolition of Black History Month - Mary Beard, herself a bestselling author, wonders why more women historians aren't - Princeton U. historian Imani Perry claims mistreatment in parking ticket arrest - Retired historian George Dennison remains on the payroll at the U. of Montana while faculty are cut - The Atlantic profiles exciting ways to teach history

http://historynewsnetwork.org/article/783

Two motorboats travelling up and down a lake at constant speeds leave opposite ends A and B at the same instant, passing each other, for the first time 600 metres from A, and on their return, 400. . . . A cyclist and a runner start off simultaneously around a race track each going at a constant speed. The cyclist goes all the way around and then catches up with the runner. He then instantly turns. . . . Imagine starting with one yellow cube and covering it all over with a single layer of red cubes, and then covering that cube with a layer of blue cubes. How many red and blue cubes would you need? Choose any two numbers. Call them a and b. Work out the arithmetic mean and the geometric mean. Which is bigger? Repeat for other pairs of numbers. What do you notice? Square numbers can be represented as the sum of consecutive odd numbers. What is the sum of 1 + 3 + ..... + 149 + 151 + 153? What would be the smallest number of moves needed to move a Knight from a chess set from one corner to the opposite corner of a 99 by 99 square board? Choose a couple of the sequences. Try to picture how to make the next, and the next, and the next... Can you describe your reasoning? Three circles have a maximum of six intersections with each other. What is the maximum number of intersections that a hundred circles If you can copy a network without lifting your pen off the paper and without drawing any line twice, then it is traversable. Decide which of these diagrams are traversable. Imagine an infinitely large sheet of square dotty paper on which you can draw triangles of any size you wish (providing each vertex is on a dot). What areas is it/is it not possible to draw? Imagine a large cube made from small red cubes being dropped into a pot of yellow paint. How many of the small cubes will have yellow paint on their faces? Some students have been working out the number of strands needed for different sizes of cable. Can you make sense of their solutions? Draw a square. A second square of the same size slides around the first always maintaining contact and keeping the same orientation. How far does the dot travel? Use the animation to help you work out how many lines are needed to draw mystic roses of different sizes. How could Penny, Tom and Matthew work out how many chocolates there are in different sized boxes? Triangle numbers can be represented by a triangular array of squares. What do you notice about the sum of identical triangle numbers? What size square corners should be cut from a square piece of paper to make a box with the largest possible volume? Can you describe this route to infinity? Where will the arrows take you next? Euler discussed whether or not it was possible to stroll around Koenigsberg crossing each of its seven bridges exactly once. Experiment with different numbers of islands and bridges. Points P, Q, R and S each divide the sides AB, BC, CD and DA respectively in the ratio of 2 : 1. Join the points. What is the area of the parallelogram PQRS in relation to the original rectangle? Can you dissect a square into: 4, 7, 10, 13... other squares? 6, 9, 12, 15... other squares? 8, 11, 14... other squares? Can you cross each of the seven bridges that join the north and south of the river to the two islands, once and once only, without retracing your steps? The opposite vertices of a square have coordinates (a,b) and (c,d). What are the coordinates of the other vertices? The picture illustrates the sum 1 + 2 + 3 + 4 = (4 x 5)/2. Prove the general formula for the sum of the first n natural numbers and the formula for the sum of the cubes of the first n natural. . . . Can you see how this picture illustrates the formula for the sum of the first six cube numbers? If all the faces of a tetrahedron have the same perimeter then show that they are all congruent. Which hexagons tessellate? A huge wheel is rolling past your window. What do you see? Show that among the interior angles of a convex polygon there cannot be more than three acute angles. Some puzzles requiring no knowledge of knot theory, just a careful inspection of the patterns. A glimpse of the classification of knots and a little about prime knots, crossing numbers and. . . . To avoid losing think of another very well known game where the patterns of play are similar. The triangle OMN has vertices on the axes with whole number co-ordinates. How many points with whole number coordinates are there on the hypotenuse MN? How many moves does it take to swap over some red and blue frogs? Do you have a method? A tilted square is a square with no horizontal sides. Can you devise a general instruction for the construction of a square when you are given just one of its sides? Show that all pentagonal numbers are one third of a triangular number. Build gnomons that are related to the Fibonacci sequence and try to explain why this is possible. On the graph there are 28 marked points. These points all mark the vertices (corners) of eight hidden squares. Can you find the eight Watch these videos to see how Phoebe, Alice and Luke chose to draw 7 squares. How would they draw 100? Jo made a cube from some smaller cubes, painted some of the faces of the large cube, and then took it apart again. 45 small cubes had no paint on them at all. How many small cubes did Jo use? Can you find a rule which connects consecutive triangular numbers? The aim of the game is to slide the green square from the top right hand corner to the bottom left hand corner in the least number of Three frogs hopped onto the table. A red frog on the left a green in the middle and a blue frog on the right. Then frogs started jumping randomly over any adjacent frog. Is it possible for them to. . . . Can you maximise the area available to a grazing goat? A blue coin rolls round two yellow coins which touch. The coins are the same size. How many revolutions does the blue coin make when it rolls all the way round the yellow coins? Investigate for a. . . . A game for 2 players Here are four tiles. They can be arranged in a 2 by 2 square so that this large square has a green edge. If the tiles are moved around, we can make a 2 by 2 square with a blue edge... Now try to. . . . In a right angled triangular field, three animals are tethered to posts at the midpoint of each side. Each rope is just long enough to allow the animal to reach two adjacent vertices. Only one animal. . . . Is it possible to rearrange the numbers 1,2......12 around a clock face in such a way that every two numbers in adjacent positions differ by any of 3, 4 or 5 hours? A 2 by 3 rectangle contains 8 squares and a 3 by 4 rectangle contains 20 squares. What size rectangle(s) contain(s) exactly 100 squares? Can you find them all? In a three-dimensional version of noughts and crosses, how many winning lines can you make?

http://nrich.maths.org/public/leg.php?code=-68&cl=3&cldcmpid=650

Late Jurassic–Late Cretaceous, 160–66 Ma |Mounted skeletal cast of Troodon inequalis, Perot Museum| Elopterygidae? Lambrecht, 1933 Troodontidae is a group of bird-like theropod dinosaurs. During most of the 20th century, troodontid fossils were few and scrappy and they have therefore been allied, at various times, with many dinosaurian lineages. More recent fossil discoveries of complete and articulated specimens (including specimens which preserve feathers, eggs and embryos, and complete juveniles), have helped to increase understanding about this group. Anatomical studies, particularly studies of the most primitive troodontids, like Sinovenator, demonstrate striking anatomical similarities with Archaeopteryx and primitive dromaeosaurids, and demonstrate that they are relatives comprising a clade called Paraves. Troodontids are a group of small, bird-like, gracile maniraptorans. All troodontids have unique features of the skull, such as large numbers of closely-spaced teeth in the lower jaw. Troodontids have sickle-claws and raptorial hands, and some of the highest non-avian encephalization quotients, suggesting that they were behaviourally advanced and had keen senses. The largest troodontid was Troodon, and the smallest was Anchiornis, which is also the smallest known non-avian dinosaur. They had unusually long legs compared to other theropods, with a large, curved claw on their retractable second toes, similar to the "sickle-claw" of the dromaeosaurids. However, the sickle-claws of troodontids were not as large or recurved as in their relatives, and in some instances could not be held off the ground and "retracted" to the same degree. In at least one troodontid, Borogovia, the second toe could not be held far off the ground at all and the claw was straight, not curved or sickle-like. Troodontids had unusually large brains among dinosaurs, comparable to those of living flightless birds. Their eyes were also large, and pointed forward, indicating that they had good binocular vision. The ears of troodontids were also unusual among theropods, having enlarged middle ear cavities, indicating acute hearing ability. The placement of this cavity near the eardrum may have aided in the detection of low-frequency sounds. Troodontid ears were also asymmetrical, with one ear placed higher on the skull than the other, a feature shared only with some owls. The specialization of the ears may indicate that troodontids hunted in a manner similar to owls, using their hearing to locate small prey. Although most paleontologists believe that they were predatory carnivores, the many small, coarsely serrated teeth, large denticle size, and U-shaped jaws of some species (particularly Troodon) suggest that some species may have been omnivorous or herbivorous. Some suggest that the large denticle size is reminiscent of the teeth of extant iguanine lizards. In contrast, a few species, such as Byronosaurus, had large numbers of needle-like teeth, which seem best-suited for picking up small prey, such as birds, lizards and small mammals. Other morphological characteristics of the teeth, such as the detailed form of the denticles and the presence of blood grooves, also seem to indicate carnivory. Though little is known directly about the predatory behavior of troodontids, Fowler and colleagues theorize that the longer legs and smaller sickle claws (as compared to dromaeosaurids) indicates a more cursorial lifestyle, though the study indicates that troodontids were still likely to have used the unguals for prey manipulation. The proportions of the metatarsals, tarsals and unguals of troodontids appear indicative of their having nimbler, but weaker feet, perhaps better adapted for capturing and subduing smaller prey. This suggests an ecological separation from the slower but more powerful Dromaeosauridae. Many troodontid nests, including eggs that contain fossilized embryos, have been described. Hypotheses about troodontid reproduction have been developed from this evidence (see Troodon). A few troodont fossils, including specimens of Mei and Sinornithoides, demonstrate that these animals roosted like birds, with their heads tucked under their forelimbs. These fossils, as well as numerous skeletal similarities to birds and related feathered dinosaurs, support the idea that troodontids probably bore a bird-like feathered coat. The discovery of several fully feathered, primitive troodontids (Jinfengopteryx and Anchiornis) lend support to this. In 2004, Mark Norell and colleagues described two partial troodontid skulls (specimen numbers IGM 100/972 and IGM 100/974) found in a nest of oviraptorid eggs in the Djadokhta Formation of Mongolia. The nest is quite certainly that of an oviraptorosaur, since an oviraptorid embryo is still preserved inside one of the eggs. The two partial troodontid skulls were first described by Norell et al. (1994) as dromaeosaurids, but reassigned to the troodontid Byronosaurus after further study. The troodontids were either hatchlings or embryos, and fragments of eggshell are adhered to them although it seems to be oviraptorid eggshell. The presence of tiny troodontids in an oviraptorid nest is an enigma. Hypotheses explaining how they came to be there include that they were the prey of the adult oviraptorid, that they were there to prey on oviraptorid hatchlings, or that some troodontids may have been nest parasites. Troodontids and bird evolution Troodontids are important in research into the origin of birds because they share many anatomical characters with early birds. Crucially, the substantially complete fossil identified as WDC DML 001 ("Lori") is a troodontid from the Late Jurassic Morrison Formation, close to the time of Archaeopteryx and several troodontid specimens from the Tiaojishan Formation of China (Anchiornis) which are even older. The discovery of these Jurassic troodonts is positive physical evidence that derived deinonychosaurs were present before the time that avians arose, and basal paravians must have evolved much earlier. This fact strongly invalidates the "temporal paradox" cited by the few remaining opponents of the idea that birds are closely related to dinosaurs. Troodontid fossils were among the first dinosaur remains described. Initially, Leidy (1856) assumed they were lacertilian (lizards), but, by 1924, they were referred to Dinosauria by Gilmore, who suggested that they were ornithischians. It was not until 1945 that C.M. Sternberg recognized Troodontidae as a theropod family. Since 1969, Troodontidae has typically been allied with Dromaeosauridae, in a clade (natural group) known as Deinonychosauria, but this was by no means a consensus. Holtz (in 1994) erected the clade Bullatosauria, uniting Ornithomimosauria (the "ostrich-dinosaurs") and Troodontidae, on the basis of characteristics including, among others, an inflated braincase (parabasisphenoid) and a long, low opening in the upper jaw (the maxillary fenestra). Features of the pelvis also suggested they were less advanced than dromaeosaurids. New discoveries of primitive troodontids from China (such as Sinovenator and Mei), however, display strong similarities between Troodontidae, Dromaeosauridae and the primitive bird Archaeopteryx, and most paleontologists, including Holtz, now consider troodontids to be much more closely related to birds than they are to ornithomimosaurs, causing the clade Bullatosauria to be abandoned. One study of theropod systematics by members of the Theropod Working Group has uncovered striking similarities among the most basal dromaeosaurids, troodontids, and Archaeopteryx. This clade is together called Paraves by Novas and Pol. The cladogram published in Hwang et al. found that Archaeopteryx represents a more basal branch of Paraves, and places dromaeosaurids and troodontids as more derived. This raises the possibility that aerodynamic behaviors could be ancestral to all of Deinonychosauria. The extensive cladistic analysis conducted by Turner et al. (2012) supported the monophyly of Troodontidae. There are multiple possibilities of the genera included in Troodontidae as well as how they are related. Very primitive species, such as Anchiornis huxleyi, have alternately been found to be early troodontids or early members of the closely related group avilalae by various studies. The cladogram below follows the results of a study by Pascal Godefroit and colleagues in 2013. In 2014, Brusatte, Lloyd, Wang and Norell published an analysis on Coelurosauria, a simplified version shown below. This analysis included more troodontid species but failed to resolve many of their interrelationships, resulting in large "polytomies" (sets of species where the branching order in the family tree is uncertain). - Junchang Lü, Li Xu, Yongqing Liu, Xingliao Zhang, Songhai Jia, and Qiang Ji (2010). "A new troodontid (Theropoda: Troodontidae) from the Late Cretaceous of central China, and the radiation of Asian troodontids." (PDF). Acta Palaeontologica Polonica 55 (3): 381–388. doi:10.4202/app.2009.0047. - Currie, P. J. (1985). "Cranial anatomy of Stenonychosaurus inequalis (Saurischia, Theropoda) and its bearing on the origin of birds". Canadian Journal of Earth Sciences 22: 1643–1658. doi:10.1139/e85-173. - Castanhinha, R.; Mateus, O. (2006). "On the left-right asymmetry in dinosaurs". Journal of Vertebrate Paleontology 26 (Supp. 3): 48A. doi:10.1080/02724634.2006.10010069. - Mackovicky, Peter J.; Norell, Mark A. (2004). "Troodontidae". In Weishampel, David B.; Dodson, Peter; and Osmólska, Halszka (eds.). The Dinosauria (2nd ed.). Berkeley: University of California Press. pp. 184–195. ISBN 0-520-24209-2. - Holtz, T.R. Jr.; Brinkman, D.L.; Chandler, C.L. (1998). "Denticle morphometrics and a possibly omnivorous feeding habit for the theropod dinosaur Troodon" (PDF). Gaia 15: 159–166. - Currie, PJ; Dong, Z (2001). "New information on Cretaceous troodontids (Dinosauria, Theropoda) from the People's Republic of China". Canadian Journal of Earth Sciences 38: 1753–1766. doi:10.1139/e01-065. - Fowler, D.W.; Freedman, E.A.; Scannella, J.B.; Kambic, R.E. (2011). "The Predatory Ecology of Deinonychus and the Origin of Flapping in Birds". PLoS ONE 6 (12): e28964. doi:10.1371/journal.pone.0028964. PMC 3237572. PMID 22194962. - Xu; Norell (2004). "A new troodontid dinosaur from China with avian-like sleeping posture". Nature 431 (7010): 838–841. doi:10.1038/nature02898. PMID 15483610. - Bever, G.S. and Norell, M.A. (2009). "The perinate skull of Byronosaurus (Troodontidae) with observations on the cranial ontogeny of paravian theropods." American Museum Novitates, 3657: 51 pp. - Norell, Mark A.; Clark, James M.; Dashzeveg, Demberelyin; Barsbold, Rhinchen; Chiappe, Luis M.; Davidson, Amy R.; McKenna, Malcolm C.; Perle, Altangerel; Novacek, Michael J. (November 4, 1994). "A theropod dinosaur embryo and the affinities of the Flaming Cliffs dinosaur eggs". Science 266 (5186): 779–782. doi:10.1126/science.266.5186.779. PMID 17730398. - Hu, D., Hou, L., Zhang, L. and Xu, X. (2009) "A pre-Archaeopteryx troodontid theropod from China with long feathers on the metatarsus." Nature, 461, 1 October 2009: 640-643. doi:10.1038/nature08322 PMID 19794491. - Novas, F. E.; Pol, D. (2005). "New evidence on deinonychosaurian dinosaurs from the Late Cretaceous of Patagonia". Nature 3285 (7028): 858–861. doi:10.1038/nature03285. PMID 15729340. - Hwang, S.H.; Norell, M.A.; Ji, Q.; Gao, K.-Q. (2002). "New specimens of Microraptor zhaoianus (Theropoda: Dromaeosauridae) from Northeastern China". American Museum Novitates 3381: 1–44. doi:10.1206/0003-0082(2002)381<0001:nsomzt>2.0.co;2. - Turner, A. H.; Makovicky, P. J.; Norell, M. A. (2012). "A Review of Dromaeosaurid Systematics and Paravian Phylogeny". Bulletin of the American Museum of Natural History 371: 1. doi:10.1206/748.1. - Pascal Godefroit, Andrea Cau, Hu Dong-Yu, François Escuillié, Wu Wenhao and Gareth Dyke (2013). "A Jurassic avialan dinosaur from China resolves the early phylogenetic history of birds". Nature. in press (7454): 359–62. doi:10.1038/nature12168. PMID 23719374. - Brusatte, S. L.; Lloyd, G. T.; Wang, S. C.; Norell, M. A. (2014). "Gradual Assembly of Avian Body Plan Culminated in Rapid Rates of Evolution across the Dinosaur-Bird Transition". Current Biology 24 (20): 2386. doi:10.1016/j.cub.2014.08.034. |Wikimedia Commons has media related to Troodontidae.|

https://en.wikipedia.org/wiki/Troodontid

A bioterrorism attack is the deliberate release of viruses, bacteria, or other germs (agents) used to cause illness or death in people, animals, or plants. These agents are typically found in nature, but it is possible that they could be changed to increase their ability to cause disease, make them resistant to current medicines, or to increase their ability to be spread into the environment. Biological agents can be spread through the air, through water, or in food. Terrorists may use biological agents because they can be extremely difficult to detect and do not cause illness for several hours to several days. Some bioterrorism agents, like the smallpox virus, can be spread from person to person and some, like anthrax, can not. For information on which bioterrorism agents can be spread from person to person, please see the alphabetical list of bioterrorism agents. Bioterrorism agents can be separated into three categories, depending on how easily they can be spread and the severity of illness or death they cause. Category A agents are considered the highest risk and Category C agents are those that are considered emerging threats for disease. These high-priority agents include organisms or toxins that pose the highest risk to the public and national security because: - They can be easily spread or transmitted from person to person - They result in high death rates and have the potential for major public health impact - They might cause public panic and social disruption - They require special action for public health preparedness. These agents are the second highest priority because: - They are moderately easy to spread - They result in moderate illness rates and low death rates - They require specific enhancements of CDC's laboratory capacity and enhanced disease monitoring. These third highest priority agents include emerging pathogens that could be engineered for mass spread in the future because: - They are easily available - They are easily produced and spread - They have potential for high morbidity and mortality rates and major health impact. You can look for the bioterrorism agent by name on the A-Z List of Bioterrorism Agents/Diseases. The CDC and the American Red Cross have teamed up to answer questions and provide advice on steps you can take to prepare yourself and your loved ones in the event of a bioterrorist attack. For preparedness information and guidelines, please see Emergency Preparedness and You . The Department of Homeland Security has established a website to provide information to the public about emergencies and emergency preparedness. For information on what to do in the event of a bioterrorist attack, please see Ready.gov. - Department of Homeland Security – National Response Plan - American Red Cross – Terrorism Preparedness - The American Medical Association – Bioterrorism: Frequently Asked Questions - The Food and Drug Administration – Drug Preparedness and Response to Bioterrorism (information on antibiotics and dosage) - Environmental Protection Agency – Water Security - National Library of Medicine/National Institutes of Health Medline Plus – Biodefense and Bioterrorism - Page last updated February 12, 2007 - Content source: CDC Emergency Risk Communication Branch (ERCB), Division of Emergency Operations (DEO), Office of Public Health Preparedness and Response (OPHPR) - Content source: CDC Emergency Risk Communication Branch (ERCB), Division of Emergency Operations (DEO), National Center for Emerging and Zoonotic Infectious Diseases (NCEZID) Get email updates To receive email updates about this page, enter your email address:

http://www.emergency.cdc.gov/bioterrorism/overview.asp

Sinus conditions and sinus anatomy The sinuses do an important job, making mucus to keep the inside of the nose moist and helping to protect against dust, allergens and pollutants. However, when things go wrong with the sinuses, such as blockages or infections, the result can be very painful. The sinuses are a connected system of hollow cavities in the skull. The largest sinus cavities are around an inch across; others are much smaller. The sinus cavities include: - The maxillary sinuses (the largest), in the cheekbones. - The frontal sinuses, in the low-centre of the forehead. - The ethmoid sinuses, between the eyes, at the nasal bridge. - The sphenoid sinuses, in bones behind the eyes and nasal cavity. The sinuses are lined with soft, pink tissue called mucosa. Normally, the sinuses are empty except for a thin layer of mucus. The inside of the nose has ridges called turbinates. Normally these structures help humidify and filter air. The nose is divided in the centre by a thin wall, called the septum. Most of the sinuses drain into the nose through a small channel or drainage pathway called the middle meatus. The purpose of the sinuses is unclear. One theory is that sinuses help humidify the air we breathe in; another is that they enhance our voices. Acute sinusitis(sinus infection): Viruses or bacteria infect the sinus cavity, causing inflammation. Increased mucus production, nasal congestion, discomfort in the cheeks, forehead or around the eyes and headaches are common symptoms. It often develops quickly and can go on for up to 12 weeks. Chronic sinusitis(or chronic rhinosinusitis): More than just a series of infections, chronic sinusitis is a persistent process of inflammation of the sinuses that goes on for more than 12 weeks and can continue for several months. Allergic rhinitis: Allergens like pollen, dust mites or pet dander cause the defences in the nose and sinuses to overreact. Mucus, nasal stuffiness, sneezing and itching result. Deviated septum: If the septum that divides the nose deviates too far to one side, airflow can be obstructed. Turbinate hypertophy: The ridges on the nasal septum are enlarged, potentially obstructing airflow. Nasal polyps: Small growths called polyps sometimes grow in the nasal cavity, in response to inflammation. Asthma, chronic sinus infections and allergic rhinitis can lead to nasal polyps.

http://www.webmd.boots.com/allergies/guide/picture-of-the-sinuses

An important prerequisite for intelligence is a good short-term memory which can store and process the information needed for ongoing processes. This 'working memory' is a kind of mental notepad – without it, we could not follow a conversation, do mental arithmetic or play any simple game. A new study has discovered neurons allowing crows to remember short-term. In the animal kingdom, the group of birds including crows and ravens, the corvids, are known for their intelligence because they have just such a working memory, but their endbrain – which is highly-developed but has a fundamentally different structure from that of mammals – has no cerebral cortex; and that is the part of the brain which in mammals produces the working memory. So how do corvids manage to have a working memory with no cerebral cortex? To answer that question, three researchers from the Institute for Neurobiology at Tübingen University taught crows to play a version of the children’s game of “pairs.” Using a computer monitor, Lena Veit, Konstantin Hartmann and Professor Andreas Nieder briefly showed the crows a random image. The crows had to remember it for one second before choosing the same image from a selection of four by tapping the remembered picture with their beaks. In order to choose the correct image, they must have stored it in a working memory – which they appeared to do quite easily. Crows were able to select the picture they had just seen from a number of images. Tübingen researchers discovered which cells were activated to briefly store and process information. Image: LS Tierphysiologie Simultaneous measurements of electric potentials in the crows’ brains showed that nerve cells in one particular area of the endbrain were responsible for this capacity to remember. Although the image had disappeared from the screen, those cells remained active during the short period of remembering – retaining the information about the image until the crow retrieved it in order to make the right choice. If a crow couldn’t remember and selected a wrong image, those particular endbrain cells were barely activated. Prolonged activation of such cells ensured that important information could be stored and later accessed. Nieder and his team concluded that cognitive abilities are possible in a range of differently-structured brains. “Clearly, a good working memory – an important characteristic of human beings – can also exist without a layered cerebral cortex. The corvids’ fundamentally differently-structured endbrain shows that evolution has found a number of independent solutions,” says Lena Veit. There are great benefits in the ability to temporarily store information. “An organism with a good working memory is intelligent; it is released from the compulsion to respond immediately to stimuli,” says Professor Nieder. “The big question is now – how do neural networks in the brain have to be composed in order to actively store and process information?” Citation: Lena Veit, Konstantin Hartmann, Andreas Nieder, Neuronal Correlates of Visual Working Memory in the Corvid Endbrain, Journal of Neuroscience, June 4, 2014, DOI:10.1523/JNEUROSCI.0612-14.2014 - PHYSICAL SCIENCES - EARTH SCIENCES - LIFE SCIENCES - SOCIAL SCIENCES Subscribe to the newsletter Stay in touch with the scientific world! Know Science And Want To Write? - Pseudoscience Creeping Into Your Conference? A Case In GMOs And Glyphosate - My Thoughts On The LIGO-VIRGO Result - Beyond Diamonds And Gems: The World's Rarest Minerals - Henri Poincaré Predicted The Existence Of Gravitational Waves As Early As June 5, 1905 - Gravitational Waves? Watch the LIGO press conference at 10:30 Eastern. - The 10 catastrophic mistakes virtual reality firms are making - When Machines Can Do Any Job, What Will Humans Do? - "Alas, all I see are a couple of links to articles in a blog by an ex tv weatherman. If their theories..." - "Where did the estimate a peak power equal to the total electrical power consumption of human mankind..." - "Already linked to one here, 2/11/16 at 21:28. Can't you follow a simple thread? ..." - "That’s history for you! Interestingly, I was scouring the web for the great physicist /..." - "That reminds me of the chemiosmotic war that followed Peter Mitchell's publication of his theory..." - Beard Microbiology: Grubby Hipsters May Be On To Something - Water Tops the List of Health Concerns for Competitive Eaters - Natural Flavors Are More Radioactive Than Artificial Ones. - Bariatric Surgery Beneficial Even for Older People - Opiates No Better at Easing Knee Osteoarthritis Pain - BRCA Tests Increasing for Younger Breast Cancer Patients - Dark matter scientists on brink of discovering elusive particles - Science responds to globalized disease threat to farms and food systems - Americans' support for science remains strong - Loss of sleep during adolescence may be a diabetes danger - Study of cognitive development in deaf children revisits longstanding debate

http://www.science20.com/news_articles/why_crows_are_intelligent_working_memory_neurons-138073

It is widely accepted that supermassive black holes (SMBHs) sit in the centre of elliptical galaxies or bulges of spiral galaxies. They suck in as much matter as possible, generating blasts of radiation. Stars, gas and everything else nearby forms a compact “halo” and then falls to a gravitationally enforced death spiral. The greedy nature and the sheer size of these black holes have led to the idea that dark matter may supply (or may have supplied) the SMBH with some mass during its evolution. But could it be that dark matter may not be significantly involved after all? This might be one cosmic phenomenon dark matter can’t be blamed for… Black hole accretion disks are compact halos created as dust, gas and other debris are pulled toward a black hole event horizon. Accretion disks radiate electromagnetic radiation, and the frequency of which depends on the mass of the black hole. The more massive it is, the higher the energy of radiation emitted into space. In the case of a SMBH, the huge mass causes very bright emission as the matter from the accretion disk falls into the event horizon (the point at which gravity becomes so strong that even light cannot escape). As accretion disk matter falls toward the event horizon, approximately 10% of the mass is converted into energy and ejected as X-rays. This is a far more efficient energy conversion rate than the most efficient nuclear fusion reaction (approximately 0.5%). This X-ray emission can then be observed, creating a quasar, signifying a SMBH is driving the active galaxy. Interestingly, an SMBH is not thought to be formed from single dead massive star. They are thought to have been created from a “seed” and then grown over billions of years. The source of the mass feeding the growing SMBH comes from its accretion disk, but it is uncertain what form the matter comes in and at what rate it “feeds” the black hole. There are several possibilities as to how the largest black holes were seeded, but two are the most widely accepted: The mechanisms affecting the rate of accretion disk growth are not so clear-cut. Some theories suggest that huge quantities (most of the black hole mass) comes from dark matter. However, as dark matter is “non-baryonic” (i.e. the opposite to baryonic matter – the matter we know, love and observe in our universe) it will emit very little radiation as it falls into the black hole event horizon. If this is the case, SMBHs would grow disproportionately when compared with radiation emitted from galactic centres (only baryonic particles will emit X-rays). New research headed by Sebastien Peirani (at the Institut d’Astrophysique de Paris, France) suggests only a very small fraction of a SMBH is composed from dark matter as it evolved. Dark matter is predicted to be collisionless and will be scattered very easily by baryonic gas clouds and stars. It seems unlikely that dark matter will be able to stay inside the black hole’s accretion disk for very long before it is repelled by all the “normal” matter being pulled toward the event horizon. By modelling a “typical” accretion disk and comparing the results with observations of quasar luminosity, the French group found that most of the matter fuelling the SMBHs is relativistic baryonic matter. At a critical distance, outside the black hole, baryonic matter from the accretion disk is accelerated to a significant fraction of the speed of light, emitting radiation. Comparing this with simulations of a collisionless disk (i.e. the characteristics of dark matter), the baryonic model fits observations the best. “Application of our results to black hole seeds hosted by halos issued from cosmological simulations indicate that dark matter contributes to no more than [approx.] 10% of the total accreted mass, confirming that the bolometric quasar luminosity is related to the baryonic accretion history of the black hole.” – Abstract from “Dark Matter Accretion into Supermassive Black Holes“

http://www.universetoday.com/13091/

1.Present the key vocab. to students 2. after students show good comprehension of the proverb, have them relate to real life situations and make sentences with it. 3. Have students act out the story of the proverb Date Last Modified: September 11, 2010 provide students text after listening activity; ask students to highlight characters that are new to them; collect new expressions from the text students highlighted and assign each student two or three new words to use dictionary to find out the relevant meanings

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In geology and earth science, a plateau (or ; plural plateaus or plateaux), also called a high plain or tableland, is an area of highland, usually consisting of relatively flat terrain that is raised significantly above the surrounding area, often with one or more sides with steep slopes. Plateaus can be formed by a number of processes, including upwelling of volcanic magma, extrusion of lava, and erosion by water and glaciers. Magma rises from the mantle causing the ground to swell upward, in really large, flat areas of rock that are uplifted. Plateaus can also be built up by lava spreading outward from cracks and weak areas in the crust. Plateaus can also be formed by the erosional processes of glaciers on mountain ranges, leaving them sitting between the mountain ranges. Water can also erode mountains and other landforms down into plateaus. Computer modeling studies suggest that high plateaus may also be partially a result from the feedback between tectonic deformation and dry climatic conditions created at the lee side of growing orogens. Plateaus are classified according to their surrounding environment. The largest and highest plateau in the world is the Tibetan Plateau, called the "roof of the world", which is still being formed by the collisions of the Indo-Australian and Eurasian tectonic plates. The Tibetan plateau covers approximately 2500000km2, at about 5000m (16,000feet) above sea level. The plateau is sufficiently high enough to reverse the Hadley cell convection cycles and to drive the monsoons of India towards the south. The second-highest plateau is the Deosai Plateau of the Deosai National Park (also known as Deoasai Plains) at an average elevation of 4114m (13,497feet). It is located in the Astore and Skardu districts of Gilgit-Baltistan, in northern Pakistan. Deosai means 'the land of giants'. The park protects an area of 3000km2. It is known for its rich flora and fauna of the Karakoram-West Tibetan Plateau alpine steppe ecoregion. In spring it is covered by sweeps of wildflowers and a wide variety of butterflies. The highest point in Deosai is Deosai Lake, or Sheosar Lake from the Shina language meaning "Blind lake" (Sheo - Blind, Sar - lake) near the Chilim Valley. The lake lies at an elevation of 4142m (13,589feet), one of the highest lakes in the world, and is 2.3km long, 1.8km wide, and 40m (130feet) deep on average. Some other major plateaus in Asia are: Armenian Highlands (~400000km2, elevation 900-2100m), Iranian plateau(~3700000km2, elevation 300-1500m), Anatolian Plateau, Mongolian Plateau (~2600000km2, elevation 1000-1500m), and the Deccan Plateau. Another very large plateau is the Antarctic Plateau, which covers most of central Antarctica, where there are no known mountains, but rather 3000m (10,000feet) or more of ice - which very slowly spreads toward the coastline via enormous glaciers. This ice cap is so massive that echolocation sound measurements of the thickness of the ice have shown that large parts of the "dry land" surface of Antarctica have been pressed below sea level. Thus, if the icecap were suddenly removed, large areas of Antarctica would be flooded by the oceans. On the other hand, were the icecap to gradually melt away, the surface of the land beneath it would gradually rebound/[Isostasy] away from the center of the Earth, and that land would ultimately rise above sea level. In northern Arizona and southern Utah the Colorado Plateau is bisected by the Colorado River and the Grand Canyon. How this came to be is that over 10 million years ago, a river was already there, though not necessarily on exactly the same course. Then, subterranean geological forces caused the land in that part of North America to gradually rise by about a centimeter per year for millions of years. An unusual balance occurred: the river that would become the Colorado River was able to erode into the crust of the Earth at a nearly equal rate to the uplift of the plateau. Now, millions of years later, the North Rim of the Grand Canyon is at an elevation of about 2450m (8,040feet) above sea level, and the South Rim of the Grand Canyon is about 2150m (7,050feet) above sea level. At its deepest, the Colorado River is about 1830m (6,000feet) below the level of the North Rim High altitude plateau in North America is the Mexican plateau. With an area of 601882km2 and average height of 1,825 m, it is the home of more than 70 millions people. See main article: Tepui. A tepui, or tepuy, is a table-top mountain or mesa found in the Guiana Highlands of South America, especially in Venezuela and western Guyana. The word tepui means "house of the gods" in the native tongue of the Pemon, the indigenous people who inhabit the Gran Sabana. Tepuis tend to be found as isolated entities rather than in connected ranges, which makes them the host of a unique array of endemic plant and animal species. Some of the most outstanding tepuis are Neblina, Autana, Auyantepui and Mount Roraima. They are typically composed of sheer blocks of Precambrian quartz arenite sandstone that rise abruptly from the jungle, giving rise to spectacular natural scenery. Auyantepui is the source of Angel Falls, the world's tallest waterfall.The parallel Sierra of Andes delimit one of the world highest plateaus: the Altiplano, (Spanish for "high plain"), Andean Plateau or Bolivian Plateau. It lies in west-central South America, where the Andes are at their widest, is the most extensive area of high plateau on Earth outside of Tibet. The bulk of the Altiplano lies within Bolivian and Peruvian territory while its southern parts lie in Chile and Argentina. The Altiplano plateau hosts several cities like Puno, Oruro, Potosí, Cuzco and La Paz, the administrative seat of Bolivia. Northeastern Altiplano is more humid than the Southwestern, the latter of which hosts several salares, or salt flats, due to its aridity. At the Bolivia-Peru border lies Lake Titicaca, the largest lake in South America. The highest African plateau is the Ethiopian Highlands which cover the central part of Ethiopia. It forms the largest continuous area of its altitude in the continent, with little of its surface falling below 1500 m (4,921 ft), while the summits reach heights of up to 4550 m (14,928 ft). It is sometimes called the Roof of Africa due to its height and large area. Another example is the Highveld which is the portion of the South African inland plateau which has an altitude above approximately 1500 m, but below 2100 m, thus excluding the Lesotho mountain regions. It is home to some of largest South-African urban agglomerations. The Western Plateau, part of the Australian Shield, is an ancient craton covering much of the continent's southwest, an area of some 700,000 square kilometres. It has an average elevation of between 305 and 460 m. The North Island Volcanic Plateau is an area of high land occupying much of the centre of the North Island of New Zealand, with volcanoes, lava plateaus, and crater lakes, the most notable of which is the country's largest lake, Lake Taupo. The plateau stretches approximately 100 km east to west and 130 km north to south. The majority of the plateau is more than 600 m above sea level.

http://everything.explained.today/Plateau/

From Latin: sector "a cutter" Definition: The number of square units it takes to exactly fill a sector of a circle. Try this Drag one of the orange dots that define the endpoints of the sector. The sector area is recalculated as you drag. What the formulae are doing is taking the area of the whole circle, and then taking a fraction of that depending on what fraction of the circle the sector fills. So for example, if the central angle was 90°, then the sector would have an area equal to one quarter of the whole circle. If you know the central angle This is the method used in the animation above. C is the central angle in r is the radius of the circle of which the sector is part. π is Pi, approximately 3.142 If you know the arc length L is the arc length. R is the radius of the circle of which the sector is part. Sector area is proportional to arc length The area enclosed by a sector is proportional to the arc length of the sector. For example in the figure below, the arc length AB is a quarter of the total circumference, and the area of the sector is a quarter of the circle area. Similarly below, the arc length is half the circumference, and the area id half the total circle. You can experiment with other proportions in the applet at the top of the page. Other circle topics Equations of a circle Angles in a circle (C) 2009 Copyright Math Open Reference. All rights reserved

http://www.mathopenref.com/arcsectorarea.html

Mitochondria play an essential role in the generation of energy in eukaryotic cells. Mitochondria are the organelles that are the main “chemical factories” of the cell where cellular aerobic respiration—using the Krebs (citric acid) cycle and respiratory electron transport to produce NADH (nicotinamide adenine dinucleotide) and ATP (adenosine triphosphate)—occurs. In the light microscope, mitochondria look like short rods or thin filaments about 0.5 to 2 microns long. A mitochondrion is made up of a smooth outer membrane and an inner membrane that is folded into tubular shapes called cristae. Many aerobic respiration reactions are catalyzed by enzymes that are bound to mitochondrial membranes. Other reactions occur in the space between the inner and outer mitochondrial membranes. Cells may contain several hundred mitochondria. Cells that are dividing and cells that are metabolically active need larger amounts of ATP and usually have large numbers of mitochondria. Size and Structure All eukaryotic cells except some primitive protozoans contain mitochondria. All mitochondria contain their own DNA(genomes). There are typically between twenty and one hundred copies of the mitochondrial genome per mitochondrion. The mitochondria of multicellular animals contain genomes of 14 to 20 kilobases (kb), present as single circles. The mitochondrial DNA of some organisms, such as some protozoa, algae, and fungi, is organized in linear molecules with ends of chromosomes (telomeres) much like nuclear chromosomes. In contrast, the mitochondrial DNA of higher plants is larger and more complex—from 200 to 2,500 kb—and is present in many different molecules. The size and organization of the mitochondrial genome vary widely from one plant species to another. Electron micrographs of mitochondrial DNA show linear and circular DNAs of a variety of sizes and complex, branched molecules that are larger than the size of the genome. Cloning the mitochondrial DNA and comparing the sequences of the clones show that the entire complexity of a plant mitochondrial genome can be represented as a “master circle.” Also, it has been learned that sequences are repeated on the master chromosome. The repeated sequences differ for different plant species. A series of recombination events between these identical repeated sequences results in a series of rearrangements of mitochondrial DNA and forms the complex, multiple molecules of varying sizes that are the physical structure of the plant mitochondrial genome. Adding to the complexity of mitochondrial DNAs in higher plants is the fact that some plants, such as corn, contain extra chromosomal mitochondrial nucleic acids. Plasmid-like DNAs (circular double-stranded molecules) and double-stranded and single-stranded RNAs have been found in some corn strains. Genes Encoded by Mitochondrial DNA |Genes Encoded by Mitochondrial DNA| The ribosomes of mitochondria are different from those of chloroplasts and the cytoplasm, using a slightly different genetic code (a sequence of three bases that codes for a particular amino acid). Mitochondrial genomes code for all of the ribosomal RNAs found in mitochondria and for most of the tRNAs. Mitochondria make only a small number of proteins that are needed for electron transport and ATP production. The other proteins needed in mitochondria are coded by nuclear DNA, translated in the cytoplasm of the cell, then transported into the mitochondria. Plant mitochondria do not encode a full set of tRNAs, and some are imported from the cytoplasm. Even though the mitochondrial genome of higher plants is much larger than that of animals, the plant mitochondrial genome codes for only a few more genes. The mitochondrial genome of Arabidopsis has been sequenced and contains thirty-two protein-coding genes, twenty-two tRNA genes, and three ribosomal RNA genes. Exchange of DNA Mitochondrial DNA from plants also differs from that of animals in that mitochondrial DNAs contain segments of DNA that originally were in nuclear and chloroplast DNAs. There appear to have been exchanges of DNAs between all three of the higher plant genomes. There is evidence that mitochondrial genes have been transferred to the nucleus and some mitochondrial tRNAs appear to be of chloroplast origin. Changes in nuclear genes have been shown to lead to changes in the copy number of the different mitochondrial DNA configurations. Mitochondria and chloroplasts contain the biochemical machinery to alter the sequence of the final messenger RNA (mRNA) product in a process called RNA editing. The most common editing is changing a cytosine to a uracil (two of the bases found on the “rungs” of DNA molecules and which are responsible for determining the nucleotide sequences that form the genetic code). Inheritance of Mitochondrial DNA Given the complex branched network of plant mitochondrial DNA, it is difficult to see how the inheritance of a complete genome is ensured. It is still not clearly understood how this complex network of DNAs is passed to daughter cells in away that assures that all of the genetic information is maintained.

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2 Answers | Add Yours One of the most significant impacts of the Trail of Tears was that it marked a point where Native Americans lost any semblance of power. The removal of different tribes to different parts of the United States disbanded them. It dislodged them from their homes and relegated them to a position of powerlessness. Native Americans would no longer be a formidable voice in negotiating with the American government. The encroachment that started with the Trail of Tears continued throughout the 19th and 20th Centuries. It is difficult to fathom that at one point in time, the land that is now firmly accepted as "American" was inhabited by another. The Native Americans would argue that no one "owns" the land. Yet, one of the most significant impacts of the Trail of Tears was that it showed that the Native Americans would never come close to exercising power over their destinies on this land. Another significant impact of the Trail of Tears is that it forged Andrew Jackson's legacy as one of the most cruel Presidents. Even the most ardent supporter of Jackson has to admit that the Trail of Tears brings out a rather dark side to the leader of "Jacksonian Democracy." Jackson earned the nickname that Native Americans gave him, "Sharp Knife," through the Trail of Tears. The idea that American government would enforce and support a policy in which the “trail where they cried" becomes a part of logistical practice is morally repugnant. The zeal with which "Old Hickory" proceeded with Native American relegation has to be seen as a stain on American History. This becomes another significant impact of the Trail of Tears. The Trail of Tears is the name given to the ethnic cleansing and the relocation of Native Americans because of the Indian Removal Act of 1830. It removed members of Indian tribes from their lands into "Indian territory" in present day Oklahoma. It led to bigger conflicts between the natives and Americans because of maltreatment and search of Gold. It led to violent confrontations and treaties that never worked for either side. Missionaries taught Indians to live with whites and to Christianize them. We’ve answered 301,149 questions. We can answer yours, too.Ask a question

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Types of volcanic eruptions Several types of volcanic eruptions—during which lava, tephra (ash, lapilli, volcanic bombs and blocks), and assorted gases are expelled from a volcanic vent or fissure—have been distinguished by volcanologists. These are often named after famous volcanoes where that type of behavior has been observed. Some volcanoes may exhibit only one characteristic type of eruption during a period of activity, while others may display an entire sequence of types all in one eruptive series. There are three different types of eruptions. The most well-observed are magmatic eruptions, which involve the decompression of gas within magma that propels it forward. Phreatomagmatic eruptions are another type of volcanic eruption, driven by the compression of gas within magma, the direct opposite of the process powering magmatic activity. The third eruptive type is the phreatic eruption, which is driven by the superheating of steam via contact with magma; these eruptive types often exhibit no magmatic release, instead causing the granulation of existing rock. Within these wide-defining eruptive types are several subtypes. The weakest are Hawaiian and submarine, then Strombolian, followed by Vulcanian and Surtseyan. The stronger eruptive types are Pelean eruptions, followed by Plinian eruptions; the strongest eruptions are called "Ultra Plinian." Subglacial and phreatic eruptions are defined by their eruptive mechanism, and vary in strength. An important measure of eruptive strength is Volcanic Explosivity Index (VEI), an order of magnitude scale ranging from 0 to 8 that often correlates to eruptive types. Volcanic eruptions arise through three main mechanisms: - Gas release under decompression causing magmatic eruptions - Thermal contraction from chilling on contact with water causing phreatomagmatic eruptions - Ejection of entrained particles during steam eruptions causing phreatic eruptions There are two types of eruptions in terms of activity, explosive eruptions and effusive eruptions. Explosive eruptions are characterized by gas-driven explosions that propels magma and tephra. Effusive eruptions, meanwhile, are characterized by the outpouring of lava without significant explosive eruption. Volcanic eruptions vary widely in strength. On the one extreme there are effusive Hawaiian eruptions, which are characterized by lava fountains and fluid lava flows, which are typically not very dangerous. On the other extreme, Plinian eruptions are large, violent, and highly dangerous explosive events. Volcanoes are not bound to one eruptive style, and frequently display many different types, both passive and explosive, even the span of a single eruptive cycle. Volcanoes do not always erupt vertically from a single crater near their peak, either. Some volcanoes exhibit lateral and fissure eruptions. Notably, many Hawaiian eruptions start from rift zones, and some of the strongest Surtseyan eruptions develop along fracture zones. Scientists believed that pulses of magma mixed together in the chamber before climbing upward—a process estimated to take several thousands of years. But Columbia University volcanologists found that the eruption of Costa Rica’s Irazú Volcano in 1963 was likely triggered by magma that took a nonstop route from the mantle over just a few months. Volcano explosivity index The volcanic explosivity index (commonly shortened VEI) is a scale, from 0 to 8, for measuring the strength of eruptions. It is used by the Smithsonian Institution's Global Volcanism Program in assessing the impact of historic and prehistoric lava flows. It operates in a way similar to the Richter scale for earthquakes, in that each interval in value represents a tenfold increasing in magnitude (it is logarithmic). The vast majority of volcanic eruptions are of VEIs between 0 and 2. Volcanic eruptions by VEI index |VEI||Plume height||Eruptive volume *||Eruption type||Frequency **||Example| |0||<100 m (330 ft)||1,000 m3 (35,300 cu ft)||Hawaiian||Continuous||Kilauea| |1||100–1,000 m (300–3,300 ft)||10,000 m3 (353,000 cu ft)||Hawaiian/Strombolian||Months||Stromboli| |2||1–5 km (1–3 mi)||1,000,000 m3 (35,300,000 cu ft) †||Strombolian/Vulcanian||Months||Galeras (1992)| |3||3–15 km (2–9 mi)||10,000,000 m3 (353,000,000 cu ft)||Vulcanian||Yearly||Nevado del Ruiz (1985)| |4||10–25 km (6–16 mi)||100,000,000 m3 (0.024 cu mi)||Vulcanian/Peléan||Few years||Eyjafjallajökull (2010)| |5||>25 km (16 mi)||1 km3 (0.24 cu mi)||Plinian||5–10 years||Mount St. Helens (1980)| |6||>25 km (16 mi)||10 km3 (2 cu mi)||Plinian/Ultra Plinian||1,000 years||Krakatoa (1883)| |7||>25 km (16 mi)||100 km3 (20 cu mi)||Ultra Plinian||10,000 years||Tambora (1815)| |8||>25 km (16 mi)||1,000 km3 (200 cu mi)||Supervolcanic||100,000 years||Lake Toba (74 ka)| | * This is the minimum eruptive volume necessary for the eruption to be considered within the category. ** Values are a rough estimate. Exceptions occur. † There is a discontinuity between the 2nd and 3rd VEI level; instead of increasing by a magnitude of 10, the value increases by a magnitude of 100 (from 10,000 to 1,000,000). Magmatic eruptions produce juvenile clasts during explosive decompression from gas release. They range in intensity from the relatively small lava fountains on Hawaii to catastrophic Ultra Plinian eruption columns more than 30 km (19 mi) high, bigger than the eruption of Mount Vesuvius in 79 that buried Pompeii. Hawaiian eruptions are a type of volcanic eruption, named after the Hawaiian volcanoes with which this eruptive type is hallmark. Hawaiian eruptions are the calmest types of volcanic events, characterized by the effusive eruption of very fluid basalt-type lavas with low gaseous content. The volume of ejected material from Hawaiian eruptions is less than half of that found in other eruptive types. Steady production of small amounts of lava builds up the large, broad form of a shield volcano. Eruptions are not centralized at the main summit as with other volcanic types, and often occur at vents around the summit and from fissure vents radiating out of the center. Hawaiian eruptions often begin as a line of vent eruptions along a fissure vent, a so-called "curtain of fire." These die down as the lava begins to concentrate at a few of the vents. Central-vent eruptions, meanwhile, often take the form of large lava fountains (both continuous and sporadic), which can reach heights of hundreds of meters or more. The particles from lava fountains usually cool in the air before hitting the ground, resulting in the accumulation of cindery scoria fragments; however, when the air is especially thick with clasts, they cannot cool off fast enough due to the surrounding heat, and hit the ground still hot, the accumulation of which forms spatter cones. If eruptive rates are high enough, they may even form splatter-fed lava flows. Hawaiian eruptions are often extremely long lived; Puʻu ʻŌʻō, a cinder cone of Kilauea, has been erupting continuously since 1983. Another Hawaiian volcanic feature is the formation of active lava lakes, self-maintaining pools of raw lava with a thin crust of semi-cooled rock; there are currently only 5 such lakes in the world, and the one at Kīlauea's Kupaianaha vent is one of them. Flows from Hawaiian eruptions are basaltic, and can be divided into two types by their structural characteristics. Pahoehoe lava is a relatively smooth lava flow that can be billowy or ropey. They can move as one sheet, by the advancement of "toes," or as a snaking lava column. A'a lava flows are denser and more viscous then pahoehoe, and tend to move slower. Flows can measure 2 to 20 m (7 to 66 ft) thick. A'a flows are so thick that the outside layers cools into a rubble-like mass, insulating the still-hot interior and preventing it from cooling. A'a lava moves in a peculiar way—the front of the flow steepens due to pressure from behind until it breaks off, after which the general mass behind it moves forward. Pahoehoe lava can sometimes become A'a lava due to increasing viscosity or increasing rate of shear, but A'a lava never turns into pahoehoe flow. Hawaiian eruptions are responsible for several unique volcanological objects. Small volcanic particles are carried and formed by the wind, chilling quickly into teardrop-shaped glassy fragments known as Pele's tears (after Pele, the Hawaiian volcano deity). During especially high winds these chunks may even take the form of long drawn-out strands, known as Pele's hair. Sometimes basalt aerates into reticulite, the lowest density rock type on earth. Although Hawaiian eruptions are named after the volcanoes of Hawaii, they are not necessarily restricted to them; the largest lava fountain ever recorded formed on the island of Izu Ōshima (on Mount Mihara) in 1986, a 1,600 m (5,249 ft) gusher that was more than twice as high as the mountain itself (which stands at 764 m (2,507 ft)). Volcanoes known to have Hawaiian activity include: - Puʻu ʻŌʻō, a parasitic cinder cone located on Kilauea on the island of Hawaiʻi which has been erupting continuously since 1983. The eruptions began with a 6 km (4 mi)-long fissure-based "curtain of fire" on 3 January. These gave way to centralized eruptions on the site of Kilauea's east rift, eventually building up the still active cone. - For a list of all of the volcanoes of Hawaii, see List of volcanoes in the Hawaiian - Emperor seamount chain. - Mount Etna, Italy. - Mount Mihara in 1986 (see above paragraph) Strombolian eruptions are a type of volcanic eruption, named after the volcano Stromboli, which has been erupting continuously for centuries. Strombolian eruptions are driven by the bursting of gas bubbles within the magma. These gas bubbles within the magma accumulate and coalesce into large bubbles, called gas slugs. These grow large enough to rise through the lava column. Upon reaching the surface, the difference in air pressure causes the bubble to burst with a loud pop, throwing magma in the air in a way similar to a soap bubble. Because of the high gas pressures associated with the lavas, continued activity is generally in the form of episodic explosive eruptions accompanied by the distinctive loud blasts. During eruptions, these blasts occur as often as every few minutes. The term "Strombolian" has been used indiscriminately to describe a wide variety of volcanic eruptions, varying from small volcanic blasts to large eruptive columns. In reality, true Strombolian eruptions are characterized by short-lived and explosive eruptions of lavas with intermediate viscosity, often ejected high into the air. Columns can measure hundreds of meters in height. The lavas formed by Strombolian eruptions are a form of relatively viscous basaltic lava, and its end product is mostly scoria. The relative passivity of Strombolian eruptions, and its non-damaging nature to its source vent allow Strombolian eruptions to continue unabated for thousands of years, and also makes it one of the least dangerous eruptive types. Strombolian eruptions eject volcanic bombs and lapilli fragments that travel in parabolic paths before landing around their source vent. The steady accumulation of small fragments builds cinder cones composed completely of basaltic pyroclasts. This form of accumulation tends to result in well-ordered rings of tephra. Strombolian eruptions are similar to Hawaiian eruptions, but there are differences. Strombolian eruptions are noisier, produce no sustained eruptive columns, do not produce some volcanic products associated with Hawaiian volcanism (specifically Pele's tears and Pele's hair), and produce fewer molten lava flows (although the eruptive material does tend to form small rivulets). Volcanoes known to have Strombolian activity include: - Parícutin, Mexico, which erupted from a fissure in a cornfield in 1943. Two years into its life, pyroclastic activity began to wane, and the outpouring of lava from its base became its primary mode of activity. Eruptions ceased in 1952, and the final height was 424 m (1,391 ft). This was the first time that scientists are able to observe the complete life cycle of a volcano. - Mount Etna, Italy, which has displayed Strombolian activity in recent eruptions, for example in 1981, 1999, 2002-2003, and 2009. - Mount Erebus in Antarctica, the southernmost active volcano in the world, having been observed erupting since 1972. Eruptive activity at Erebus consists of frequent Strombolian activity. - Stromboli itself. The namesake of the mild explosive activity that it possesses has been active throughout historical time; essentially continuous Strombolian eruptions, occasionally accompanied by lava flows, have been recorded at Stromboli for more than a millennium. Vulcanian eruptions are a type of volcanic eruption, named after the volcano Vulcano, which means the word Volcano. It was named so following Giuseppe Mercalli's observations of its 1888-1890 eruptions. In Vulcanian eruptions, highly viscous magma within the volcano make it difficult for vesiculate gases to escape. Similar to Strombolian eruptions, this leads to the buildup of high gas pressure, eventually popping the cap holding the magma down and resulting in an explosive eruption. However, unlike Strombolian eruptions, ejected lava fragments are not aerodynamic; this is due to the higher viscosity of Vulcanian magma and the greater incorporation of crystalline material broken off from the former cap. They are also more explosive than their Strombolian counterparts, with eruptive columns often reaching between 5 and 10 km (3 and 6 mi) high. Lastly, Vulcanian deposits are andesitic to dacitic rather than basaltic. Initial Vulcanian activity is characterized by a series of short-lived explosions, lasting a few minutes to a few hours and typified by the ejection of volcanic bombs and blocks. These eruptions wear down the lava dome holding the magma down, and it disintegrates, leading to much more quiet and continuous eruptions. Thus an early sign of future Vulcanian activity is lava dome growth, and its collapse generates an outpouring of pyroclastic material down the volcano's slope. Deposits near the source vent consist of large volcanic blocks and bombs, with so-called "bread-crust bombs" being especially common. These deeply cracked volcanic chunks form when the exterior of ejected lava cools quickly into a glassy or fine-grained shell, but the inside continues to cool and vesiculate. The center of the fragment expands, cracking the exterior. However the bulk of Vulcanian deposits are fine grained ash. The ash is only moderately dispersed, and its abundance indicates a high degree of fragmentation, the result of high gas contents within the magma. In some cases these have been found to be the result of interaction with meteoric water, suggesting that Vulcanian eruptions are partially hydrovolcanic. Volcanoes that have exhibited Vulcanian activity include: - Sakurajima, Japan has been the site of Vulcanian activity near-continuously since 1955. - Tavurvur, Papua New Guinea, one of several volcanoes in the Rabaul Caldera. - Irazú Volcano in Costa Rica exhibited Vulcanian activity in its 1965 eruption. Peléan eruptions (or nuée ardente) are a type of volcanic eruption, named after the volcano Mount Pelée in Martinique, the site of a massive Peléan eruption in 1902 that is one of the worst natural disasters in history. In Peléan eruptions, a large amount of gas, dust, ash, and lava fragments are blown out the volcano's central crater, driven by the collapse of rhyolite, dacite, and andesite lava dome collapses that often create large eruptive columns. An early sign of a coming eruption is the growth of a so-called Peléan or lava spine, a bulge in the volcano's summit preempting its total collapse. The material collapses upon itself, forming a fast-moving pyroclastic flow (known as a block-and-ash flow) that moves down the side of the mountain at tremendous speeds, often over 150 km (93 mi) per hour. These massive landslides make Peléan eruptions one of the most dangerous in the world, capable of tearing through populated areas and causing massive loss of life. The 1902 eruption of Mount Pelée caused tremendous destruction, killing more than 30,000 people and competely destroying the town of St. Pierre, the worst volcanic event in the 20th century. Peléan eruptions are characterized most prominently by the incandescent pyroclastic flows that they drive. The mechanics of a Peléan eruption are very similar to that of a Vulcanian eruption, except that in Peléan eruptions the volcano's structure is able to withstand more pressure, hence the eruption occurs as one large explosion rather than several smaller ones. Volcanoes known to have Peléan activity include: - Mount Pelée, Martinique. The 1902 eruption of Mount Pelée completely devastated the island, destroying the town of St. Pierre and leaving only 3 survivors. The eruption was directly preceded by lava dome growth. - Mayon Volcano, the Philippines most active volcano. It has been the site of many different types of eruptions, Peléan included. Approximately 40 ravines radiate from the summit and provide pathways for frequent pyroclastic flows and mudslides to the lowlands below. Mayon's most violent eruption occurred in 1814 and was responsible for over 1200 deaths. - The 1951 Peléan eruption of Mount Lamington. Prior to this eruption the peak had not even been recognized as a volcano. Over 3,000 people were killed, and it has become a benchmark for studying large Peléan eruptions. Mount Lamington following the devastating 1951 eruption. Plinian eruptions (or Vesuvian) are a type of volcanic eruption, named for the historical eruption of Mount Vesuvius in 79 of Mount Vesuvius that buried the Roman towns of Pompeii and Herculaneum and, specifically, for its chronicler Pliny the Younger. The process powering Plinian eruptions starts in the magma chamber, where dissolved volatile gases are stored in the magma. The gases vesiculate and accumulate as they rise through the magma conduit. These bubbles agglutinate and once they reach a certain size (about 75% of the total volume of the magma conduit) they explode. The narrow confines of the conduit force the gases and associated magma up, forming an eruptive column. Eruption velocity is controlled by the gas contents of the column, and low-strength surface rocks commonly crack under the pressure of the eruption, forming a flared outgoing structure that pushes the gases even faster. These massive eruptive columns are the distinctive feature of a Plinian eruption, and reach up 2 to 45 km (1 to 28 mi) into the atmosphere. The densest part of the plume, directly above the volcano, is driven internally by gas expansion. As it reaches higher into the air the plume expands and becomes less dense, convection and thermal expansion of volcanic ash drive it even further up into the stratosphere. At the top of the plume, powerful prevailing winds drive the plume in a direction away from the volcano. These highly explosive eruptions are associated with volatile-rich dacitic to rhyolitic lavas, and occur most typically at stratovolcanoes. Eruptions can last anywhere from hours to days, with longer eruptions being associated with more felsic volcanoes. Although they are associated with felsic magma, Plinian eruptions can just as well occur at basaltic volcanoes, given that the magma chamber differentiates and has a structure rich in silicon dioxide. Plinian eruptions are similar to both Vulcanian and Strombolian eruptions, except that rather than creating discrete explosive events, Plinian eruptions form sustained eruptive columns. They are also similar to Hawaiian lava fountains in that both eruptive types produce sustained eruption columns maintained by the growth of bubbles that move up at about the same speed as the magma surrounding them. Regions affected by Plinian eruptions are subjected to heavy pumice airfall affecting an area 0.5 to 50 km3 (0 to 12 cu mi) in size. The material in the ash plume eventually finds its way back to the ground, covering the landscape in a thick layer of many cubic kilometers of ash. However the most dangerous eruptive feature are the pyroclastic flows generated by material collapse, which move down the side of the mountain at extreme speeds of up to 700 km (435 mi) per hour and with the ability to extend the reach of the eruption hundreds of kilometers. The ejection of hot material from the volcano's summit melts snowbanks and ice deposits on the volcano, which mixes with tephra to form lahars, fast moving mudslides with the consistency of wet concrete that move at the speed of a river rapid. Major Plinian eruptive events include: - The AD 79 eruption of Mount Vesuvius buried the Roman towns of Pompeii and Herculaneum under a layer of ash and tephra. It is the model Plinian eruption. Mount Vesuvius has erupted several times since then. Its last eruption was in 1944 and caused problems for the allied armies as they advanced through Italy. It was the report by Pliny that Younger that lead scientists to refer to vesuvian eruptions as "Plinian". - The 1980 eruption of Mount St. Helens in Washington, which ripped apart the volcano's summit, was a Plinian eruption of Volcanic Explosivity Index (VEI) 5. - The strongest types of eruptions, with a VEI of 8, are so-called "Ultra-Plinian" eruptions, such as the most recent one at Lake Toba 74 thousand years ago, which put out 2800 times the material erupted by Mount St. Helens in 1980. - Hekla in Iceland, an example of basaltic Pilian volcanism being its 1947-48 eruption. The past 800 years have been a pattern of violent initial eruptions of pumice followed by prolonged extrusion of basaltic lava from the lower part of the volcano. - Pinatubo in the Philippines on 15 June 1991, which produced 5 km3 (1 cu mi) of dacitic magma, a 40 km (25 mi) high eruption column, and released 17 megatons of sulfur dioxide. Phreatomagmatic eruptions are eruptions that arise from interactions between water and magma. They are driven from thermal contraction (as opposed to magmatic eruptions, which are driven by thermal expansion) of magma when it comes in contact with water. This temperature difference between the two causes violent water-lava interactions that make up the eruption. The products of phreatomagmatic eruptions are believed to be more regular in shape and finer grained than the products of magmatic eruptions because of the differences in eruptive mechanisms. There is debate about the exact nature of phreatomagmatic eruptions, and some scientists believe that fuel-coolant reactions may be more critical to the explosive nature than thermal contraction. Fuel coolant reactions may fragment the volcanic material by propagating stress waves, widening cracks and increasing surface area that ultimetly lead to rapid cooling and explosive contraction-driven eruptions. A Surtseyan eruption (or hydrovolcanic) is a type of volcanic eruption caused by shallow-water interactions between water and lava, named so after its most famous example, the eruption and formation of the island of Surtsey off the coast of Iceland in 1963. Surtseyan eruptions are the "wet" equivalent of ground-based Strombolian eruptions, but because of where they are taking place they are much more explosive. This is because as water is heated by lava, it flashes in steam and expands violently, fragmenting the magma it is in contact with into fine-grained ash. Surtseyan eruptions are the hallmark of shallow-water volcanic oceanic islands, however they are not specifically confined to them. Surtseyan eruptions can happen on land as well, and are caused by rising magma that comes into contact with an aquifer (water-bearing rock formation) at shallow levels under the volcano. The products of Surtseyan eruptions are generally oxidized palagonite basalts (though andesitic eruptions do occur, albeit rarely), and like Strombolian eruptions Surtseyan eruptions are generally continuous or otherwise rhythmic. A distinct defining feature of a Surtseyan eruption is the formation of a pyroclastic surge (or base surge), a ground hugging radial cloud that develops along with the eruption column. Base surges are caused by the gravitational collapse of a vaporous eruptive column, one that is denser overall then a regular volcanic column. The densest part of the cloud is nearest to the vent, resulting a wedge shape. Associated with these laterally moving rings are dune-shaped depositions of rock left behind by the lateral movement. These are occasionally disrupted by bomb sags, rock that was flung out by the explosive eruption and followed a ballistic path to the ground. Accumulations of wet, spherical ash known as accretionary lapilli is another common surge indicator. Over time Surtseyan eruptions tend to form maars, broad low-relief volcanic craters dug into the ground, and tuff rings, circular structures built of rapidly quenched lava. These structures are associated with a single vent eruption, however if eruptions arise along fracture zones a rift zone may be dug out; these eruptions tend to be more violent then the ones forming a tuff ring or maars, an example being the 1886 eruption of Mount Tarawera. Littoral cones are another hydrovolcanic feature, generated by the explosive deposition of basaltic tephra (although they are not truly volcanic vents). They form when lava accumulates within cracks in lava, superheats and explodes in a steam explosion, breaking the rock apart and depositing it on the volcano's flank. Consecutive explosions of this type eventually generate the cone. Volcanoes known to have Surtseyan activity include: - Surtsey, Iceland. The volcano built itself up from depth and emerged above the Atlantic Ocean off the coast of Iceland in 1963. Initial hydrovolcanics were highly explosive, but as the volcano grew out rising lava started to interact less with the water and more with the air, until finally Surtseyan activity waned and became more Strombolian in character. - Ukinrek Maars in Alaska, 1977, and Capelinhos in the Azores, 1957, both examples of above-water Surtseyan activity. - Mount Tarawera in New Zealand erupted along a rift zone in 1886, killing 150 people. - Ferdinandea, a seamount in the Mediterranean Sea, breached sea level in July 1831 and was the source of a dispute over sovereignty between Italy, France, and Great Britain. The volcano did not build tuff cones strongly enough to withstand erosion, and disappeared back below the waves soon after it appeared. - The underwater volcano Hunga Tonga in Tonga breached sea level in 2009. Both of its vents exhibited Surtseyan activity for much of the time. It was also the site of an earlier eruption in May 1988. Submarine eruptions are a type of volcanic eruption that occurs underwater. An estimated 75% of the total volcanic eruptive volume is generated by submarine eruptions near mid ocean ridges alone, however because of the problems associated with detecting deep sea volcanics, they remained virtually unknown until advances in the 1990s made it possible to observe them. Submarine eruptions may produce seamounts which may break the surface to form volcanic islands and island chains. Submarine volcanism is driven by various processes. Volcanoes near plate boundaries and mid-ocean ridges are built by the decompression melting of mantle rock that rises on an upwelling portion of a convection cell to the crustal surface. Eruptions associated with subducting zones, meanwhile, are driven by subducting plates that add volatiles to the rising plate, lowering its melting point. Each process generates different rock; mid-ocean ridge volcanics are primarily basaltic, whereas subduction flows are mostly calc-alkaline, and more explosive and viscous. Spreading rates along mid-ocean ridges vary widely, from 2 cm (0.8 in) per year at the Mid-Atlantic Ridge, to up to 16 cm (6 in) along the East Pacific Rise. Higher spreading rates are a probably cause for higher levels of volcanism. The technology for studying seamount eruptions did not exist until advancements in hydrophone technology made it possible to "listen" to acoustic waves, known as T-waves, released by submarine earthquakes associated with submarine volcanic eruptions. The reason for this is that land-based seismometers cannot detect sea-based earthquakes below a magnitude of 4, but acoustic waves travel well in water and long periods of time. A system in the North Pacific, maintained by the United States Navy and originally intended for the detection of submarines, has detected an event on average every 2 to 3 years. The most common underwater flow is pillow lava, a circular lava flow named after its unusual shape. Less common are glassy, marginal sheet flows, indicative of larger-scale flows. Volcaniclastic sedimentary rocks are common in shallow-water environments. As plate movement starts to carry the volcanoes away from their eruptive source, eruption rates start to die down, and water erosion grinds the volcano down. The final stages of eruption caps the seamount in alkalic flows. There are about 100,000 deepwater volcanoes in the world, although most are beyond the active stage of their life. Some exemplary seamounts are Loihi Seamount, Bowie Seamount, Davidson Seamount, and Axial Seamount. Subglacial eruptions are a type of volcanic eruption characterized by interactions between lava and ice, often under a glacier. The nature of glaciovolcanism dictates that it occurs at areas of high latitude and high altitude. It has been suggested that subglacial volcanoes that are not actively erupting often dump heat into the ice covering them, producing meltwater. This meltwater mix means that subglacial eruptions often generate dangerous jökulhlaups (floods) and lahars. The study of glaciovolcanism is still a relatively new field. Early accounts described the unusual flat-topped steep-sided volcanoes (called tuyas) in Iceland that were suggested to have formed from eruptions below ice. The first English-language paper on the subject was published in 1947 by William Henry Mathews, describing the Tuya Butte field in northwest British Columbia, Canada. The eruptive process that builds these structures, originally inferred in the paper, begins with volcanic growth below the glacier. At first the eruptions resemble those that occur in the deep sea, forming piles of pillow lava at the base of the volcanic structure. Some of the lava shatters when it comes in contact with the cold ice, forming a glassy breccia called hyaloclastite. After a while the ice finally melts into a lake, and the more explosive eruptions of Surtseyan activity begins, building up flanks made up of mostly hyaloclastite. Eventually the lake boils off from continued volcanism, and the lava flows become more effusive and thicken as the lava cools much more slowly, often forming columnar jointing. Well-preserved tuyas show all of these stages, for example Hjorleifshofdi in Iceland. Products of volcano-ice interactions stand as various structures, whose shape is dependent on complex eruptive and environmental interactions. Glacial volcanism is a good indicator of past ice distribution, making it an important climatic marker. Since they are imbedded in ice, as ice retracts worldwide there are concerns that tuyas and other structures may destabalize, resulting in mass landslides. Evidence of volcanic-glacial interactions are evident in Iceland and parts of British Columbia, and it is even possible that they play a role in deglaciation. Glaciovolcanic products have been identified in Iceland, the Canadian province of British Columbia, the U.S. states of Hawaii and Alaska, the Cascade Range of western North America, South America and even on the planet Mars. Volcanoes known to have subglacial activity include: - Mauna Kea in tropical Hawaii. There is evidence of past subglacial eruptive activity on the volcano in the form of a subglacial deposit on its summit. The eruptions originated about 10,000 years ago, during the last ice age, when the summit of Mauna Kea was covered in ice. - In 2008, the British Antarctic Survey reported a volcanic eruption under the Antarctica ice sheet 2,200 years ago. It is believed to be that this was the biggest eruption in Antarctica in the last 10,000 years. Volcanic ash deposits from the volcano were identified through an airborne radar survey, buried under later snowfalls in the Hudson Mountains, close to Pine Island Glacier. - Iceland, well known for both glaciers and volcanoes, is often a site of subglacial eruptions. An example an eruption under the Vatnajökull ice cap in 1996, which occurred under an estimated 2,500 ft (762 m) of ice. - As part of the search for life on Mars, scientists have suggested that there may be subglacial volcanoes on the red planet. Several potential sites of such volcanism have been reviewed, and compared extensively with similar features in Iceland: Viable microbial communities have been found living in deep (-2800 m) geothermal groundwater at 349 K and pressures >300 bar. Furthermore, microbes have been postulated to exist in basaltic rocks in rinds of altered volcanic glass. All of these conditions could exist in polar regions of Mars today where subglacial volcanism has occurred. Phreatic eruptions (or steam-blast eruptions) are a type of eruption driven by the expansion of steam. When cold ground or surface water come into contact with hot rock or magma it superheats and explodes, fracturing the surrounding rock and thrusting out a mixture of steam, water, ash, volcanic bombs, and volcanic blocks. The distinguishing feature of phreatic explosions is that they only blast out fragments of pre-existing solid rock from the volcanic conduit; no new magma is erupted. Because they are driven by the cracking of rock strata under pressure, phreatic activity does not always result in an eruption; if the rock face is strong enough to withstand the explosive force, outright eruptions may not occur, although cracks in the rock will probably develop and weaken it, furthering future eruptions. Often a precursor of future volcanic activity, phreatic eruptions are generally weak, although there have been exceptions. Some phreatic events may be triggered by earthquake activity, another volcanic precursor, and they may also travel along dike lines. Phreatic eruptions form base surges, lahars, avalanches, and volcanic block "rain." They may also release deadly toxic gas able to suffocate anyone in range of the eruption. Volcanoes known to exhibit phreatic activity include: - Mount St. Helens, which exhibited phreatic activity just prior to its catastrophic 1980 eruption (which was itself Plinian). - Taal Volcano, Philippines, 1965. - La Soufrière of Guadeloupe (Lesser Antilles), 1975-1976 activity. - Soufrière Hills volcano on Montserrat, West Indies, 1995–2012. - Poás Volcano, has frequent geyser like phreatic eruptions from its crater lake. - Mount Bulusan, well known for its sudden phreatic eruptions. - Mount Ontake, all historical eruptions of this volcano have been phreatic including the deadly 2014 eruption. - List of currently erupting volcanoes - List of Quaternary volcanic eruptions - Prediction of volcanic activity - Timetable of major worldwide volcanic eruptions - Heiken, G. and Wohletz, K. Volcanic Ash. University of California Press. p. 246. - "VHP Photo Glossary: Effusive Eruption". USGS. 29 December 2009. Retrieved 3 August 2010. - "Volcanoes of Canada: Volcanic eruptions". Geological Survey of Canada. Natural Resources Canada. 2 April 2009. Retrieved 3 August 2010. - "How Volcanoes Work: Hawaiian Eruptions". San Diego State University. Retrieved 2 August 2010. - "How Volcanoes Work: Hydrovolcic Eruptions". San Diego State University. Retrieved 4 August 2010. - Ruprecht P, Plank T. Feeding andesitic eruptions with a high-speed connection from the mantle. Nature. 2013;500(7460):68-72. - "How Volcanoes Work: Eruption Variabilty". San Diego State University. Retrieved 3 August 2010. - "How Volcanoes Work: Basaltic Lava". San Diego State University. Retrieved 2 August 2010. - "Oshima". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 2 August 2010. - "How Volcanoes Work: Strombolian Eruptions". San Diego State University. Retrieved 29 July 2010. - Mike Burton, Patrick Allard, Filippo Muré, Alessandro La Spina (2007). "Magmatic Gas Composition Reveals the Source Depth of Slug-Driven Strombolian Explosive Activity". Science (American Association for the Advancement of Science) 317 (5835): 227–230. Bibcode:2007Sci...317..227B. doi:10.1126/science.1141900. ISSN 1095-9203. Retrieved 30 July 2010. - Cain, Fraser (22 April 2010). "Strombolian Eruption". Universe Today. Retrieved 30 July 2010. - Seach, John. "Mt Etna Volcano Eruptions - John Seach". Old eruptions. Volcanolive. Retrieved 30 July 2010. - Seach, John. "Mt Etna Volcano Eruptions - John Seach". Recent eruptions. Volcanolive. Retrieved 30 July 2010. - "Erebus". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 31 July 2010. - Kyle, P. R. (Ed.), Volcanological and Environmental Studies of Mount Erebus, Antarctica, Antarctic Research Series, American Geophysical Union, Washington DC, 1994. - "Stromboli". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 31 July 2010. - "How Volcanoes Work: Vulcanian Eruptions". San Diego State University. Retrieved 1 August 2010. - Cain, Fraser. "Vulcanian Eruptions". Universe Today. Retrieved 1 August 2010. - "How Volcanoes Work: Sakurajima Volcano". San Diego State University. Retrieved 1 August 2010. - "VHP Photo Glossary: Vulcanian eruption". USGS. Retrieved 1 August 2010. - Cain, Fraser. "Pelean Eruption". Universe Today. Retrieved 2 August 2010. - Donald Hyndman and David Hyndman (April 2008). Natural Hazards and Disasters. Cengage Learning. pp. 134–135. ISBN 978-0-495-31667-1. Retrieved 2 August 2010. - Nelson, Stephan A. (30 September 2007). "Volcanoes, Magma, and Volcanic Eruptions". Tulane University. Retrieved 2 August 2010. - Richard V. Fisher and Grant Heiken (1982). "Mt. Pelée, Martinique: May 8 and 20 pyroclastic flows and surges". Journal of Volcanology and Geothermal Research 13 (3-4): 339–371. Bibcode:1982JVGR...13..339F. doi:10.1016/0377-0273(82)90056-7. - "How Volcanoes Work: Mount Pelée Eruption (1902)". San Diego State University. Retrieved 1 August 2010. - "Mayon". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 2 August 2010. - "Lamington: Photo Gallery". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 2 August 2010. - "How Volcanoes Work: Plinian Eruptions". San Diego State University. Retrieved 3 August 2010. - "How Volcanoes Work: Eruption Model". San Diego State University. Retrieved 3 August 2010. - Cain, Fraser. "Plinian Eruption". Universe Today. Retrieved 3 August 2010. - "How Volcanoes Work: Calderas". San Diego State University. Retrieved 3 August 2010. - Stephen Self, Jing-Xia Zhao, Rick E. Holasek, Ronnie C. Torres, and Alan J. King. "The Atmospheric Impact of the 1991 Mount Pinatubo Eruption". USGS. Retrieved 3 August 2010. - A.B. Starostin, A.A. Barmin, and O.E. Melnik (May 2005). "A transient model for explosive and phreatomagmatic eruptions". Journal of Volcanology and Geothermal Research. Volcanic Eruption Mechanisms - Insights from intercomparison of models of conduit processes 143 (1-3): 133–151. Bibcode:2005JVGR..143..133S. doi:10.1016/j.jvolgeores.2004.09.014. Retrieved 4 August 2010. - "X. Classification of Volcanic Eruptions: Surtseyan Eruptions". Lecture Notes. University of Alabama. Retrieved 5 August 2010. - Alwyn Scarth and Jean-Claude Tanguy (31 May 2001). Volcanoes of Europe. Oxford University Press. p. 264. ISBN 978-0-19-521754-4. Retrieved 5 August 2010. - "Hunga Tonga-Hunga Ha'apai: Index of Monthly Reports". Global Volcanism Program. Smithsonian National Museum of Natural History. Retrieved 5 August 2010. - Chadwick, Bill (10 January 2006). "Recent Submarine Volcanic Eruptions". Vents Program. NOAA. Retrieved 5 August 2010. - Hubert Straudigal and David A Clauge. "The Geological History of Deep-Sea Volcanoes: Biosphere, Hydrosphere, and Lithosphere Interactions" (PDF). Oceanography. Seamounts Special Issue (Oceanography Society) 32 (1). Retrieved 4 August 2010. - Paul Wessel, David T. Sandwell, Seung-Sep Kim. "The Global Seamount Census" (PDF). Oceanography. Seamounts Special Issue (Oceanography Society) 23 (1). ISSN 1042-8275. Retrieved 25 June 2010. - "Glaciovolcanism - University of British Columbia". University of British Columbia. Retrieved 5 August 2010. - Black, Richard (20 January 2008). "Ancient Antarctic eruption noted". BBC News. Retrieved 5 August 2010. - Alden, Andrew. "Tuya or Subglacial Volcano, Iceland". about.com. Retrieved 5 August 2010. - "Kinds of Volcanic Eruptions". Volcano World. Oregon State University. Retrieved 5 August 2010. - "Iceland's subglacial eruption". Hawaiian Volcano Observatory. USGS. 11 October 1996. Retrieved 5 August 2010. - "Subglacial Volcanoes On Mars". Space Daily. 27 June 2001. Retrieved 5 August 2010. - Leonid N. Germanovich and Robert P. Lowell (1995). "The mechanism of phreatic eruptions". Journal of Geophysical Research. Solid Earth (American Geophysical Union) 100 (B5): 8417–8434. Bibcode:1995JGR...100.8417G. doi:10.1029/94JB03096. Retrieved 7 August 2010. - "VHP Photo Glossary: Phreatic eruption". USGS. 17 July 2008. Retrieved 6 August 2010. - Watson, John (5 February 1997). "Types of volcanic eruptions". USGS. Retrieved 7 August 2010. - "Phreatic Eruptions - John Seach". Volcano World. Retrieved 6 August 2010. - Grant Heiken and Kenneth Wohletz (1985). Volcanic Ash. University of California Press. p. 258. ISBN 978-0-520-05241-3. Retrieved 5 August 2010. - A.B. Starostin, A.A. Barmin and O.E. Melnik (May 2005). "A transient model for explosive and phreatomagmatic eruptions". Journal of Volcanology and Geothermal Research. Volcanic Eruption Mechanisms - Insights from intercomparison of models of conduit processes 143 (1-3): 133–151. Bibcode:2005JVGR..143..133S. doi:10.1016/j.jvolgeores.2004.09.014. Retrieved 5 August 2010. - Pyle, D. M. (January 1989). "The thickness, volume and grainsize of tephra fall deposits". Bulletin of Volcanology 51 (1): 1–15. Bibcode:1989BVol...51....1P. doi:10.1007/BF01086757. - Colleen M. Riley, William I. Rose, Gregg J. S. Bluth (28 October 2003). "Quantitative shape measurements of distal volcanic ash" (PDF). Journal of Geophysical Research 108 (B10). Bibcode:2003JGRB..108.2504R. doi:10.1029/2001JB000818. - William Henry Mathews (September 1947). ""Tuyas," flat-topped volcanoes in northern British Columbia". American Journal of Science 245 (9): 560–570. doi:10.2475/ajs.245.9.560.. This is the original landmark paper by William Henry Mathews that first described tuyas and subglacial eruptions. |Wikimedia Commons has media related to Diagrams of volcanic eruptions.| - USGS Hawaiian Volcano Observatory (HVO) homepage. USGS. - Distinguishing eruptive types. - How Volcanoes Work. San Diego State University.

https://en.wikipedia.org/wiki/Types_of_volcanic_eruptions

First, to recap, we know that there is a roughly 11 year sunspot cycle where we go from essentially no sunspots on the surface of the Sun to a period of high sunspot activity and then back to essentially no sunspots again. It also turns out that sunspot position on the Sun is not random. Graphing where sunspots appear during the sunspot cycle yields the famous "Butterfly Diagram" showing that, at the start of the sunspot cycle, sunspots initially appear near ±30° of latitude and then, as the cycle progresses, sunspots appear closer and closer to the equator. Keep in mind that the sunspots themselves don't move (although they appear to from Earth during the rotation of the Sun) but form, typically exist for a couple of weeks or more, and then dissipate. Where they pop up on the surface of the Sun, however, is what changes during the sunspot cycle and is reflected on the butterfly diagram above. What else have we learned about sunspots? Well, a sunspot can be over 50,000 km across (for reference, the diameter of the Earth is about 12,750 km) and consists of two visible parts - the darker central umbra (Latin for shadow) and the lighter surrounding penumbra (the Latin prefix means almost or nearly). The image below shows a large sunspot group (AR 1339 from November 4) with a filtered telescope. Note the clearly visible umbra and penumbra for each sunspot. Sunspots are darker than the surrounding Sun because they're cooler. The surface of the Sun is around 6000 K (over 10,000° F) and sunspots are 1500 K (~2250°) or more cooler. So, while they appear dark compared to the rest of the Sun, if you could somehow take a sunspot off the Sun and place it by itself in space, it would glow brightly in the sky! The next image shows a sunspot a bit more dramatically. This image was taken in August of 2010 at the Big Bear Solar Observatory in the San Bernadino Mountains of California. This image was taken with a special filter called a hydrogen alpha (Ha) filter. It's a filter that only passes a narrow bandwidth of visible light at a wavelength of 656 nm (6.56 x 10-7 m). This is the energy given off by electrons in a hydrogen atom falling from the 3rd to the 2nd orbital. The Sun, being a big ball of mostly hydrogen gas, gives off a lot of energy at this wavelength and these filters bring out a lot of detail on the "surface" (photosphere) of the Sun. This image shows a granulation around the sunspot. Those are the tops of convection cells where hot gases are "bubbling" up from deeper in the Sun. Here's a neat animation of solar activity. This convection of hot, ionized gas (plasma) generates the Sun's magnetic field. Because the Sun is a big ball of gas, it doesn't rotate at the same speed everywhere - it takes about a 9 days less to rotate at the equator (~25 days) than it does near the poles (~34 days). This differential rotation leads to magnetic flux tubes in the convection zone of the Sun getting twisted up (they actually behave much like rubber bands). This inhibits convection and leads to the development of a cooler sunspot (don't ask me to explain this any better since I'm not a solar physicist). Magnetic lines of force also pop up above the photosphere often leading to the formation of two sunspots of opposite magnetic polarity (one where the magnetic field emerges from the photosphere, the other where it reenters the Sun). Sunspots have about 1000 times more magnetic energy than surrounding areas of the Sun. Below is a magnetogram image of the Sun for today (November 13). Black indicates areas where the Sun's magnetic lines of force are coming toward us and white indicates areas where the Sun's magnetic lines of force are moving away from us. Compare this to the visible image of the Sun for today and you can see that major black and white areas correlate with the positions of sunspots. Anyway, models to explain this 22-year cycle are all based on the differential rotation of the Sun affecting the internal convection and thus changing the magnetic field over time. If the Sun rotated faster or slower, or the convection zone was thicker or thinner, or convection was faster or slower, this cycle would be different. The details are messy and I don't understand them myself (phrases like "...regeneration of the poloidal field by lifting and twisting a toroidal flux tube by helical turbulence..." follow by a page of equations are typical in the literature. So, you may be thinking, sunspots are cool looking features on the surface of the Sun that we don't fully understand but do they have any significance for us here on Earth? Yes, as a matter of fact they do. That will be the topic of the next post.

http://hudsonvalleygeologist.blogspot.com/2011/11/why-do-we-have-sunspots.html

Temporal range: Early Eocene to present |Eastern screech owl| some 25, see text Striginae sensu Sibley & Ahlquist The true owls or typical owls (family Strigidae) are one of the two generally accepted families of owls, the other being the barn owls (Tytonidae). The Sibley-Ahlquist taxonomy unites the Caprimulgiformes with the owl order; here, the typical owls are a subfamily Striginae. This is unsupported by more recent research (see Cypselomorphae for details), but the relationships of the owls in general are still unresolved. This large family comprises around 189 living species in 25 genera. The typical owls have a cosmopolitan distribution and are found on every continent except Antarctica. While typical owls (hereafter referred to simply as owls) vary greatly in size, with the smallest species, the elf owl, being a hundredth the size of the largest, the Eurasian eagle-owl and Blakiston's fish owl, owls generally share an extremely similar body plan. They tend to have large heads, short tails, cryptic plumage and round facial discs around the eyes. The family is generally arboreal (with a few exceptions like the burrowing owl) and obtain their food on the wing. The wings are large, broad, rounded and long. Like for other birds of prey, in many owl species females are larger than males. Because of their nocturnal habits they tend not to exhibit sexual dimorphism in their plumage. The feathers are soft and the base of each is downy, allowing for silent flight. The toes and tarsus are feathered in some species, and more so in species at higher latitudes. Numerous species of owl in the genus Glaucidium and the northern hawk-owl have eye patches on the backs of their heads, apparently to convince other birds they are being watched at all times. Numerous nocturnal species have ear-tufts, feathers on the sides of the head that are thought to have a camouflage function, breaking up the outline of a roosting bird. The feathers of the facial disc are arranged in order to increase sound delivered to the ears. Hearing in owls is highly sensitive and the ears are asymmetrical allowing the owl to localise a sound. In addition to hearing owls have massive eyes relative to their body size. Contrary to popular belief, however, owls cannot see well in extreme dark and are able to see fine in the day. Owls are generally nocturnal and spend much of the day roosting. They are often perceived as tame since they will allow people to approach quite closely before taking flight, but they are instead attempting to avoid detection. The cryptic plumage and inconspicuous locations adopted are an effort to avoid predators and mobbing by small birds. - Genus Megascops – screech-owls, some 20 species - Genus Otus – scops-owls; probably paraphyletic, about 45 species - Genus Pyrroglaux – Palau owl - Genus Margarobyas – bare-legged owl or Cuban screech-owl - Genus Ptilopsis – white-faced owls, 2 species - Genus Mimizuku – giant scops-owl or Mindanao eagle-owl - Genus Bubo – horned owls, eagle-owls and fish-owls; paraphyletic with Nyctea, Ketupa and Scotopelia, some 25 species - Genus Strix – earless owls, some 19 species, including 4 that were previously classified as Ciccaba - Genus Ciccaba – the 4 species have been transferred to Strix - Genus Lophostrix – crested owl - Genus Jubula – maned owl - Genus Pulsatrix – spectacled owls, 3 species - Genus Surnia – northern hawk-owl - Genus Glaucidium – pygmy owls, about 30–35 species - Genus Xenoglaux – long-whiskered owlet - Genus Micrathene – elf owl - Genus Athene – 2–4 species (depending on whether Speotyto and Heteroglaux are included or not) - Genus Aegolius – saw-whet owls, 4 species - Genus Ninox – Australasian hawk-owls, some 20 species - Genus Uroglaux – Papuan hawk-owl - Genus Pseudoscops – Jamaican owl and possibly striped owl - Genus Asio – eared owls, 6–7 species - Genus Nesasio – fearful owl - Genus Mascarenotus – Mascarene owls, 3 species (extinct c. 1850) - Genus Sceloglaux – laughing owl (extinct 1914?) - Genus Grallistrix – stilt-owls, 4 species - Genus Ornimegalonyx – Caribbean giant owls, 1–2 species - Cuban giant owl, Ornimegalonxy oteroi - Ornimegalonyx sp. – probably subspecies of O. oteroi - Genus Asphaltoglaux - Mioglaux (Late Oligocene? – Early Miocene of WC Europe) – includes "Bubo" poirreiri - Intulula (Early/Middle Miocene of WC Europe) – includes "Strix/Ninox" brevis - Alasio (Middle Miocene of Vieux-Collonges, France) – includes "Strix" collongensis - "Otus/Strix" wintershofensis – fossil (Early/Middle Miocene of Wintershof West, Germany) – may be close to extant genus Ninox - "Strix" edwardsi – fossil (Middle Miocene of Grive-Saint-Alban, France) - "Asio" pygmaeus – fossil (Early Pliocene of Odessa, Ukraine) - Strigidae gen. et sp. indet. UMMP V31030 (Rexroad Late Pliocene of Kansas, USA) – Strix/Bubo? - Ibiza owl, Strigidae gen. et sp. indet. – prehistoric (Late Pleistocene/Holocene of Es Pouàs, Ibiza) The supposed fossil heron "Ardea" lignitum (Late Pliocene of Germany) was apparently a strigid owl, possibly close to Bubo. The Early–Middle Eocene genus Palaeoglaux from west-central Europe is sometimes placed here, but given its age it is probably better considered its own family for the time being. - Marks, J. S.; Cannings, R.J. and Mikkola, H. (1999). "Family Strigidae (Typical Owls)". In del Hoyo, J.; Elliot, A. & Sargatal, J. (eds.) (1999). Handbook of the Birds of the World. Volume 5: Barn-Owls to Hummingbirds. Lynx Edicions. ISBN 84-87334-25-3 - Earhart, Caroline M. and Johnson, Ned K. (1970). "Size Dimorphism and Food Habits of North American Owls". Condor 72 (3): 251–264. doi:10.2307/1366002. - Kelso L & Kelso E (1936). "The Relation of Feathering of Feet of American Owls to Humidity of Environment and to Life Zones". Auk 53 (1): 51–56. doi:10.2307/4077355. - Olson, p. 131 - Feduccia, J. Alan; Ford, Norman L. (1970). "Some birds of prey from the Upper Pliocene of Kansas" (PDF). Auk 87 (4): 795–797. doi:10.2307/4083714. - Sánchez Marco, Antonio (2004). "Avian zoogeographical patterns during the Quaternary in the Mediterranean region and paleoclimatic interpretation" (PDF). Ardeola 51 (1): 91–132. - Olson, p. 167 - Olson, Storrs L. (1985). The fossil record of birds. In: Farner, D.S.; King, J.R. & Parkes, Kenneth C. (eds.): Avian Biology 8: 79–238. Academic Press, New York. |Wikimedia Commons has media related to Strigidae.|

https://en.wikipedia.org/wiki/Strigidae

We take our ABC's for granted, learning 26 letters in a precise order from our youngest days. When introduced to a second or third language later in life we may realize that even similar tongues to English contain slightly different alphabets--the Spanish ñ, the French ç--despite the fact that they evolved from the same roots. Historical variation in the English alphabet seems largely glossed over in contemporary education, but identifying some of the "missing letters" can help explain a few historical puzzles. First, there's ampersand, considered the 27th letter of the English alphabet until about 150 years ago. It's name comes from its position at the end of the ABC's: The word “ampersand” came many years later when “&” was actually part of the English alphabet. In the early 1800s, school children reciting their ABCs concluded the alphabet with the &. It would have been confusing to say “X, Y, Z, and.” Rather, the students said, “and per se and.” “Per se” means “by itself,” so the students were essentially saying, “X, Y, Z, and by itself and.” Over time, “and per se and” was slurred together into the word we use today: ampersand. When a word comes about from a mistaken pronunciation, it’s called a mondegreen. Before the introduction of the Latin alphabet after the Roman conquest of Britain, Anglo-Saxon had an alphabet all its own known as furthorc. In the ensuing battle of cultural power politics Anglo-Saxon lost out. Collateral damage included the letter "thorn," pictured at right, pronounced with the hard "th" sound. It was replaced by the humble Y, always ready to do double duty in that ambiguous no-man's-land between consonants and vowels. This explains the anachronistic use of Y in titles like "Ye Olde English Shoppe"--it's just another spelling of "the." On Friday we'll take a look at another missing letter, the long s (resembling "f"). For a sneak peek and a list of nine other extinct English letters, check out this article from MentalFloss (via @johndcook).

http://mattdickenson.com/2013/02/20/micro-institutions-everywhere-the-english-alphabet/

Slavery in ancient Greece Slavery was a very common practice in Ancient Greek history, as in other places of the time. It is estimated that the majority of Athenian citizens owned at least one slave; most ancient writers considered slavery natural and even necessary. This paradigm was notably questioned in Socratic dialogues; the Stoics produced the first recorded condemnation of slavery. Modern historiographical practice distinguishes chattel (personal possession) slavery from land-bonded groups such as the penestae of Thessaly or the Spartan helots, who were more like medieval serfs (an enhancement to real estate). The chattel slave is an individual deprived of liberty and forced to submit to an owner, who may buy, sell or lease them like any other chattel. The academic study of slavery in ancient Greece is beset by significant methodological problems. Documentation is disjointed and very fragmented, focusing primarily on Athens. No treatises are specifically devoted to the subject, and jurisprudence was interested in slavery only inasmuch as it provided a source of revenue. Comedies and tragedies represented stereotypes while iconography made no substantial differentiation between slaves and craftsmen. - 1 Terminology - 2 Origins of slavery - 3 Economic role - 4 Demographics - 5 Status of slaves - 6 Slavery conditions - 7 Views of Greek slavery - 8 Notes - 9 References - 10 Further reading - 11 External links The ancient Greeks had several words for slaves, which leads to textual ambiguity when they are studied out of their proper context. In Homer, Hesiod and Theognis of Megara, the slave was called δμώς / dmōs. The term has a general meaning but refers particularly to war prisoners taken as booty (in other words, property). During the classical period, the Greeks frequently used ἀνδράποδον / andrapodon, (literally, "one with the feet of a man") as opposed to τετράποδον / tetrapodon, "quadruped", or livestock. The most common word is δοῦλος / doulos, used in opposition to "free man" (ἐλεύθερος / eleútheros); an earlier form of the former appears in Mycenaean inscriptions as do-e-ro, "male slave" (or "servant", "bondman"; Linear B: 𐀈𐀁𐀫), or do-e-ra, "female slave" (or "maid-servant", "bondwoman"; Linear B: ). The verb δουλεὐω (which survives in Modern Greek, meaning "work") can be used metaphorically for other forms of dominion, as of one city over another or parents over their children. Finally, the term οἰκέτης / oiketēs was used, meaning "one who lives in house", referring to household servants. Other terms used were less precise and required context: - θεράπων / therapōn – At the time of Homer, the word meant "squire" (Patroclus was referred to as the therapōn of Achilles and Meriones that of Idomeneus); during the classical age, it meant "servant". - ἀκόλουθος / akolouthos – literally, "the follower" or "the one who accompanies". Also, the diminutive ἀκολουθίσκος, used for page boys. - παῖς / pais – literally "child", used in the same way as "houseboy", also used in a derogatory way to call adult slaves. - σῶμα / sōma – literally "body", used in the context of emancipation. Origins of slavery Slaves were present through the Mycenaean civilization, as documented in numerous tablets unearthed in Pylos 140. Two legal categories can be distinguished: "slaves (εοιο)" and "slaves of the god (θεοιο)", the god in this case probably being Poseidon. Slaves of the god are always mentioned by name and own their own land; their legal status is close to that of freemen. The nature and origin of their bond to the divinity is unclear. The names of common slaves show that some of them came from Kythera, Chios, Lemnos or Halicarnassus and were probably enslaved as a result of piracy. The tablets indicate that unions between slaves and freemen were common and that slaves could work and own land. It appears that the major division in Mycenaean civilization was not between a free individual and a slave but rather if the individual was in the palace. · There is no continuity between the Mycenaean era and the time of Homer, where social structures reflected those of the Greek dark ages. The terminology differs: the slave is no longer do-e-ro (doulos) but dmōs. In the Iliad, slaves are mainly women taken as booty of war, while men were either ransomed or killed on the battlefield. In the Odyssey, the slaves also seem to be mostly women. These slaves were servants and sometimes concubines. There were some male slaves, especially in the Odyssey, a prime example being the swineherd Eumaeus. The slave was distinctive in being a member of the core part of the oikos ("family unit", "household"): Laertes eats and drinks with his servants; in the winter, he sleeps in their company. The term dmōs is not considered pejorative, and Eumaeus, the "divine" swineherd, bears the same Homeric epithet as the Greek heroes. Slavery remained, however, a disgrace. Eumaeus himself declares, "Zeus, of the far-borne voice, takes away the half of a man's virtue, when the day of slavery comes upon him". · It is difficult to determine when slave trading began in the archaic period. In Works and Days (8th century BC), Hesiod owns numerous dmōes although their status is unclear. The presence of douloi is confirmed by lyric poets such as Archilochus or Theognis of Megara. According to epigraphic evidence, the homicide law of Draco (c. 620 BC) mentioned slaves. According to Plutarch, Solon (c. 594-593 BC) forbade slaves from practising gymnastics and pederasty. By the end of the period, references become more common. Slavery becomes prevalent at the very moment when Solon establishes the basis for Athenian democracy. Classical scholar Moses Finley likewise remarks that Chios, which, according to Theopompus, was the first city to organize a slave trade, also enjoyed an early democratic process (in the 6th century BC). He concludes that "one aspect of Greek history, in short, is the advance hand in hand, of freedom and slavery." All activities were open to slaves with the exception of politics. For the Greeks, politics was the only activity worthy of a citizen, the rest being relegated wherever possible to non-citizens. It was status that was of importance, not activity. The principal use of slavery was in agriculture, the foundation of the Greek economy. Some small landowners might own one slave, or even two. An abundant literature of manuals for landowners (such as the Economy of Xenophon or that of Pseudo-Aristotle) confirms the presence of dozens of slaves on the larger estates; they could be common labourers or foremen. The extent to which slaves were used as a labour force in farming is disputed. It is certain that rural slavery was very common in Athens, and that ancient Greece did not know of the immense slave populations found on the Roman latifundia. Slave labour was prevalent in mines and quarries, which had large slave populations, often leased out by rich private citizens. The strategos Nicias leased a thousand slaves to the silver mines of Laurium in Attica; Hipponicos, 600; and Philomidès, 300. Xenophon indicates that they received one obolus per slave per day, amounting to 60 drachmas per year. This was one of the most prized investments for Athenians. The number of slaves working in the Laurium mines or in the mills processing ore has been estimated at 30,000. Xenophon suggested that the city buy a large number of slaves, up to three state slaves per citizen, so that their leasing would assure the upkeep of all the citizens. Slaves were also used as craftsmen and tradespersons. As in agriculture, they were used for labour that was beyond the capability of the family. The slave population was greatest in workshops: the shield factory of Lysias employed 120 slaves, and the father of Demosthenes owned 32 cutlers and 20 bedmakers. Slaves were also employed in the home. The domestic's main role was to stand in for his master at his trade and to accompany him on trips. In time of war he was batman to the hoplite. The female slave carried out domestic tasks, in particular bread baking and textile making. Only the poorest citizens did not possess a domestic slave. It is difficult to estimate the number of slaves in ancient Greece, given the lack of a precise census and variations in definitions during that era. It is certain that Athens had the largest slave population, with as many as 80,000 in the 6th and 5th centuries BC, on average three or four slaves per household. In the 5th century BC, Thucydides remarked on the desertion of 20,890 slaves during the war of Decelea, mostly tradesmen. The lowest estimate, of 20,000 slaves, during the time of Demosthenes, corresponds to one slave per family. Between 317 BC and 307 BC, the tyrant Demetrius Phalereus ordered a general census of Attica, which arrived at the following figures: 21,000 citizens, 10,000 metics and 400,000 slaves. The orator Hypereides, in his Against Areistogiton, recalls that the effort to enlist 15,000 male slaves of military age led to the defeat of the Southern Greeks at the Battle of Chaeronea (338 BC), which corresponds to the figures of Ctesicles. According to the literature, it appears that the majority of free Athenians owned at least one slave. Aristophanes, in Plutus, portrays poor peasants who have several slaves; Aristotle defines a house as containing freemen and slaves. Conversely, not owning even one slave was a clear sign of poverty. In the celebrated discourse of Lysias For the Invalid, a cripple pleading for a pension explains "my income is very small and now I'm required to do these things myself and do not even have the means to purchase a slave who can do these things for me." However, the huge slave populations of the Romans were unknown in ancient Greece. When Athenaeus cites the case of Mnason, friend of Aristotle and owner of a thousand slaves, this appears to be exceptional. Plato, owner of five slaves at the time of his death, describes the very rich as owning 50 slaves. Sources of supply There were four primary sources of slaves: war, in which the defeated would become slaves to the victorious unless a more objective outcome was reached; piracy (at sea); banditry (on land); and international trade. By the rules of war of the period, the victor possessed absolute rights over the vanquished, whether they were soldiers or not. Enslavement, while not systematic, was common practice. Thucydides recalls that 7,000 inhabitants of Hyccara in Sicily were taken prisoner by Nicias and sold for 120 talents in the neighbouring village of Catania. Likewise in 348 BC the population of Olynthus was reduced to slavery, as was that of Thebes in 335 BC by Alexander the Great and that of Mantineia by the Achaean League. The existence of Greek slaves was a constant source of discomfort for free Greeks. The enslavement of cities was also a controversial practice. Some generals refused, such as the Spartans Agesilaus II and Callicratidas. Some cities passed accords to forbid the practice: in the middle of the 3rd century BC, Miletus agreed not to reduce any free Knossian to slavery, and vice versa. Conversely, the emancipation by ransom of a city that had been entirely reduced to slavery carried great prestige: Cassander, in 316 BC, restored Thebes. Before him, Philip II of Macedon enslaved and then emancipated Stageira. Piracy and banditry Piracy and banditry provided a significant and consistent supply of slaves, though the significance of this source varied according to era and region. Pirates and brigands would demand ransom whenever the status of their catch warranted it. Whenever ransom was not paid or not warranted, captives would be sold to a trafficker. In certain areas, piracy was practically a national specialty, described by Thucydides as "the old-fashioned" way of life. Such was the case in Acarnania, Crete, and Aetolia. Outside of Greece, this was also the case with Illyrians, Phoenicians, and Etruscans. During the Hellenistic period, Cilicians and the mountain peoples from the coasts of Anatolia could also be added to the list. Strabo explains the popularity of the practice among the Cilicians by its profitability; Delos, not far away, allowed for "moving a myriad of slaves daily". The growing influence of the Roman Republic, a large consumer of slaves, led to development of the market and an aggravation of piracy. In the 1st century BC, however, the Romans largely eradicated piracy to protect the Mediterranean trade routes. There was slave trade between kingdoms and states of the wider region. The fragmentary list of slaves confiscated from the property of the mutilators of the Hermai mentions 32 slaves whose origin have been ascertained: 13 came from Thrace, 7 from Caria, and the others came from Cappadocia, Scythia, Phrygia, Lydia, Syria, Ilyria, Macedon and Peloponnese. Local professionals sold their own people to Greek slave merchants. The principal centres of the slave trade appear to have been Ephesus, Byzantium, and even faraway Tanais at the mouth of the Don. Some "barbarian" slaves were victims of war or localised piracy, but others were sold by their parents. There is a lack of direct evidence of slave traffic, but corroborating evidence exists. Firstly, certain nationalities are consistently and significantly represented in the slave population, such as the corps of Scythian archers employed by Athens as a police force—originally 300, but eventually nearly a thousand. Secondly, the names given to slaves in the comedies often had a geographical link; thus Thratta, used by Aristophanes in The Wasps, The Acharnians, and Peace, simply signified Thracian woman. Finally, the nationality of a slave was a significant criterion for major purchasers; the ancient advice was not to concentrate too many slaves of the same origin in the same place, in order to limit the risk of revolt. It is also probable that, as with the Romans, certain nationalities were considered more productive as slaves than others. The price of slaves varied in accordance with their ability. Xenophon valued a Laurion miner at 180 drachmas; while a workman at major works was paid one drachma per day. Demosthenes' father's cutlers were valued at 500 to 600 drachmas each. Price was also a function of the quantity of slaves available; in the 4th century BC they were abundant and it was thus a buyer's market. A tax on sale revenues was levied by the market cities. For instance, a large slave market was organized during the festivities at the temple of Apollo at Actium. The Acarnanian League, which was in charge of the logistics, received half of the tax proceeds, the other half going to the city of Anactorion, of which Actium was a part. Buyers enjoyed a guarantee against latent defects; the transaction could be invalidated if the bought slave turned out to be crippled and the buyer had not been warned about it. Curiously, it appears that the Greeks did not "breed" their slaves, at least during the Classical Era, though the proportion of houseborn slaves appears to have been rather large in Ptolemaic Egypt and in manumission inscriptions at Delphi. Sometimes the cause of this was natural; mines, for instance, were exclusively a male domain. On the other hand, there were many female domestic slaves. The example of African slaves in the American South on the other hand demonstrates that slave populations can multiply. This incongruity remains relatively unexplained. Xenophon advised that male and female slaves should be lodged separately, that "…nor children born and bred by our domestics without our knowledge and consent—no unimportant matter, since, if the act of rearing children tends to make good servants still more loyally disposed, cohabiting but sharpens ingenuity for mischief in the bad." The explanation is perhaps economic; even a skilled slave was cheap, so it may have been cheaper to purchase a slave than to raise one. Additionally, childbirth placed the slave-mother's life at risk, and the baby was not guaranteed to survive to adulthood. Houseborn slaves (oikogeneis) often constituted a privileged class. They were, for example, entrusted to take the children to school; they were "pedagogues" in the first sense of the term. Some of them were the offspring of the master of the house, but in most cities, notably Athens, a child inherited the status of its mother. Status of slaves The Greeks had many degrees of enslavement. There was a multitude of categories, ranging from free citizen to chattel slave, and including Penestae or helots, disenfranchised citizens, freedmen, bastards, and metics. The common ground was the deprivation of civic rights. - had no rights - Right to own property - Authority over the work of another - Power of punishment over another - Legal rights and duties (liability to arrest and/or arbitrary punishment, or to litigate) - Familial rights and privileges (marriage, inheritance, etc.) - Possibility of social mobility (manumission or emancipation, access to citizen rights) - Religious rights and obligations - Military rights and obligations (military service as servant, heavy or light soldier, or sailor) Athenian slaves were the property of their master (or of the state), who could dispose of them as he saw fit. He could give, sell, rent, or bequeath them. A slave could have a spouse and children, but the slave family was not recognized by the state, and the master could scatter the family members at any time. Slaves had fewer judicial rights than citizens and were represented by their master in all judicial proceedings. A misdemeanour that would result in a fine for the free man would result in a flogging for the slave; the ratio seems to have been one lash for one drachma. With several minor exceptions, the testimony of a slave was not admissible except under torture. Slaves were tortured in trials because they often remained loyal to their master. A famous example of trusty slave was Themistocles's Persian slave Sicinnus (the counterpart of Ephialtes of Trachis), who, despite his Persian origin, betrayed Xerxes and helped Athenians in the Battle of Salamis. Despite torture in trials, the Athenian slave was protected in an indirect way: if he was mistreated, the master could initiate litigation for damages and interest (δίκη βλάβης / dikē blabēs). Conversely, a master who excessively mistreated a slave could be prosecuted by any citizen (γραφὴ ὕβρεως / graphē hybreōs); this was not enacted for the sake of the slave, but to avoid violent excess (ὕβρις / hubris). Isocrates claimed that "not even the most worthless slave can be put to death without trial"; the master's power over his slave was not absolute. Draco's law apparently punished with death the murder of a slave; the underlying principle was: "was the crime such that, if it became more widespread, it would do serious harm to society?" The suit that could be brought against a slave's killer was not a suit for damages, as would be the case for the killing of cattle, but a δίκη φονική (dikē phonikē), demanding punishment for the religious pollution brought by the shedding of blood. In the 4th century BC, the suspect was judged by the Palladion, a court which had jurisdiction over unintentional homicide; the imposed penalty seems to have been more than a fine but less than death—maybe exile, as was the case in the murder of a Metic. However, slaves did belong to their master's household. A newly-bought slave was welcomed with nuts and fruits, just like a newly-wed wife. Slaves took part in most of the civic and family cults; they were expressly invited to join the banquet of the Choes, second day of the Anthesteria, and were allowed initiation into the Eleusinian Mysteries. A slave could claim asylum in a temple or at an altar, just like a free man. The slaves shared the gods of their masters and could keep their own religious customs if any. Slaves could not own property, but their masters often let them save up to purchase their freedom, and records survive of slaves operating businesses by themselves, making only a fixed tax-payment to their masters. Athens also had a law forbidding the striking of slaves: if a person struck what appeared to be a slave in Athens, that person might find himself hitting a fellow-citizen, because many citizens dressed no better. It astonished other Greeks that Athenians tolerated back-chat from slaves. Athenian slaves fought together with Athenian freemen at the battle of Marathon, and the monuments memorialize them. It was formally decreed before the battle of Salamis that the citizens should "save themselves, their women, children, and slaves". Slaves had special sexual restrictions and obligations. For example, a slave could not engage free boys in pederastic relationships ("A slave shall not be the lover of a free boy nor follow after him, or else he shall receive fifty blows of the public lash."), and they were forbidden from the palaestrae ("A slave shall not take exercise or anoint himself in the wrestling-schools."). Both laws are attributed to Solon. Fathers wanting to protect their sons from unwanted advances provided them with a slave guard, called a paidagogos, to escort the boy in his travels. The sons of vanquished foes would be enslaved and often forced to work in male brothels, as in the case of Phaedo of Elis, who at the request of Socrates was bought and freed from such an enterprise by the philosopher's rich friends. On the other hand, it is attested in sources that the rape of slaves was persecuted, at least occasionally. Slaves in Gortyn In Gortyn, in Crete, according to a code engraved in stone dating to the 6th century BC, slaves (doulos or oikeus) found themselves in a state of great dependence. Their children belonged to the master. The master was responsible for all their offences, and, inversely, he received amends for crimes committed against his slaves by others. In the Gortyn code, where all punishment was monetary, fines were doubled for slaves committing a misdemeanour or felony. Conversely, an offence committed against a slave was much less expensive than an offence committed against a free person. As an example, the rape of a free woman by a slave was punishable by a fine of 200 staters (400 drachms), while the rape of a non-virgin slave by another slave brought a fine of only one obolus (a sixth of a drachm). Slaves did have the right to possess a house and livestock, which could be transmitted to descendants, as could clothing and household furnishings. Their family was recognized by law: they could marry, divorce, write a testament and inherit just like free men. A specific case: debt slavery Prior to its interdiction by Solon, Athenians practiced debt enslavement: a citizen incapable of paying his debts became "enslaved" to the creditor. The exact nature of this dependency is a much controversial issue among modern historians: was it truly slavery or another form of bondage? However, this issue primarily concerned those peasants known as "hektēmoroi" working leased land belonging to rich landowners and unable to pay their rents. In theory, those so enslaved would be liberated when their original debts were repaid. The system was developed with variants throughout the Near East and is cited in the Bible. Solon put an end to it with the σεισάχθεια / seisachtheia, liberation of debts, which prevented all claim to the person by the debtor and forbade the sale of free Athenians, including by themselves. Aristotle in his Constitution of the Athenians quotes one of Solon's poems: And many a man whom fraud or law had sold Far from his god-built land, an outcast slave, I brought again to Athens; yea, and some, Exiles from home through debt’s oppressive load, Speaking no more the dear Athenian tongue, But wandering far and wide, I brought again; And those that here in vilest slavery (douleia) Crouched ‘neath a master’s (despōtes) frown, I set them free. Though much of Solon's vocabulary is that of "traditional" slavery, servitude for debt was at least different in that the enslaved Athenian remained an Athenian, dependent on another Athenian, in his place of birth. It is this aspect which explains the great wave of discontent with slavery of the 6th century BC, which was not intended to free all slaves but only those enslaved by debt. The reforms of Solon left two exceptions: the guardian of an unmarried woman who had lost her virginity had the right to sell her as a slave, and a citizen could "expose" (abandon) unwanted newborn children. The practice of manumission is confirmed to have existed in Chios from the 6th century BC. It probably dates back to an earlier period, as it was an oral procedure. Informal emancipations are also confirmed in the classical period. It was sufficient to have witnesses, who would escort the citizen to a public emancipation of his slave, either at the theatre or before a public tribunal. This practice was outlawed in Athens in the middle of the 6th century BC to avoid public disorder. The practice became more common in the 4th century BC and gave rise to inscriptions in stone which have been recovered from shrines such as Delphi and Dodona. They primarily date to the 2nd and 1st centuries BC, and the 1st century AD. Collective manumission was possible; an example is known from the 2nd century BC in the island of Thasos. It probably took place during a period of war as a reward for the slaves' loyalty, but in most cases the documentation deals with a voluntary act on the part of the master (predominantly male, but in the Hellenistic period also female). The slave was often required to pay for himself an amount at least equivalent to his street value. To this end they could use their savings or take a so-called "friendly" loan (ἔρανος / eranos) from their master, a friend or a client like the hetaera Neaira did. Emancipation was often of a religious nature, where the slave was considered to be "sold" to a deity, often Delphian Apollo, or was consecrated after his emancipation. The temple would receive a portion of the monetary transaction and would guarantee the contract. The manumission could also be entirely civil, in which case the magistrate played the role of the deity. The slave's freedom could be either total or partial, at the master's whim. In the former, the emancipated slave was legally protected against all attempts at re-enslavement—for instance, on the part of the former master's inheritors. In the latter case, the emancipated slave could be liable to a number of obligations to the former master. The most restrictive contract was the paramone, a type of enslavement of limited duration during which time the master retained practically absolute rights. In regard to the city, the emancipated slave was far from equal to a citizen by birth. He was liable to all types of obligations, as one can see from the proposals of Plato in The Laws: presentation three times monthly at the home of the former master, forbidden to become richer than him, etc. In fact, the status of emancipated slaves was similar to that of metics, the residing foreigners, who were free but did not enjoy a citizen's rights. Spartan citizens used helots, a dependent group collectively owned by the state. It is uncertain whether they had chattel slaves as well. There are mentions of people manumitted by Spartans, which was supposedly forbidden for helots, or sold outside of Lakonia: the poet Alcman; a Philoxenos from Cytherea, reputedly enslaved with all his fellow citizens when his city was conquered, later sold to an Athenian; a Spartan cook bought by Dionysius the Elder or by a king of Pontus, both versions being mentioned by Plutarch; and the famous Spartan nurses, much appreciated by Athenian parents. Some texts mention both slaves and helots, which seems to indicate that they were not the same thing. Plato in Alcibiades I cites "the ownership of slaves, and notably helots" among the Spartan riches, and Plutarch writes about "slaves and helots". Finally, according to Thucydides, the agreement that ended the 464 BC revolt of helots stated that any Messenian rebel who might hereafter be found within the Peloponnese was "to be the slave of his captor", which means that the ownership of chattel slaves was not illegal at that time. Most historians thus concur that chattel slaves were indeed used in the Greek city-state of Sparta, at least after the Lacedemonian victory of 404 BC against Athens, but not in great numbers and only among the upper classes. As was in the other Greek cities, chattel slaves could be purchased at the market or taken in war. It is difficult to appreciate the condition of Greek slaves. According to Aristotle, the daily routine of slaves could be summed up in three words: "work, discipline, and feeding". Xenophon's advice is to treat slaves as domestic animals, that is to say punish disobedience and reward good behaviour. For his part, Aristotle prefers to see slaves treated as children and to use not only orders but also recommendations, as the slave is capable of understanding reasons when they are explained. Greek literature abounds with scenes of slaves being flogged; it was a means of forcing them to work, as were control of rations, clothing, and rest. This violence could be meted out by the master or the supervisor, who was possibly also a slave. Thus, at the beginning of Aristophanes' The Knights (4–5), two slaves complain of being "bruised and thrashed without respite" by their new supervisor. However, Aristophanes himself cites what is a typical old saw in ancient Greek comedy: "He also dismissed those slaves who kept on running off, or deceiving someone, or getting whipped. They were always led out crying, so one of their fellow slaves could mock the bruises and ask then: 'Oh you poor miserable fellow, what's happened to your skin? Surely a huge army of lashes from a whip has fallen down on you and laid waste your back?'" The condition of slaves varied very much according to their status; the mine slaves of Laureion and the pornai (brothel prostitutes) lived a particularly brutal existence, while public slaves, craftsmen, tradesmen and bankers enjoyed relative independence. In return for a fee (ἀποφορά / apophora) paid to their master, they could live and work alone. They could thus earn some money on the side, sometimes enough to purchase their freedom. Potential emancipation was indeed a powerful motivator, though the real scale of this is difficult to estimate. Ancient writers considered that Attic slaves enjoyed a "peculiarly happy lot": Pseudo-Xenophon deplores the liberties taken by Athenian slaves: "as for the slaves and Metics of Athens, they take the greatest licence; you cannot just strike them, and they do not step aside to give you free passage". This alleged good treatment did not prevent 20,000 Athenian slaves from running away at the end of the Peloponnesian War at the incitement of the Spartan garrison at Attica in Decelea. These were principally skilled artisans (kheirotekhnai), probably among the better-treated slaves. The title of a 4th-century comedy by Antiphanes, The Runaway-catcher (Δραπεταγωγός), suggests that slave flight was not uncommon. Conversely, there are no records of a large-scale Greek slave revolt comparable to that of Spartacus in Rome. It can probably be explained by the relative dispersion of Greek slaves, which would have prevented any large-scale planning. Slave revolts were rare, even in Rome. Individual acts of rebellion of slaves against their master, though scarce, are not unheard of; a judicial speech mentions the attempted murder of his master by a boy slave, not 12 years old. Views of Greek slavery Very few authors of antiquity call slavery into question. To Homer and the pre-classical authors, slavery was an inevitable consequence of war. Heraclitus states that "War is the father of all, the king of all ... he turns some into slaves and sets others free". During the classical period, the main justification for slavery was economic. From a philosophical point of view, the idea of "natural" slavery emerged at the same time; thus, as Aeschylus states in The Persians, the Greeks "[o]f no man are they called the slaves or vassals", while the Persians, as Euripides states in Helen, "are all slaves, except one" — the Great King. Hippocrates theorizes about this latent idea at the end of the 5th century BC. According to him, the temperate climate of Anatolia produced a placid and submissive people. This explanation is reprised by Plato, then Aristotle in Politics, where he develops the concept of "natural slavery": "for he that can foresee with his mind is naturally ruler and naturally master, and he that can do these things with his body is subject and naturally a slave." As opposed to an animal, a slave can comprehend reason but "…has not got the deliberative part at all." In parallel, the concept that all men, whether Greek or barbarian, belonged to the same race was being developed by the Sophists and thus that certain men were slaves although they had the soul of a freeman and vice versa. Aristotle himself recognized this possibility and argued that slavery could not be imposed unless the master was better than the slave, in keeping with his theory of "natural" slavery. The Sophists concluded that true servitude was not a matter of status but a matter of spirit; thus, as Menander stated, "be free in the mind, although you are slave: and thus you will no longer be a slave". This idea, repeated by the Stoics and the Epicurians, was not so much an opposition to slavery as a trivialisation of it. The Greeks could not comprehend an absence of slaves. Slaves exist even in the "Cloudcuckooland" of Aristophanes' The Birds as well as in the ideal cities of Plato's Laws or Republic. The utopian cities of Phaleas of Chalcedon and Hippodamus of Miletus are based on the equal distribution of property, but public slaves are used respectively as craftsmen and land workers. The "reversed cities" placed women in power or even saw the end of private property, as in Lysistrata or Assemblywomen, but could not picture slaves in charge of masters. The only societies without slaves were those of the Golden Age, where all needs were met without anyone having to work. In this type of society, as explained by Plato, one reaped generously without sowing. In Telekleides' Amphictyons barley loaves fight with wheat loaves for the honour of being eaten by men. Moreover, objects move themselves—dough kneads itself, and the jug pours itself. Similarly, Aristotle said that slaves would not be necessary "if every instrument could accomplish its own work... the shuttle would weave and the plectrum touch the lyre without a hand to guide them", like the legendary constructs of Daedalus and Hephaestus. Society without slaves is thus relegated to a different time and space. In a "normal" society, one needs slaves. Slavery in Greek antiquity has long been an object of apologetic discourse among Christians, who are typically awarded the merit of its collapse. From the 16th century the discourse became moralizing in nature. The existence of colonial slavery had significant impact on the debate, with some authors lending it civilizing merits and others denouncing its misdeeds. Thus Henri-Alexandre Wallon in 1847 published a History of Slavery in Antiquity among his works for the abolition of slavery in the French colonies. In the 19th century, a politico-economic discourse emerged. It concerned itself with distinguishing the phases in the organisation of human societies and correctly identifying the place of Greek slavery. The influence of Marx is decisive; for him the ancient society was characterized by development of private ownership and the dominant (and not secondary as in other pre-capitalist societies) character of slavery as a mode of production. The Positivists represented by the historian Eduard Meyer (Slavery in Antiquity, 1898) were soon to oppose the Marxist theory. According to him slavery was the foundation of Greek democracy. It was thus a legal and social phenomenon, and not economic. Current historiography developed in the 20th century; led by authors such as Joseph Vogt, it saw in slavery the conditions for the development of elites. Conversely, the theory also demonstrates an opportunity for slaves to join the elite. Finally, Vogt estimates that modern society, founded on humanist values, has surpassed this level of development. In 2011, Greek slavery remains the subject of historiographical debate, on two questions in particular: can it be said that ancient Greece was a "slave society", and did Greek slaves comprise a social class? - A traditional pose in funerary steles, see for instance Felix M. Wassermann, "Serenity and Repose: Life and Death on Attic Tombstones" The Classical Journal, Vol. 64, No. 5, p.198. - J.M.Roberts, The New Penguin History of the World, p.176–177, 223 - Chantraine, s.v. δμώς. - For instance Odyssey 1:398, where Telemachus mentions "the slaves that goodly Odysseus won for [him]". - Used once by Homer in Iliad 7:475 to refer to prisoners taken in war; the line was athetized by Aristarchus of Samothrace following Zenodotus and Aristophanes of Byzantium, see Kirk, p.291. - Chantraine, s.v. ἀνερ. - Definition from LSJ. - Mycenean transliterations can be confusing and do not directly reflect pronunciation; for clarification see the article about Linear B. - Chantraine, s.v. δοῦλος. See also Mactoux (1981). - Chantraine, s.v. οἰκος. - Iliad, 16:244 and 18:152. - Iliad, 23:113. - Chantraine, s.v. θεράπων. - Chantraine, s.v. ἀκόλουθος. - Chantraine, s.v. παῖς. - Cartledge, p.137. - Chantraine, s.v. σῶμα. - Garlan, p.32. - Burkert, p.45. - Garlan, p.35. - Mele, pp.115–155. - Garlan, p.36. - For instance Chryseis (1:12–3, 29–30, 111–5), Briseis (2:688–9), Diomede (6:654–5), Iphis (6:666–8) and Hecamede (11:624–7). - See in the Iliad the pleas of Adrastus the Trojan (1:46–50), the sons of Antimachus (11:131–5) and Lycaon (21:74–96), all begging for mercy in exchange of a ransom. - There are 50 of them in Ulysses' house (22:421) and in Alcinous' house (7:103). - Before his fight with Achilles, Hector predicts for his wife Andromache a life of bondage and mentions weaving and water-fetching (6:454–8). In the Odyssey, servants tend the fire (20:123), prepare the suitors' feast (1:147), grind wheat (7:104, 20:108–9), make the bed (7:340–2) and take care of the guests. - In the Iliad, Chryseis sleeps with Agamemnon, Briseis and Diomede with Achilles, Iphis with Patroclus. In the Odyssey, twelve female servants sleep with the suitors (20:6–8) against Euryclea's direct orders (22:423–425). - Odyssey, 16:140–1. - Odyssey, 11:188–91. - Odyssey, 14:3. - Garlan, p.43. - Odyssey 17:322–323. Online version of Butcher-Lang 1879 translation. - For instance Works and Days, 405. - "κατὰ ταὐτὰ φόνοθ δίκας εἷναι δοῦλον κτείναντι ἢ ἐλεὐτερον." Dareste, Haussoulier and Reinach, 4, 5, 8. - Life of Solon, 1:6. - Apud Athenaeus, 6:265bc = FGrH 115, fgt.122. - Finley (1997), pp.170–171. - Finley (1997), p.180. - Finley (1997), p.148. - Finley (1997), p.149. - Jameson argues in favour of a very large use of slaves; Wood (1983 and 1988) disputes it. - Finley (1997), p.150. - Poroi (On Revenues), 4. - Lauffer, p.916. - Demosthenes, 12:8–19. - Demosthenes, Against Aphobos, 11:9. - Finley (1997), pp.151–152. - Jones, pp.76–79. - Ctesicles, apud Athenaeus 6:272c. - Ctesicles was the author of a history preserved as two fragments in the Athenaeus. - Politics, 252a26–b15. - Lysias, For the invalid, 3. - Athenaeus, 6:264d. - Republic, 9:578d–e. - Thucydide, 8:40, 2. - See Ducrey for further reading. - Thucydides, 6:62 and 7:13. - Garlan, p. 57. - Plutarch, Life of Agesilaus, 7:6. - Xenophon, Hellenica, 1:6, 14. - Diodorus Siculus, 19:53,2. - Plutarch, Life of Alexander, 7:3. - The Greeks made little differentiation between pirates and bandits, both being called lēstai or peiratai. Brulé (1978a), p.2. - See Ormerod, Brulé (1978b) and Gabrielsen for further reading. - Finley (1997), p.230. - Thucydides, 1:5, 3. - Strabo, 14:5, 2. - Brulé (1978a), p.6. - Brulé (1978a), pp.6–7. - Pritchett and Pippin (1956), p.278 and Pritchett (1961), p.27. - Herodotus, 5:6; Philostratus II, Life of Apollonius Tyana, 18:7, 12. - Plassart, pp.151–213. - During the Classical and Hellenistic periods, it was the master who named the slave; this could be the master's name, an ethnic name as mentioned above, a name from their native area (Manes for Lydian, Midas for a Phrygian, etc.), a historical name (Alexander, Cleopatra, etc.). In short, a slave could carry practically any name, but barbarian names could only be given to slaves. Masson, pp.9—21. - Plato, Laws, 777cd; Pseudo-Aristotle, Economics, 1:5. - Garlan, p.61. - Circa 216 BC. Inscriptiones Graecae IX 1², 2, 583. - Hypereides, Against Athenogenes, 15 and 22. - Garlan, p.59. - Finley (1997), p.155. - The Economist, IX. Trans H. G. Dakyns, accessed 16 May 2006. - Pritchett and Pippin, pp.276–281. - Garlan, p.58. Finley (1997), p.154–155 remains doubtful. - Garlan, p.58. - Carlier, p.203. - Finley (1997), p.147. - Finley (1997), pp.165–89. - Garlan, p.47. - Antiphon, First Tetralogy, 2:7, 4:7; Demosthenes, Against Pantenos, 51 (2) and Against Evergos, 14, 15, 60. - For instance Lycurgus, Against Leocrates, 29. - Aeschines, Against Timarchus, 17. - Panathenaicus, 181.http://www.perseus.tufts.edu/hopper/text?doc=Perseus%3Atext%3A1999.01.0144%3Aspeech%3D12%3Asection%3D181 - Morrow, p.212. - Lycurgus, Against Leocrates, 66. - Morrow, p.213. - Aristotle, Constitution of the Athenians, 57:3. - Burkert, p.259. - Carlier, p.204. - Old Oligarch, Constitution of the Athenians, 10. - Pausanias, 1:29, 6. - Plutarch, Life of Themistocles, 10:4–5. - Aeschines, Against Timarchos 1.138–139 - Diogenes Laertius, Lives of the Philosophers, 2.105 - Wilhelm Kroll "Knabenliebe" in Pauly-Wissowa, Realencyclopaedie der klassischen Altertumswissenschaft, vol. 11, cols. 897–906 - Lévy (1995), p.178. - Finley (1997), p.200. - Finley (1997), p.201. - Lévy (1995), p.179. - Aristotle, Constitution of the Athenians, See also 1:2 and Plutarch, Life of Solon, 13:2. - Literally, "six-parters" or "sixthers", because they owed either one-sixth or five-sixths (depending on the interpretation) of their harvest. See Von Fritz for further reading. - Deuteronomy, 15:12–17. - Constitution of the Athenians 12:4. Trans. by Sir Frederic Kenyon, accessed 15 May 2006. - Finley, p.174. - Finley (1997), p.160. - Plutarch, Life of Solon 23.2. - Brulé (1992), p.83. - Garlan, p.79. - Garlan, p.80. - Dunant and Pouilloux, pp.35–37, no.173. - Demosthenes, Against Neaira, 59:29–32. - See Foucart for further reading. - Garlan, p.82. - Garlan, p.83. - Garlan, p.84. - Laws, 11:915 a–c. - Garlan, p.87. - Herakleides Lembos, fgt. 9 Dilts and Suidas, s.v. Ἀλκμάν. - Suidas, s.v. Φιλόξενος. - Life of Lycurgus, 12:13. - Life of Lycurgus, 16:5; Life of Alcibiades, 5:3. - "…ἀνδραπόδων κτήσει τῶν τε ἄλλων καὶ τῶν εἱλωτικῶν", Alcibiades I, 122d. - "…δοὐλοις καὶ Εἴλωσι", Comp. Lyc. et Num., 2. - Oliva, pp.172–173; Ducat, p.55; Lévy (2003), pp.112–113. - Economics, 1344a35. - Xenophon, Economics, 13:6. - Politics, I, 3, 14. - Peace, v.743–749. Trans. Ian Johnston, 2006, accessed 17 May 06. - Garlan, p.147. - Garlan, p.148. - Finley (1997), p. 165. - Morrow, p.210. See Plato, The Republic, 8:563b; Demosthenes, Third Philippic, 3; Aeschines, Against Timarchos, 54; Aristophanes, Assemblywomen, 721–22 and Plautus, Stichus, 447–50. - Constitution of the Athenians, I, 10. - Thucydides (7:27). - Apud Athenaeus, 161e. - Cartledge, p.139. - Garlan, p.180. - Finley (1997), p.162–3. - Antiphon, On the Murder of Herodes, 69. - Heraclitus, frag.53. - Mactoux (1980), p.52. - The Persians, v.242. Trans. ed. Herbert Weir Smyth, accessed 17 May 2006. - Helen, v.276. - Hippocratic corpus, Of Airs, Waters, and Places (Peri aeron hydaton topon), 23. - Republic, 4:435a–436a. - Politics, 7:1327b. - Politics, 1:2, 2. Trans. H. Rackham, accessed 17 May 2006. - Politics, 1:13, 17. - John D. Bury and Russell Meiggs (4th ed. 1975): A History of Greece to the Death of Alexander the Great. New York: St. Martin's Press, page 375 - For instance Hippias of Elis apud Platon, Protagoras, 337c; Antiphon, Pap. Oxyr., 9:1364. - An idea already expressed by Euripides, Ion, 854–856frag.831. - Politics, 1:5, 10. - Menander, frag. 857. - Garlan, p.130. - Republic, 10:469b sq. and 470c. - Apud Aristotle, Politics, 1267b. - Apud Aristotle, Politics, 1268a. - Politics, 271a–272b. - Apud Athenaeus, 268 b–d. - Aristotle, Politics, Book 1 Part 4 - Garlan, p.8. - Garlan, p.10–13. - Garlan, p.13–14. - Garlan, p.19–20. - Garlan, p.201. - This article draws heavily on the Esclavage en Grèce antique article in the French-language Wikipedia, which was accessed in the version of 17 May 2006. - (French) Brulé, P. (1978a) "Signification historique de la piraterie grecque ", Dialogues d'histoire ancienne no.4 (1978), pp. 1–16. - (French) Brulé, P. (1992) "Infanticide et abandon d'enfants", Dialogues d'histoire ancienne no.18 (1992), pp. 53–90. - Burkert, W. Greek Religion. Oxford: Blackwell Publishing, 1985. ISBN 0-631-15624-0, originally published as Griechische Religion der archaischen und klassischen Epoche. Stuttgart: Verlag W. Kohlhammer, 1977. - (French) Carlier, P. Le IVe siècle grec jusqu'à la mort d'Alexandre. Paris: Seuil, 1995. ISBN 2-02-013129-3 - Cartledge, P.. "Rebels and Sambos in Classical Greece", Spartan Reflections. Berkeley: University of California Press, 2003, p. 127–152 ISBN 0-520-23124-4 - (French) Chantraine, P. Dictionnaire étymologique de la langue grecque. Paris: Klincksieck, 1999 (new edition). ISBN 2-252-03277-4 - (French) Dareste R., Haussoullier B., Reinach Th. Recueil des inscriptions juridiques grecques, vol.II. Paris: E. Leroux, 1904. - (French) Ducat, Jean. Les Hilotes, BCH suppl.20. Paris: publications of the École française d'Athènes, 1990 ISBN 2-86958-034-7 - (French) Dunant, C. and Pouilloux, J. Recherches sur l'histoire et les cultes de Thasos II. Paris: publications of the École française d'Athènes, 1958. - Finley, M. (1997). Économie et société en Grèce ancienne. Paris: Seuil, 1997 ISBN 2-02-014644-4, originally published as Economy and Society in Ancient Greece. London: Chatto and Windus, 1981. - Garlan, Y. Les Esclaves en Grèce ancienne. Paris: La Découverte, 1982. 1982 ISBN 2-7071-2475-3, translated in English as Slavery in Ancient Greece. Ithaca, N.Y.: Cornell University Press, 1988 (1st edn. 1982) ISBN 0-8014-1841-0 - Kirk, G.S. (editor). The Iliad: a Commentary, vol.II (books 5–8). Cambridge: Cambridge University Press, 1990. ISBN 0-521-28172-5 - Jameson, M.H. "Agriculture and Slavery in Classical Athens", Classical Journal, no.73 (1977–1978), pp. 122–145. - Jones, A.H.M.. Athenian Democracy. Oxford: Blackwell Publishing, 1957. - (German) Lauffer, S. "Die Bergwerkssklaven von Laureion", Abhandlungen no.12 (1956), pp. 904–916. - (French) Lévy, E. (1995). La Grèce au Ve siècle de Clisthène à Socrate. Paris: Seuil, 1995 ISBN 2-02-013128-5 - (French) Lévy, E. (2003). Sparte. Paris: Seuil, 2003 ISBN 2-02-032453-9 - (French) Mactoux, M.-M. (1980). Douleia: Esclavage et pratiques discursives dans l'Athènes classique. Paris: Belles Lettres, 1980. ISBN 2-251-60250-X - (French) Mactoux, M.-M. (1981). "L'esclavage comme métaphore : douleo chez les orateurs attiques", Proceedings of the 1980 GIREA Workshop on Slavery, Kazimierz, 3–8 November 1980, Index, 10, 1981, pp. 20–42. - (French) Masson, O. "Les noms des esclaves dans la Grèce antique", Proceedings of the 1971 GIREA Workshop on Slavery, Besançon, 10–11 mai 1971. Paris: Belles Lettres, 1973, pp. 9–23. - (French) Mele, A. "Esclavage et liberté dans la société mycénienne", Proceedings of the 1973 GIREA Workshop on Slavery, Besançon 2–3 mai 1973. Paris: Les Belles Lettres, 1976. - Morrow, G.R. "The Murder of Slaves in Attic Law", Classical Philology, Vol. 32, No. 3 (Jul., 1937), pp. 210–227. - Oliva, P. Sparta and her Social Problems. Prague: Academia, 1971. - (French) Plassart, A. "Les Archers d'Athènes," Revue des études grecques, XXVI (1913), pp. 151–213. - Pomeroy, S.B. Goddesses, Whores, Wives and Slaves. New York: Schoken, 1995. ISBN 0-8052-1030-X - Pritchett, W.K. and Pippin, A. (1956). "The Attic Stelai, Part II", Hesperia, Vol.25, No.3 (Jul.–Sep., 1956), pp. 178–328. - Pritchett (1961). "Five New Fragments of the Attic Stelai", Hesperia, Vol.30, No. 1 (Jan.–Mar., 1961), pp. 23–29. - Wood, E.M. (1983). "Agriculture and Slavery in Classical Athens", American Journal of Ancient History No.8 (1983), pp. 1–47. - Von Fritz, K. "The Meaning of ἙΚΤΗΜΟΡΟΣ", The American Journal of Philology, Vol.61, No.1 (1940), pp. 54–61. - Wood, E.M. (1988). Peasant-Citizen and Slave: The Foundations of Athenian Democracy. New York: Verso, 1988 ISBN 0-86091-911-0. - General studies - Bellen, H., Heinen H., Schäfer D., Deissler J., Bibliographie zur antiken Sklaverei. I: Bibliographie. II: Abkurzungsverzeichnis und Register, 2 vol. Stuttgart: Steiner, 2003. ISBN 3-515-08206-9 - Bieżuńska-Małowist I. La Schiavitù nel mondo antico. Naples: Edizioni Scientifiche Italiane, 1991. - Finley, M.: - Garnsey, P. Ideas of Slavery from Aristotle to Augustine. Cambridge: Cambridge University Press, 1996. ISBN 0-521-57433-1 - De Ste-Croix, G.E.M. The Class Struggle in the Ancient Greek World. London: Duckworth; Ithaca, N.Y.: Cornell University Press, 1981. ISBN 0-8014-1442-3 - Vidal-Naquet, P.: - "Women, Slaves and Artisans", third part of The Black Hunter : Forms of Thought and Forms of Society in the Greek World. Baltimore: Johns Hopkins University Press, 1988 (1st edn. 1981). ISBN 0-8018-5951-4 - with Vernant J.-P. Travail et esclavage en Grèce ancienne. Bruxelles: Complexe, "History" series, 2006 (1st edn. 1988). ISBN 2-87027-246-4 - Wiedemann, T. Greek and Roman Slavery. London: Routledge, 1989 (1st edn. 1981). ISBN 0-415-02972-4 - Westermann, W.L. The Slave Systems of Greek and Roman Antiquity. Philadelphia: The American Philosophical Society, 1955. - Specific studies - Brulé, P. (1978b). La Piraterie crétoise hellénistique, Belles Lettres, 1978. ISBN 2-251-60223-2 - Brulé, P. and Oulhen, J. (dir.). Esclavage, guerre, économie en Grèce ancienne. Hommages à Yvon Garlan. Rennes: Presses universitaires de Rennes, "History" series, 1997. ISBN 2-86847-289-3 - Ducrey, P. Le traitement des prisonniers de guerre en Grèce ancienne. Des origines à la conquête romaine. Paris: De Boccard, 1968. - Foucart, P. "Mémoire sur l'affranchissement des esclaves par forme de vente à une divinité d'après les inscriptions de Delphes", Archives des missions scientifiques et littéraires, 2nd series, vol.2 (1865), pp. 375–424. - Hunt, P. Slaves, Warfare, and Ideology in the Greek Historians. Cambridge: Cambridge University Press, 1998. ISBN 0-521-58429-9 - Ormerod, H.A. Piracy in the Ancient World. Liverpool: Liverpool University Press, 1924. - Gabrielsen, V. "La piraterie et le commerce des esclaves", in E. Erskine (ed.), Le Monde hellénistique. Espaces, sociétés, cultures. 323-31 av. J.-C.. Rennes: Presses Universitaires de Rennes, 2004, pp. 495–511. ISBN 2-86847-875-1 |Wikimedia Commons has media related to Slavery in Ancient Greece.| - (French) GIREA, The International Group for Research on Slavery in Antiquity (in French) - Greek law bibliographic database at Nomoi - Documents on Greek slavery on the Ancient History Sourcebook. - Manumission records of women at Delphi at attalus.org - (French) Index thématiques de l'esclavage et de la dépendance Subject index on slavery and related topics, by author. - (French) Bibliothèque numérique ISTA Free library

https://en.wikipedia.org/wiki/Slavery_in_Ancient_Greece

Plaçage was a recognized extralegal system in French and Spanish slave colonies of North America (including the Caribbean) by which ethnic European men entered into the equivalent of common-law marriages with women of color, of African, Native American and mixed-race descent. The term comes from the French placer meaning "to place with". The women were not legally recognized as wives but were known as placées; their relationships were recognized among the free people of color as mariages de la main gauche or left-handed marriages. They became institutionalized with contracts or negotiations that settled property on the woman and her children, and in some cases gave them freedom if enslaved. The system flourished throughout the French and Spanish colonial periods, reaching its zenith during the latter, between 1769 and 1803. It was most practiced in New Orleans, where planter society had created enough wealth to support the system. It also took place in the Latin-influenced cities of Natchez and Biloxi, Mississippi; Mobile, Alabama; St. Augustine and Pensacola, Florida; as well as Saint-Domingue (now the Republic of Haiti). Plaçage became associated with New Orleans as part of its cosmopolitan society. History and development of the plaçage system The plaçage system developed from the predominance of men among early colonial populations, who took women as consorts from Native Americans and enslaved Africans. Later there developed a class of free people of color in Louisiana, and especially New Orleans, during the colonial years, from whom wealthy men would choose. In this period there was a shortage of European women, as the colonies were dominated in the early day by male explorers and colonists. Given the harsh conditions in Louisiana, persuading women to follow the men was not easy. France recruited willing farm- and city-dwelling women, known as casket or casquette girls, because they brought all their possessions to the colonies in a small trunk or casket. France also sent women convicted along with their debtor husbands, and in 1719, deported 209 women felons "who were of a character to be sent to the French settlement in Louisiana." (France also relocated young women orphans known as King's Daughters (French: filles du roi) to their colonies for marriage: to both Canada and Louisiana.) Historian Joan Martin maintains that there is little documentation that "casket girls", considered among the ancestors of white French Creoles, were brought to Louisiana. The Ursuline order of nuns supposedly chaperoned the casket girls until they married, but the order has denied they followed this practice. Martin suggests this was a myth, and that interracial relationships occurred from the beginning of the encounter among Europeans, Native Americans and Africans. She also writes that some Creole families who today identify as white had ancestors during the colonial period who were African or multiracial, and whose descendants married white over generations. Through warfare and raids, Native American women were often captured to be traded, sold, or taken as wives. At first, the colony generally imported male Africans to use as slave labor because of the heavy work of clearing to develop plantations. Over time, it also imported African female slaves. Marriage between the races was forbidden according to the Code Noir of the eighteenth century, but interracial sex continued. The upper class European men during this period often did not marry until their late twenties or early thirties. Premarital sex with an intended white bride, especially if she was of high rank, was not permitted socially. Free people of color White male colonists, often the younger sons of noblemen, military men, and planters, who needed to accumulate some wealth before they could marry, took women of color as consorts before marriage. Merchants and administrators also followed this practice if they were wealthy enough. A white man might rape a slave as young as twelve. When the women bore children, they were sometimes emancipated along with their children. Both the woman and her children might take the surnames of the man. When Creole men reached an age when they were expected to marry, some also kept their relationships with their placées, but this was less common. A wealthy white Creole man could have two (or more) families: one legal, and the other not. Their mixed-race children became the nucleus of the class of free people of color or gens de couleur libres in Louisiana and Saint-Domingue. After the Haitian Revolution in the late 18th and early 19th centuries, many refugees came to New Orleans, adding a new wave of French-speaking free people of color. During the period of French and Spanish rule, the gens de couleur came to constitute a third class in New Orleans and other former French cities - between the white Creoles and the mass of black slaves. They had certain status and rights, and often acquired education and property. Later their descendants became leaders in New Orleans, holding political office in the city and state, and becoming part of what developed as the African-American middle class in the United States. By 1788, 1500 Creole women of color and black women were being maintained by white men. Certain customs had evolved. It was common for a wealthy, married Creole to live primarily outside New Orleans on his plantation with his white family. He often kept a second address in the city to use for entertaining and socializing among the white elite. He had built or bought a house for his placée and their children. She and her children were part of the society of Creoles of color. The white world might not recognize the placée as a wife legally and socially, but she was recognized as such among the Creoles of color. Some of the women acquired slaves and plantations. Particularly during the Spanish colonial era, a woman might be listed as owning slaves; these were sometimes relatives who she intended to free after earning enough money to buy their freedom. While in New Orleans (or other cities), the man would cohabit with the placée as an official 'boarder' at her Creole cottage or house. Many were located near Rampart Street in New Orleans-—once the demarcation line or wall between the city and the frontier. Other popular neighborhoods for Creoles of color were the Faubourg Marigny and Tremé. If the man was not married, he might keep a separate residence, preferably next door or in the same or next block as his placée. He often took part in and arranged for the upbringing and education of their children. For a time both boys and girls were educated in France, as there were no schools in New Orleans for mixed-race children. As supporting such a plaçage arrangement(s) ran into thousands of dollars per year, it was limited to the wealthy. Inheritance and work Upon the death of her protector, the placée and her family could, on legal challenge, expect up to a third of the man's property. Some white lovers tried, and succeeded, in making their mixed-race children primary heirs over other white descendants or relatives. The women in these relationships often worked to develop assets: acquiring property, running a legitimate rooming-house, or a small business as a hairdresser, marchande (female street or country merchant/vendor), or a seamstress. She could also become a placée to another white Creole. She sometimes taught her daughters to become placées, by education and informal schooling in dress, comportment, and ways to behave. A mother negotiated with a young man for the dowry or property settlement, sometimes by contract, for her daughter if a white Creole were interested in her. A former placée could also marry or to cohabit with a Creole man of color and have more children. Contrary to popular misconceptions, placées were not and did not become prostitutes. Creole men of color objected to the practice as denigrating the virtue of Creole women of color, but some, as descendants of white males, benefited by the transfer of social capital. Martin writes, "They did not choose to live in concubinage; what they chose was to survive." In the late 19th and early 20th centuries, after Reconstruction and with the reassertion of white supremacy across the former Confederacy, the white Creole historians, Charles Gayarré and Alcée Fortier, wrote histories that did not address plaçage in much detail. They suggested that little race mixing had occurred during the colonial period, and that the placées had seduced or led white Creole men astray. They wrote that the French Creoles (in the sense of having long been native to Louisiana) were ethnic Europeans who were threatened by the spectre of race-mixing like other Southern whites. Gayarré, when younger, was said to have taken a woman of color as his placée and she had their children, to his later shame. He married a white woman late in life. His earlier experience inspired his novel, Fernando de Lemos. Marie Thérèse Metoyer Marie Thérèse Metoyer dite Coincoin became an icon of black female entrepreneurship in colonial Louisiana. She was born at the frontier outpost of Natchitoches on Cane River in August 1742 as a slave of the post founder, the controversial explorer Louis Juchereau de St. Denis. She would be, for twenty years, the placée of a French colonial merchant-turned-planter, Claude Thomas Pierre Métoyer two years her junior. At the onset of their plaçage, she was already the mother of five children; she would bear ten more to Métoyer. In 1778, he freed her after the parish priest filed charges against Coincoin as a "public concubine" and threatened to have her sold at New Orleans if they did not end their relationship. As a free woman, she remained with Métoyer until 1788, when his growing fortune persuaded him to take a wife who could provide legal heirs. (He chose another Marie Thérèse, a white Créole of French and German birth.) In setting Coincoin aside, Métoyer donated to her his interest in 80 arpents, about 68 acres (280,000 m2) of unpatented land, adjacent to his plantation, to help support their free-born offspring. On that modest tract, Coincoin planted tobacco, a valuable commodity in the struggling colony. She and her children trapped bears and wild turkeys for sales of meat, hide, and oil locally and at the New Orleans market. She also manufactured medicine, a skill shared by her freed-slave sister Marie Louise dite Mariotte and likely one acquired from their African-born parents. With this money, she progressively bought the freedom of four of her first five children and several grandchildren, before investing in three African-born slaves to provide the physical labor that became more difficult as she aged. After securing a colonial patent on her homestead in 1794, she petitioned for and was given a land concession from the Spanish crown. On that piney-woods tract of 800 arpents (667 ac) on Old Red River, about 5 mi from her farmstead, she set up a vacherie (a ranch) and engaged a Spaniard to tend her cattle. Shortly before her death in 1816, Coincoin sold her homestead and divided her remaining property (her piney-woods land, the three African slaves, and their offspring) among her own progeny. As often happened among the children of plaçages, Coincoin's one surviving daughter by Métoyer, Marie Susanne, became a placée also. As a young woman, apparently with the blessing of both parents, she entered into a relationship with a newly arrived physician, Joseph Conant from New Orleans. When he left Cane River, soon after the birth of their son, she formed a second and lifelong plaçage with a Cane River planter, Jean Baptiste Anty. As a second-generation entrepreneur, Susanne became far more successful than her mother and died in 1838 leaving an estate of $61,600 (equivalent to $1,500,000 in 2009 currency). Modern archaeological work at the site of Coincoin's farmstead is documenting some of the aspects of her domestic life. A mid-nineteenth century dwelling, now dubbed the Coincoin-Prudhomme House although it was not the actual site of her residence, commemorates her within the Cane River National Heritage Area. Popular lore also has, erroneously, credited her with the ownership of a Cane River plantation founded by her son Louis Metoyer, known today as Melrose Plantation, and its historic buildings Yucca House and African House. Her eldest half-French son, Nicolas Augustin Métoyer, founded St. Augustine Parish (Isle Brevelle) Church, the spiritual center of Cane River's large community of Creoles of color who trace their heritage to Coincoin. Eulalie de Mandéville There were many other examples of white Creole fathers who reared and carefully and quietly placed their daughters of color with the sons of known friends or family members. This occurred with Eulalie de Mandéville, the elder half-sister of color to the eccentric nobleman, politician, and land developer Bernard Xavier de Marigny de Mandéville. Taken from her slave mother as a baby, and partly raised by a white grandmother, 22-year-old Eulalie was "placed" by her father, Count Pierre Enguerrand Philippe, Écuyer de Mandéville, Sieur de Marigny, with Eugène de Macarty, a member of the famous French-Irish clan in 1796. Their alliance resulted in five children and lasted almost fifty years. Macarty, like some white Creoles, never married a white woman. (In contrast to the Macartys' stable relationship, Eugène's brother Augustin de Macarty was married and was said to have had numerous, complex affairs with Creole women of color. When he died, several women made claims on behalf of their children against his estate. On his deathbed in 1845, Eugène de Macarty married Eulalie. He willed her all of his money and property, then worth $12,000. His white relatives, including his niece, Marie Delphine de Macarty LaLaurie, contested the will. The court upheld his will. After Eulalie's death, their surviving children defeated another attempt by Macarty's relatives to claim his estate, by then worth more than $150,000. Eulalie de Mandéville de Macarty became a successful marchande and ran a dairy. She died in 1848. Rosette Rochon was born in 1767 in colonial Mobile, the daughter of Pierre Rochon, a shipbuilder from a Québécois family (family name was Rocheron in Québec), and his mulâtresse slave-consort Marianne, who bore him five other children. Once Rosette reached a suitable age, she became the consort of a Monsieur Hardy, with whom she relocated to the colony of Saint Domingue. During her sojourn there, Hardy must have died or relinquished his relationship with her; for in 1797 during the Haitian Revolution, she escaped to New Orleans, where she later became the placée of Joseph Forstal and Charles Populus, both wealthy white New Orleans Creoles. Rochon came to speculate in real estate in the French Quarter; she eventually owned rental property, opened grocery stores, made loans, bought and sold mortgages, and owned and rented out (hired out) slaves. She also traveled extensively back and forth to Haiti, where her son by Hardy had become a government official in the new republic. Her social circle in New Orleans once included Marie Laveau, Jean Lafitte, and the free black contractors and real estate developers Jean-Louis Doliolle and his brother Joseph Doliolle. In particular, Rochon became one of the earliest investors in the Faubourg Marigny, acquiring her first lot from Bernard de Marigny in 1806. Bernard de Marigny, the Creole speculator, refused to sell the lots he was subdividing from his family plantation to anyone who spoke English. While this turned out to be a losing financial decision, Marigny felt more comfortable with the French-speaking, Catholic free people of color (having relatives, lovers, and even children on this side of the color line). Consequently, much of Faubourg Marigny was built by free black artisans for free people of color or for French-speaking white Creoles. Rochon remained largely illiterate, dying in 1863 at the age of 96, leaving behind an estate valued at $100,000 (today, an estate worth a million dollars). Marie Laveau (also spelled Leveau, Laveaux), known as the voodoo queen of New Orleans, was born between 1795 and 1801 as the daughter of a white Haitian plantation owner, Charles Leveaux, and his mixed black and Indian placée Marguerite Darcantel (or D'Arcantel). Because there were so many whites as well as free people of color in Haiti with the same names, Leveaux could also have been a free man of color who owned slaves and property as well. All three may have escaped Haiti along with thousands of other Creole whites and Creoles of color during the slave uprisings that culminated in the French colony's becoming the only independent black republic in the New World. At 17, Marie married a Creole man of color popularly known as Jacques Paris (however, in some documents, he is known as Santiago Paris). Paris either died, disappeared or deliberately abandoned her (some accounts also relate that he was a merchant seaman or sailor in the navy) after she produced a daughter. Laveau was styling herself as the Widow Paris and was a hairdresser for white matrons (she was also reckoned to be an herbalist and yellow fever nurse) when she met Louis-Christophe Dumesnil de Glapion and in the early 1820s, they became lovers. Marie was just beginning her spectacular career as a voodoo practitioner (she would not be declared a 'queen' until about 1830), and Dumesnil de Glapion was a fiftyish white Creole veteran of the Battle of New Orleans with relatives on both sides of the color line. Recently, it's been alleged that Dumesnil de Glapion was so in love with Marie, he refused to live separately from his placée according to racial custom. In an unusual decision, Dumesnil de Glapion passed as a man of color in order to live with her under respectable circumstances—thus explaining the confusion many historians have had whether he was truly white or black. Although it is popularly thought that Marie presented Dumesnil de Glapion with fifteen children, only five are listed in vital statistics and of these, two daughters—one the famous Marie Euchariste or Marie Leveau II—lived to adulthood. Marie Euchariste closely resembled her mother and startled many who thought that Marie Leveau had been resurrected by the black arts, or could be at two places at once, beliefs that the daughter did little to correct. Sebastopol This plantation house and property was built and cultivated by Don Pedro Morin in the 1830s in St. Bernard Parish, Louisiana. It was bought twenty years later by Colonel Ignatius Szymanski a Polish American who later served in the Confederate Army, and renamed Sebastopol. At his death, Colonel Szymanski willed this estate to his placée Eliza Romain, a free woman of color, and to their son John Szymanski. The term quadroon is a fractional one referring to a person with one white and one mulatto parent, some courts would have considered one-fourth Black. The quadroon balls were social events designed to encourage mixed-race women to form liaisons with wealthy white men through a system of concubinage known as plaçage. (Guillory 68-9). Monique Guillory writes about quadroon balls that took place in New Orleans, the city most strongly associated with these events. She approaches the balls in context of the history of a building the structure of which is now the Bourbon Orleans Hotel. Inside is the Orleans Ballroom, a legendary, if not entirely factual, location for the earliest quadroon balls. In 1805, a man named Albert Tessier began renting a dance hall where he threw twice weekly dances for free quadroon women and white men only (80). These dances were elegant and elaborate, designed to appeal to wealthy white men. Although race mixing was prohibited by New Orleans law, it was common for white gentleman to attend the balls, sometimes stealing away from white balls to mingle with the city's quadroon female population. The principal desire of quadroon women attending these balls was to become plaçee as the mistress of a wealthy gentleman, usually a young white Creole or a visiting European (81). These arrangements were a common occurrence, Guillory suggests, because the highly educated, socially refined quadroons were prohibited from marrying white men and were unlikely to find Black men of their own status. A quadroon's mother usually negotiated with an admirer the compensation that would be received for having the woman as his mistress. Typical terms included some financial payment to the parent, financial and/or housing arrangements for the quadroon herself, and, many times, paternal recognition of any children the union produced. Guillory points out that some of these matches were as enduring and exclusive as marriages. A beloved quadroon mistress had the power to destabilize white marriages and families, something she was much resented for. According to Guillory, the system of plaçage had a basis in the economics of mixed race. The plaçage of black women with white lovers, Guillory writes, could take place only because of the socially determined value of their light skin, the same light skin that commanded a higher price on the slave block, where light skinned girls fetched much higher prices than did prime field hands (82). Guillory posits the quadroon balls as the best among severely limited options for these near-white women, a way for them to control their sexuality and decide the price of their own bodies. She contends, "The most a mulatto mother and a quadroon daughter could hope to attain in the rigid confines of the black/white world was some semblance of economic independence and social distinction from the slaves and other blacks" (83). She notes that many participants in the balls were successful in actual businesses when they could no longer rely on an income from the plaçage system. She speculates they developed business acumen from the process of marketing their own bodies. Treatment in fiction - Isabel Allende, Island Beneath the Sea. A novel about a mixed-race slave who is brought to Saint-Domingue and is eventually taken to New Orleans with her master's family. Her quadroon daughter is introduced to society as a placée. - George Washington Cable, The Grandissimes, A Story of Creole Life (1880) by . He also wrote the short stories, "Títe Poulette", "Madame John's Legacy" and "Madame Delphine," which portrayed the placée as societal outcast. - William Faulkner*, Absalom, Absalom!. A young man is engaged to a woman until it is found out that he is already involved with a placée in New Orleans and has a child with her. - Edna Ferber, Saratoga Trunk. The book was later adapted as a film by the same name, starring Ingrid Bergman and Gary Cooper. But it, like the film, falls apart after the action and the heroine move on to Saratoga Springs, New York. - Barbara Hambly, The Benjamin January Mysteries. This series of novels features Benjamin January, a free man of color, in New Orleans in the 1830s. His mother and half-sister are also featured; both are placées. His wife is the daughter of a placée. - Anne Rice, The Feast of All Saints. A coming of age novel about a young man making his way in Creole New Orleans. it was adapted as a film by the same name. - Patricia Vaughn, Shadows on the Bayou., A historical romance following the life of Sylvia Dupont, a young woman raised to be a placée. Dupont marries a free man of color and struggles with the consequences. - Marcus Gardley, The House That Will Not Stand premiered at Berkeley Repertory Theatre, January 31 - March 16, 2014. - Beyoncé Knowles, in her "Formation" music video, features visuals of placées. This song was released February 6, 2016. - Chained to the Rock of Adversity, To Be Free, Black & Female in the Old South, edited by Virginia Meacham Gould, University of Georgia Press, 1998 - Katy F. Morlas, "La Madame et la Mademoiselle," graduate thesis in history, Louisiana State University and Agricultural and Mechanical College, 2003 - Joan M. Martin, Placage and the Louisiana Gens de Couleur Libre, in Creole, edited by Sybil Kein, Louisiana State University Press, Baton Rouge, 2000. - Monique Guillory, "Under One Roof: The Sins and Sanctity of the New Orleans Quadroon Balls," in Race Consciousness, edited by Judith Jackson Fossett and Jeffrey A. Tucker, New York University Press, 1997. - Mills, Elizabeth Shown. "Marie Thérèse Coincoin (1742–1816): Slave, Slave Owner, and Paradox," Chapter 1 in Janet Allred and Judy Gentry, ed., Louisiana Women: Their Lives and Times (Athens, Ga.: University of Georgia Press, 2009), chap. 1, pages 10-29; <Louisiana Women: Their Lives and Times - Google books - Mills, Gary B. The Forgotten People: Cane River's Creoles of Color. Baton Rouge: Louisiana State University Press, 1977. - Mills, Elizabeth Shown. "Which Marie Louise is 'Mariotte'? Sorting Slaves with Common Names." National Genealogical Society Quarterly 94 (September 2006): 183–204; archived online at Historic Pathways . - Morlas, ibid. - Violet Harrington Bryan, "Marcus Christian's Treatment of Les Gens de Couleur Libre," in Creole, edited by Sybil Kein, Louisiana State University Press, Baton Rouge, 2000. - Caryn Cosse Bell, "The Real Marie Laveau," review of Voodoo Queen: The Spirited Lives of Marie Laveau, by Martha Ward, University Press of Mississippi, Jackson, 2004. - Recent books - The Free People of Color of New Orleans, An Introduction, by Mary Gehman and Lloyd Dennis, Margaret Media, Inc., 1994. - Africans in Colonial Louisiana: The Development of Afro-Creole Culture in the Eighteenth Century, by Gwendolyn Midlo Hall, Louisiana State University Press, 1995. - Creole New Orleans, Race and Americanization, by Arnold R. Hirsch and Joseph Logsdon, Louisiana State University Press, 1992. - Bounded Lives, Bounded Places: Free Black Society in Colonial New Orleans, by Kimberly S. Hanger. - Afristocracy: Free Women of Color and the Politics of Race, Class, and Culture, by Angela Johnson-Fisher, Verlag, 2008. - Contemporary accounts - Travels by His Highness Duke Bernhard of Saxe-Weimar-Eisenach through North America in the years 1825 and 1826, by Bernhard, Duke of Saxe-Weimar-Eisenach; William Jeronimus and C.J. Jeronimus, University Press of America, 2001. (The Duke relates his visits to quadroon balls as a tourist in New Orleans.) - Voyage to Louisiana, (An abridged translation from the original French by Stuart O. Landry) by C.C. Robin, Pelican Publishing Co., 1966. (Robin visited Louisiana just after its purchase by the Americans and resided there for two years.) - Mon Cher, Creole genealogical newsletter, dated June 20, 2003, on the genealogy of Marie Laveau, also related to the Trudeaus, page 5. <--broken link, April 2015. - Information about the life of Marie Thérèse Coincoin Metoyer. - History of 918 Barracks Street in the French Quarter, where Eugène Macarty purchased and then built another home for his placée, Eulalie Mandeville (fwc; for free woman of color) and their children. - Website of Louisiana Creoles of color. - Website of the Musée Rosette Rochon, located on 1515 Pauger Street, Marigny, New Orleans. This house, which survived Hurricane Katrina, is the only extant residence built by Mme. Rochon.

https://en.wikipedia.org/wiki/Pla%C3%A7age

Definitions for harlem renaissance This page provides all possible meanings and translations of the word harlem renaissance a period in the 1920s when African-American achievements in art and music and literature flourished The Harlem Renaissance was a cultural movement that spanned the 1920s. At the time, it was known as the "New Negro Movement", named after the 1925 anthology by Alain Locke. Though it was centered in the Harlem neighborhood of New York City, many French-speaking black writers from African and Caribbean colonies who lived in Paris were also influenced by the Harlem Renaissance. The Harlem Renaissance is generally considered to have spanned from about 1919 until the early or mid-1930s. Many of its ideas lived on much longer. The zenith of this "flowering of Negro literature", as James Weldon Johnson preferred to call the Harlem Renaissance, was placed between 1924 and 1929. Sample Sentences & Example Usage The influence of his work by Harlem Renaissance artists is evident. ”Aberjhani is also known as author of Encyclopedia of the Harlem Renaissance, The Bridge of Silver Wings, and The Wisdom of W.E.B. Dubois. He publishes often in various publications, print and online. His poetry has an intensely intimate courage, the sort we would all wish to have, but too often hold protectively back.” “We are drawn to the Harlem Renaissance because of the hope for black uplift and interracial interaction and empathy that it embodied and because there is a certain element of romanticism associated with the era’s creativity, its seemingly larger-than-life heroes and heroines, and its most brilliantly lit terrain, Harlem, USA.” Images & Illustrations of harlem renaissance Find a translation for the harlem renaissance definition in other languages: Select another language: Discuss these harlem renaissance definitions with the community: Word of the Day Would you like us to send you a FREE new word definition delivered to your inbox daily? Use the citation below to add this definition to your bibliography: "harlem renaissance." Definitions.net. STANDS4 LLC, 2016. Web. 13 Feb. 2016. <http://www.definitions.net/definition/harlem renaissance>.

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Advancing Basic Science for Humanity 08/15/2012 - Phoenix Cluster Sets Record Pace at Forming Stars (Originally published by NASA) August 15, 2012 Astronomers have found an extraordinary galaxy cluster -- one of the largest objects in the Universe -- that is breaking several important cosmic records. Observations of this cluster, known as the Phoenix Cluster, with NASA's Chandra X-ray Observatory, the NSF’s South Pole Telescope and eight other world-class observatories, may force astronomers to rethink how these colossal structures, and the galaxies that inhabit them, evolve. Stars are forming in the Phoenix Cluster at the highest rate ever observed for the middle of a galaxy cluster. The object is also the most powerful producer of X-rays of any known cluster, and among the most massive of clusters. The data also suggest that the rate of hot gas cooling in the central regions of the cluster is the largest ever observed. This galaxy cluster has been dubbed the "Phoenix Cluster" because it is located in the constellation of the Phoenix, and because of its remarkable properties. The cluster is located about 5.7 billion light years from Earth. "The mythology of the Phoenix -- a bird rising from the dead -- is a great way to describe this revived object," said Michael McDonald, a Hubble Fellow in the Kavli Institute for Astrophysics and Space Research at the Massachusetts Institute of Technology and the lead author of a paper appearing in the August 16th issue of the journal Nature. "While galaxies at the center of most clusters may have been dormant for billions of years, the central galaxy in this cluster seems to have come back to life with a new burst of star formation." Like other galaxy clusters, Phoenix contains a vast reservoir of hot gas -- containing more normal matter than all of the galaxies in the cluster combined -- that can only be detected with X-ray telescopes like Chandra. The prevailing wisdom had once been that this hot gas should cool over time and sink to the galaxy at the center of the cluster, forming huge numbers of stars. However, most galaxy clusters have formed very few stars over the last few billion years. Astronomers think that the supermassive black hole in the central galaxy of clusters pumps energy into the system, preventing cooling of gas from causing a burst of star formation. The famous Perseus Cluster is an example of a black hole bellowing out energy and preventing the gas from cooling to form stars at a high rate. Repeated outbursts from the black hole in the center of Perseus, in the form of powerful jets, created giant cavities and produced sound waves with an incredibly deep B-flat note 57 octaves below middle C. "We thought that these very deep sounds might be found in galaxy clusters everywhere," said co-author Ryan Foley, a Clay Fellow at the Harvard-Smithsonian Center for Astrophysics in Cambridge, Mass. "The Phoenix Cluster is showing us this is not the case - or at least there are times the music stops. Jets from the giant black hole at the center of a cluster are apparently not powerful enough to prevent the cluster gas from cooling.” With its black hole not producing powerful enough jets, the center of the Phoenix Cluster is buzzing with stars that are forming about 20 times faster than in the Perseus cluster. This rate is the highest seen in the center of a galaxy cluster but not the highest seen anywhere in the Universe. However, the overall record-holding galaxies, located outside clusters, have rates only about twice as high. The frenetic pace of star birth and cooling of gas in Phoenix are causing both the galaxy and the black hole to add mass very quickly -- an important phase that the researchers predict will be relatively short-lived. "The galaxy and its black hole are undergoing unsustainable growth," said co-author Bradford Benson, a Kavli Fellow in the Kavli Institute for Cosmological Physics at the University of Chicago. "This growth spurt can't last longer than about a hundred million years, otherwise the galaxy and black hole would become much bigger than their counterparts in the nearby Universe." Remarkably, the Phoenix Cluster and its central galaxy and supermassive black hole are already among the most massive known objects of their type. Because of their tremendous size, galaxy clusters are crucial objects for studying cosmology and galaxy evolution and so finding one with such extreme properties like the Phoenix Cluster is important. "This spectacular star burst is a very significant discovery because it suggests we have to rethink how the massive galaxies in the centers of clusters grow," said Martin Rees of Cambridge University, who was not involved with the study. "The cooling of hot gas might be a much more important source of stars than previously thought." The Phoenix Cluster was originally detected by the National Science Foundation's South Pole Telescope, and later was observed in optical light by the Gemini Observatory in Chile as well as the Blanco 4-meter and Magellan telescopes, also in Chile. The hot gas and its rate of cooling were estimated from Chandra data. To measure the star formation rate in the Phoenix Cluster, several space-based telescopes were used including NASA's WISE and GALEX, and ESA's Herschel. NASA's Marshall Space Flight Center in Huntsville, Ala., manages the Chandra program for NASA's Science Mission Directorate in Washington. The Smithsonian Astrophysical Observatory controls Chandra's science and flight operations from Cambridge, Mass.

http://www.kavlifoundation.org/kavli-news/NASA-phoenix-cluster-sets-record-pace-forming-stars

Romanesque architecture is an architectural style of medieval Europe characterized by semi-circular arches. There is no consensus for the beginning date of the Romanesque style, with proposals ranging from the 6th to the late 10th century, this later date being the most commonly held. It developed in the 12th century into the Gothic style, marked by pointed arches. Examples of Romanesque architecture can be found across the continent, making it the first pan-European architectural style since Imperial Roman Architecture. The Romanesque style in England is traditionally referred to as Norman architecture. Combining features of ancient Roman and Byzantine buildings and other local traditions, Romanesque architecture is known by its massive quality, thick walls, round arches, sturdy pillars, groin vaults, large towers and decorative arcading. Each building has clearly defined forms, frequently of very regular, symmetrical plan; the overall appearance is one of simplicity when compared with the Gothic buildings that were to follow. The style can be identified right across Europe, despite regional characteristics and different materials. Many castles were built during this period, but they are greatly outnumbered by churches. The most significant are the great abbey churches, many of which are still standing, more or less complete and frequently in use. The enormous quantity of churches built in the Romanesque period was succeeded by the still busier period of Gothic architecture, which partly or entirely rebuilt most Romanesque churches in prosperous areas like England and Portugal. The largest groups of Romanesque survivors are in areas that were less prosperous in subsequent periods, including parts of southern France, northern Spain and rural Italy. Survivals of unfortified Romanesque secular houses and palaces, and the domestic quarters of monasteries are far rarer, but these used and adapted the features found in church buildings, on a domestic scale. - 1 Definition - 2 Scope - 3 History - 4 Characteristics - 5 Ecclesiastical architecture - 5.1 Plan - 5.2 Section - 5.3 Church and cathedral east ends - 5.4 Church and cathedral façades and external decoration - 5.5 Church towers - 5.6 Portals - 5.7 Interiors - 5.8 Other structures - 5.9 Decoration - 5.10 Transitional style and the continued use of Romanesque forms - 6 Romanesque castles, houses and other buildings - 7 Romanesque Revival - 8 Notes - 9 See also - 10 References - 11 Further reading - 12 External links According to the Oxford English Dictionary, the word "Romanesque" means "descended from Roman" and was first used in English to designate what are now called Romance languages (first cited 1715). The French term "romane" was first used in the architectural sense by archaeologist Charles de Gerville in a letter of 18 December 1818 to Auguste Le Prévost to describe what Gerville sees as a debased Roman architecture.[Notes 2] In 1824 Gerville's friend Arcisse de Caumont adopted the label "roman" to describe the "degraded" European architecture from the 5th to the 13th centuries, in his Essai sur l'architecture religieuse du moyen-âge, particulièrement en Normandie, at a time when the actual dates of many of the buildings so described had not been ascertained: The name Roman (esque) we give to this architecture, which should be universal as it is the same everywhere with slight local differences, also has the merit of indicating its origin and is not new since it is used already to describe the language of the same period. Romance language is degenerated Latin language. Romanesque architecture is debased Roman architecture. The first use in a published work is in William Gunn's An Inquiry into the Origin and Influence of Gothic Architecture (London 1819). The word was used by Gunn to describe the style that was identifiably Medieval and prefigured the Gothic, yet maintained the rounded Roman arch and thus appeared to be a continuation of the Roman tradition of building. The term is now used for the more restricted period from the late 10th to 12th centuries. The term "Pre-romanesque" is sometimes applied to architecture in Germany of the Carolingian and Ottonian periods and Visigothic, Mozarab and Asturian constructions between the 8th and the 10th centuries in the Iberian Peninsula while "First Romanesque" is applied to buildings in north of Italy and Spain and parts of France that have Romanesque features but pre-date the influence of the Abbey of Cluny. Portal, Church of Santa Maria, Viu de Llevata, Catalonia, Spain Cloister of the Basilica di San Giovanni in Laterano, Rome Bell tower of Angoulême Cathedral, Charente, SW France Window and Lombard band of the Rotunda of San Tomè, Almenno San Bartolomeo Buildings of every type were constructed in the Romanesque style, with evidence remaining of simple domestic buildings, elegant town houses, grand palaces, commercial premises, civic buildings, castles, city walls, bridges, village churches, abbey churches, abbey complexes and large cathedrals. Of these types of buildings, domestic and commercial buildings are the most rare, with only a handful of survivors in the United Kingdom, several clusters in France, isolated buildings across Europe and by far the largest number, often unidentified and altered over the centuries, in Italy. Many castles exist, the foundations of which date from the Romanesque period. Most have been substantially altered, and many are in ruins. By far the greatest number of surviving Romanesque buildings are churches. These range from tiny chapels to large cathedrals, and although many have been extended and altered in different styles, a large number remain either substantially intact or sympathetically restored, demonstrating the form, character and decoration of Romanesque church architecture. Saint Nicholas Rotunda in Cieszyn, Poland The Civic Hall in Massa Marittima, Italy Abbey Church of St James, Lebeny, Hungary (1208) The keep of Conisbrough Castle, England. Romanesque architecture was the first distinctive style to spread across Europe since the Roman Empire. With the decline of Rome, Roman building methods survived to an extent in Western Europe, where successive Merovingian, Carolingian and Ottonian architects continued to build large stone buildings such as monastery churches and palaces. In the more northern countries Roman building styles and techniques had never been adopted except for official buildings, while in Scandinavia they were unknown. Although the round arch continued in use, the engineering skills required to vault large spaces and build large domes were lost. There was a loss of stylistic continuity, particularly apparent in the decline of the formal vocabulary of the Classical Orders. In Rome several great Constantinian basilicas continued in use as an inspiration to later builders. Some traditions of Roman architecture also survived in Byzantine architecture with the 6th-century octagonal Byzantine Basilica of San Vitale in Ravenna being the inspiration for the greatest building of the Dark Ages in Europe, the Emperor Charlemagne's Palatine Chapel, Aachen, Germany, built around the year AD 800. Dating shortly after the Palatine Chapel is a remarkable 9th-century Swiss manuscript known as the Plan of Saint Gall and showing a very detailed plan of a monastic complex, with all its various monastic buildings and their functions labelled. The largest building is the church, the plan of which is distinctly Germanic, having an apse at both ends, an arrangement not generally seen elsewhere. Another feature of the church is its regular proportion, the square plan of the crossing tower providing a module for the rest of the plan. These features can both be seen at the Proto-Romanesque St. Michael's Church, Hildesheim, 1001–1030. Architecture of a Romanesque style also developed simultaneously in the north of Italy, parts of France and in the Iberian Peninsula in the 10th century and prior to the later influence of the Abbey of Cluny. The style, sometimes called First Romanesque or Lombard Romanesque, is characterised by thick walls, lack of sculpture and the presence of rhythmic ornamental arches known as a Lombard band. Santa Maria in Cosmedin, Rome (8th – early 12th century) has a basilical plan and reuses ancient Roman columns. St. Michael's Church, Hildesheim has similar characteristics to the church in the Plan of Saint Gall. Charlemagne was crowned by the Pope in Old St. Peter's Basilica on Christmas Day in the year 800, with an aim to re-establishing the old Roman Empire. Charlemagne's political successors continued to rule much of Europe, with a gradual emergence of the separate political states that were eventually to become welded into nations, either by allegiance or defeat, the Kingdom of Germany giving rise to the Holy Roman Empire. The invasion of England by William, Duke of Normandy, in 1066, saw the building of both castles and churches that reinforced the Norman presence. Several significant churches that were built at this time were founded by rulers as seats of temporal and religious power, or places of coronation and burial. These include the Abbaye-Saint-Denis, Speyer Cathedral and Westminster Abbey (where little of the Norman church now remains). At a time when the remaining architectural structures of the Roman Empire were falling into decay and much of its learning and technology lost, the building of masonry domes and the carving of decorative architectural details continued unabated, though greatly evolved in style since the fall of Rome, in the enduring Byzantine Empire. The domed churches of Constantinople and Eastern Europe were to greatly affect the architecture of certain towns, particularly through trade and through the Crusades. The most notable single building that demonstrates this is St Mark's Basilica, Venice, but there are many lesser-known examples, particularly in France, such as the church of Saint-Front, Périgueux and Angoulême Cathedral. Much of Europe was affected by feudalism in which peasants held tenure from local rulers over the land that they farmed in exchange for military service. The result of this was that they could be called upon, not only for local and regional spats, but to follow their lord to travel across Europe to the Crusades, if they were required to do so. The Crusades, 1095–1270, brought about a very large movement of people and, with them, ideas and trade skills, particularly those involved in the building of fortifications and the metal working needed for the provision of arms, which was also applied to the fitting and decoration of buildings. The continual movement of people, rulers, nobles, bishops, abbots, craftsmen and peasants, was an important factor in creating a homogeneity in building methods and a recognizable Romanesque style, despite regional differences. Life became generally less secure after the Carolingian period. This resulted in the building of castles at strategic points, many of them being constructed as strongholds of the Normans, descendants of the Vikings who invaded northern France under Rollo in 911. Political struggles also resulted in the fortification of many towns, or the rebuilding and strengthening of walls that remained from the Roman period. One of the most notable surviving fortifications is that of the city of Carcassonne. The enclosure of towns brought about a lack of living space within the walls, and resulted in a style of town house that was tall and narrow, often surrounding communal courtyards, as at San Gimignano in Tuscany. In Germany, the Holy Roman Emperors built a number of residences, fortified, but essentially palaces rather than castles, at strategic points and on trade routes. The Imperial Palace of Goslar (heavily restored in the 19th century) was built in the early 11th century by Otto III and Henry III, while the ruined Palace at Gelnhausen was received by Frederick Barbarossa prior to 1170. The movement of people and armies also brought about the building of bridges, some of which have survived, including the 12th-century bridge at Besalú, Catalonia, the 11th-century Puente de la Reina, Navarre and the Pont-Saint-Bénézet, Avignon. Across Europe, the late 11th and 12th centuries saw an unprecedented growth in the number of churches. A great number of these buildings, both large and small, remain, some almost intact and in others altered almost beyond recognition in later centuries. They include many very well known churches such as Santa Maria in Cosmedin in Rome, the Baptistery in Florence and San Zeno Maggiore in Verona. In France, the famous abbeys of Aux Dames and Les Hommes at Caen and Mont Saint-Michel date from this period, as well as the abbeys of the pilgrimage route to Santiago de Compostela. Many cathedrals owe their foundation to this date, with others beginning as abbey churches, and later becoming cathedrals. In England, of the cathedrals of ancient foundation, all were begun in this period with the exception of Salisbury, where the monks relocated from the Norman church at Old Sarum, and several, such as Canterbury, which were rebuilt on the site of Saxon churches. In Spain, the most famous church of the period is Santiago de Compostela. In Germany, the Rhine and its tributaries were the location of many Romanesque abbeys, notably Mainz, Worms, Speyer and Bamberg. In Cologne, then the largest city north of the Alps, a very important group of large city churches survives largely intact. As monasticism spread across Europe, Romanesque churches sprang up in Scotland, Scandinavia, Poland, Hungary, Sicily, Serbia and Tunisia. Several important Romanesque churches were built in the Crusader kingdoms. The system of monasticism in which the religious become members of an order, with common ties and a common rule, living in a mutually dependent community, rather than as a group of hermits living in proximity but essentially separate, was established by the monk Benedict in the 6th century. The Benedictine monasteries spread from Italy throughout Europe, being always by far the most numerous in England. They were followed by the Cluniac order, the Cistercians, Carthusians and Augustinian Canons. During the Crusades, the military orders of the Knights Hospitaller and the Knights Templar were founded. The monasteries, which sometimes also functioned as cathedrals, and the cathedrals that had bodies of secular clergy often living in community, were a major source of power in Europe. Bishops and the abbots of important monasteries lived and functioned like princes. The monasteries were the major seats of learning of all sorts. Benedict had ordered that all the arts were to be taught and practiced in the monasteries. Within the monasteries books were transcribed by hand, and few people outside the monasteries could read or write. In France, Burgundy was the centre of monasticism. The enormous and powerful monastery at Cluny was to have lasting effect on the layout of other monasteries and the design of their churches. Unfortunately, very little of the abbey church at Cluny remains; the "Cluny II" rebuilding of 963 onwards has completely vanished, but we have a good idea of the design of "Cluny III" from 1088 to 1130, which until the Renaissance remained the largest building in Europe. However, the church of St. Sernin at Toulouse, 1080–1120, has remained intact and demonstrates the regularity of Romanesque design with its modular form, its massive appearance and the repetition of the simple arched window motif. Many parish churches across Europe, such as this in Vestre Slidre, Norway, are of Romanesque foundation Many cathedrals such as Trier Cathedral, Germany, date from this period, with many later additions. Pilgrimage and Crusade One of the effects of the Crusades, which were intended to wrest the Holy Places of Palestine from Islamic control, was to excite a great deal of religious fervour, which in turn inspired great building programs. The Nobility of Europe, upon safe return, thanked God by the building of a new church or the enhancement of an old one. Likewise, those who did not return from the Crusades could be suitably commemorated by their family in a work of stone and mortar. The Crusades resulted in the transfer of, among other things, a great number of Holy Relics of saints and apostles. Many churches, like Saint-Front, Périgueux, had their own home grown saint while others, most notably Santiago de Compostela, claimed the remains and the patronage of a powerful saint, in this case one of the Twelve Apostles. Santiago de Compostela, located near Galicia (present day Spain) became one of the most important pilgrimage destinations in Europe. Most of the pilgrims travelled the Way of St. James on foot, many of them barefooted as a sign of penance. They moved along one of the four main routes that passed through France, congregating for the journey at Jumièges, Paris, Vézelay, Cluny, Arles and St. Gall in Switzerland. They crossed two passes in the Pyrenees and converged into a single stream to traverse north-western Spain. Along the route they were urged on by those pilgrims returning from the journey. On each of the routes abbeys such as those at Moissac, Toulouse, Roncesvalles, Conques, Limoges and Burgos catered for the flow of people and grew wealthy from the passing trade. Saint-Benoît-du-Sault, in the Berry province, is typical of the churches that were founded on the pilgrim route. The general impression given by Romanesque architecture, in both ecclesiastical and secular buildings, is one of massive solidity and strength. In contrast with both the preceding Roman and later Gothic architecture, in which the load-bearing structural members are, or appear to be, columns, pilasters and arches, Romanesque architecture, in common with Byzantine architecture, relies upon its walls, or sections of walls called piers. Romanesque architecture is often divided into two periods known as the "First Romanesque" style and the "Romanesque" style. The difference is chiefly a matter of the expertise with which the buildings were constructed. The First Romanesque employed rubble walls, smaller windows and unvaulted roofs. A greater refinement marks the Second Romanesque, along with increased use of the vault and dressed stone. The walls of Romanesque buildings are often of massive thickness with few and comparatively small openings. They are often double shells, filled with rubble. The building material differs greatly across Europe, depending upon the local stone and building traditions. In Italy, Poland, much of Germany and parts of the Netherlands, brick is generally used. Other areas saw extensive use of limestone, granite and flint. The building stone was often used in comparatively small and irregular pieces, bedded in thick mortar. Smooth ashlar masonry was not a distinguishing feature of the style, particularly in the earlier part of the period, but occurred chiefly where easily worked limestone was available. Because of the massive nature of Romanesque walls, buttresses are not a highly significant feature, as they are in Gothic architecture. Romanesque buttresses are generally of flat square profile and do not project a great deal beyond the wall. In the case of aisled churches, barrel vaults, or half-barrel vaults over the aisles helped to buttress the nave, if it was vaulted. In the cases where half-barrel vaults were used, they effectively became like flying buttresses. Often aisles extended through two storeys, rather than the one usual in Gothic architecture, so as to better support the weight of a vaulted nave. In the case of Durham Cathedral, flying buttresses have been employed, but are hidden inside the triforium gallery. Castle Rising, England, shows flat buttresses and reinforcing at the corners of the building typical in both castles and churches. Abbaye Cerisy le Foret, Normandy, France, has a compact appearance with aisles rising through two storeys buttressing the vault. St Albans Cathedral England, demonstrates the typical alterations made to the fabric of many Romanesque buildings in different styles and materials Arches and openings The arches used in Romanesque architecture are nearly always semicircular, for openings such as doors and windows, for vaults and for arcades. Wide doorways are usually surmounted by a semi-circular arch, except where a door with a lintel is set into a large arched recess and surmounted by a semi-circular "lunette" with decorative carving. These doors sometimes have a carved central jamb. Narrow doors and small windows might be surmounted by a solid stone lintel. Larger openings are nearly always arched. A characteristic feature of Romanesque architecture, both ecclesiastic and domestic, is the pairing of two arched windows or arcade openings, separated by a pillar or colonette and often set within a larger arch. Ocular windows are common in Italy, particularly in the facade gable and are also seen in Germany. Later Romanesque churches may have wheel windows or rose windows with plate tracery. There are a very small number of buildings in the Romanesque style, such as Autun Cathedral in France and Monreale Cathedral in Sicily in which pointed arches have been used extensively, apparently for stylistic reasons. It is believed that in these cases there is a direct imitation of Islamic architecture. At other late Romanesque churches such as Durham Cathedral, and Cefalù Cathedral, the pointed arch was introduced as a structural device in ribbed vaulting. Its increasing application was fundamental to the development of Gothic architecture. An arcade is a row of arches, supported on piers or columns. They occur in the interior of large churches, separating the nave from the aisles, and in large secular interiors spaces, such as the great hall of a castle, supporting the timbers of a roof or upper floor. Arcades also occur in cloisters and atriums, enclosing an open space. Arcades can occur in storeys or stages. While the arcade of a cloister is typically of a single stage, the arcade that divides the nave and aisles in a church is typically of two stages, with a third stage of window openings known as the clerestory rising above them. Arcading on a large scale generally fulfils a structural purpose, but it is also used, generally on a smaller scale, as a decorative feature, both internally and externally where it is frequently "blind arcading" with only a wall or a narrow passage behind it. The facade of Notre Dame du Puy, le Puy en Velay, France, has a more complex arrangement of diversified arches: Doors of varying widths, blind arcading, windows and open arcades. Collegiate Church of Saint Gertrude, Nivelles, Belgium uses fine shafts of Belgian marble to define alternating blind openings and windows. Upper windows are similarly separated into two openings by colonettes. Worms Cathedral, Germany, displays a great variety of openings and arcades including wheel and rose windows, many small simple windows, galleries and Lombard courses. The south portal of the Abbey of Saint-Pierre, Moissac, France, has a square door divided by an ornate doorpost, surmounted by a carved tympanum and set within a vast arched porch. In Romanesque architecture, piers were often employed to support arches. They were built of masonry and square or rectangular in section, generally having a horizontal moulding representing a capital at the springing of the arch. Sometimes piers have vertical shafts attached to them, and may also have horizontal mouldings at the level of the base. Although basically rectangular, piers can often be of highly complex form, with half-segments of large hollow-core columns on the inner surface supporting the arch, or a clustered group of smaller shafts leading into the mouldings of the arch. Piers that occur at the intersection of two large arches, such as those under the crossing of the nave and transept, are commonly cruciform in shape, each arch having its own supporting rectangular pier at right angles to the other. Columns are an important structural feature of Romanesque architecture. Colonnettes and attached shafts are also used structurally and for decoration. Monolithic columns cut from a single piece of stone were frequently used in Italy, as they had been in Roman and Early Christian architecture. They were also used, particularly in Germany, when they alternated between more massive piers. Arcades of columns cut from single pieces are also common in structures that do not bear massive weights of masonry, such as cloisters, where they are sometimes paired. In Italy, during this period, a great number of antique Roman columns were salvaged and reused in the interiors and on the porticos of churches. The most durable of these columns are of marble and have the stone horizontally bedded. The majority are vertically bedded and are sometimes of a variety of colours. They may have retained their original Roman capitals, generally of the Corinthian or Roman Composite style. Some buildings, like Santa Maria in Cosmedin (illustrated above) and the atrium at San Clemente in Rome, may have an odd assortment of columns in which large capitals are placed on short columns and small capitals are placed on taller columns to even the height. Architectural compromises of this type are seen where materials have been salvaged from a number of buildings. Salvaged columns were also used to a lesser extent in France. In most parts of Europe, Romanesque columns were massive, as they supported thick upper walls with small windows, and sometimes heavy vaults. The most common method of construction was to build them out of stone cylinders called drums, as in the crypt at Speyer Cathedral. Hollow core columns Where really massive columns were called for, such as those at Durham Cathedral, they were constructed of ashlar masonry and the hollow core was filled with rubble. These huge untapered columns are sometimes ornamented with incised decorations. A common characteristic of Romanesque buildings, occurring both in churches and in the arcades that separate large interior spaces of castles, is the alternation of piers and columns. The most simple form that this takes is to have a column between each adjoining pier. Sometimes the columns are in multiples of two or three. At St. Michael's, Hildesheim, an A B B A alternation occurs in the nave while an A B A alternation can be seen in the transepts. At Jumièges there are tall drum columns between piers each of which has a half-column supporting the arch. There are many variations on this theme, most notably at Durham Cathedral where the mouldings and shafts of the piers are of exceptional richness and the huge masonry columns are deeply incised with geometric patterns. Often the arrangement was made more complex by the complexity of the piers themselves, so that it was not piers and columns that alternated, but rather, piers of entirely different form from each other, such as those of Sant' Ambrogio, Milan, where the nature of the vault dictated that the alternate piers bore a great deal more weight than the intermediate ones and are thus very much larger. Mainz Cathedral, Germany, has rectangular piers and possibly the earliest example of an internal elevation of 3 stages. (Gothic vault) Malmesbury Abbey, England, has hollow core columns, probably filled with rubble. (Gothic vault) The cathedral of Santiago de Compostela, Spain, has large drum columns with attached shafts supporting a barrel vault. Durham Cathedral, England, has decorated masonry columns alternating with piers of clustered shafts supporting the earliest pointed high ribs. The foliate Corinthian style provided the inspiration for many Romanesque capitals, and the accuracy with which they were carved depended very much on the availability of original models, those in Italian churches such as Pisa Cathedral or church of Sant'Alessandro in Lucca and southern France being much closer to the Classical than those in England. The Corinthian capital is essentially round at the bottom where it sits on a circular column and square at the top, where it supports the wall or arch. This form of capital was maintained in the general proportions and outline of the Romanesque capital. This was achieved most simply by cutting a rectangular cube and taking the four lower corners off at an angle so that the block was square at the top, but octagonal at the bottom, as can be seen at St. Michael's Hildesheim. This shape lent itself to a wide variety of superficial treatments, sometimes foliate in imitation of the source, but often figurative. In Northern Europe the foliate capitals generally bear far more resemblance to the intricacies of manuscript illumination than to Classical sources. In parts of France and Italy there are strong links to the pierced capitals of Byzantine architecture. It is in the figurative capitals that the greatest originality is shown. While some are dependent on manuscripts illustrations of Biblical scenes and depictions of beasts and monsters, others are lively scenes of the legends of local saints. The capitals, while retaining the form of a square top and a round bottom, were often compressed into little more than a bulging cushion-shape. This is particularly the case on large masonry columns, or on large columns that alternate with piers as at Durham.(See illustrated above) Capital of Corinthian form with anthropomorphised details, Pisa Campanile Capital of Corinthian form with Byzantine decoration and carved dosseret, San Martín de Tours, Palencia Capital of amorphous form surmounting a cluster of shafts. The figurative carving shows a winged devil directing Herod to slaughter the Innocents. Monastery of San Juan de Duero, Soria Vaults and roofs The majority of buildings have wooden roofs, generally of a simple truss, tie beam or king post form. In the case of trussed rafter roofs, they are sometimes lined with wooden ceilings in three sections like those that survive at Ely and Peterborough cathedrals in England. In churches, typically the aisles are vaulted, but the nave is roofed with timber, as is the case at both Peterborough and Ely. In Italy where open wooden roofs are common, and tie beams frequently occur in conjunction with vaults, the timbers have often been decorated as at San Miniato al Monte, Florence. Vaults of stone or brick took on several different forms and showed marked development during the period, evolving into the pointed ribbed arch characteristic of Gothic architecture. The simplest type of vaulted roof is the barrel vault in which a single arched surface extends from wall to wall, the length of the space to be vaulted, for example, the nave of a church. An important example, which retains Medieval paintings, is the vault of Saint-Savin-sur-Gartempe, France, of the early 12th century. However, the barrel vault generally required the support of solid walls, or walls in which the windows were very small. Groin vaults occur in early Romanesque buildings, notably at Speyer Cathedral where the high vault of about 1060 is the first employment in Romanesque architecture of this type of vault for a wide nave. In later buildings employing ribbed vaultings, groin vaults are most frequently used for the less visible and smaller vaults, particularly in crypts and aisles. A groin vault is almost always square in plan and is constructed of two barrel vaults intersecting at right angles. Unlike a ribbed vault, the entire arch is a structural member. Groin vaults are frequently separated by transverse arched ribs of low profile as at Speyer and Santiago de Compostela. At Sainte Marie Madeleine, Vézelay, the ribs are square in section, strongly projecting and polychrome. Ribbed vaults came into general use in the 12th century. In ribbed vaults, not only are there ribs spanning the vaulted area transversely, but each vaulted bay has diagonal ribs, following the same course as the groins in a groin vault. However, whereas in a groin vault, the vault itself is the structural member, in a ribbed vault, it is the ribs that are the structural members, and the spaces between them can be filled with lighter, non-structural material. Because Romanesque arches are nearly always semi-circular, the structural and design problem inherent in the ribbed vault is that the diagonal span is larger and therefore higher than the transverse span. The Romanesque builders used a number of solutions to this problem. One was to have the centre point where the diagonal ribs met as the highest point, with the infill of all the surfaces sloping upwards towards it, in a domical manner. This solution was employed in Italy at San Michele, Pavia, and Sant' Ambrogio, Milan. The solution employed in England was to stilt the transverse ribs, maintaining a horizontal central line to the roof like that of a barrel vault. The diagonal ribs could also be depressed, a solution used on the sexpartite vaults at both the Saint-Étienne, (Abbaye-aux-Hommes) and Sainte-Trinité, (Abbaye-les-Dames) at Caen, France, in the late 11th and early 12th centuries. Pointed arched vault The problems encountered in the structure and appearance of vaults was solved late in the Romanesque period with the introduction of pointed arched ribs which allowed the height of both diagonal and transverse ribs to be varied in proportion to each other. Pointed ribs made their first appearance in the transverse ribs of the vaults at Durham Cathedral in northern England, dating from 1128. Durham is a cathedral of massive Romanesque proportions and appearance, yet its builders introduced several structural features that were new to architectural design and were later to be hallmark features of the Gothic. Another Gothic structural feature employed at Durham is the flying buttress. However, these are hidden beneath the roofs of the aisles. The earliest pointed vault in France is that of the narthex of La Madeleine, Vézelay, dating from 1130. They were subsequently employed with the development of the Gothic style at the east end of the Basilica of St Denis in Paris in 1140. An early ribbed vault in the Romanesque architecture of Sicily is that of the chancel at the Cathedral of Cefalù. Domes in Romanesque architecture are generally found within crossing towers at the intersection of a church's nave and transept, which conceal the domes externally. Called a tiburio, this tower-like structure often has a blind arcade near the roof. Romanesque domes are typically octagonal in plan and use corner squinches to translate a square bay into a suitable octagonal base. Octagonal cloister vaults appear "in connection with basilicas almost throughout Europe" between 1050 and 1100. The precise form differs from region to region. The painted barrel vault at the Abbey Church of Saint-Savin-sur-Gartempe is supported on tall marbled columns. The nave of Lisbon Cathedral is covered by a ribbed barrel vaul and has an upper, arched gallery (triforium). The Church of St Philibert, Tournus, has a series of transverse barrel vaults supported on arches. The aisle of the Abbey Church at Mozac has a groin vault supported on transverse arches. The aisles at Peterborough Cathedral have quadripartite ribbed vaults. (The nave has an ancient painted wooden ceiling.) The ribbed vaults at Saint-Étienne, Caen, are sexpartite and span two bays of the nave. The crossing of Speyer Cathedral, Germany, has a dome on squinches. Many parish churches, abbey churches and cathedrals are in the Romanesque style, or were originally built in the Romanesque style and have subsequently undergone changes. The simplest Romanesque churches are aisless halls with a projecting apse at the chancel end, or sometimes, particularly in England, a projecting rectangular chancel with a chancel arch that might be decorated with mouldings. More ambitious churches have aisles separated from the nave by arcades. Abbey and cathedral churches generally follow the Latin Cross plan. In England, the extension eastward may be long, while in Italy it is often short or non-existent, the church being of T plan, sometimes with apses on the transept ends as well as to the east. In France the church of St Front, Périgueux, appears to have been modelled on St. Mark's Basilica, Venice, or the Byzantine Church of the Holy Apostles and is of a Greek cross plan with five domes. In the same region, Angoulême Cathedral is an aisless church of the Latin cross plan, more usual in France, but is also roofed with domes. In Germany, Romanesque churches are often of distinctive form, having apses at both east and west ends, the main entrance being central to one side. It is probable that this form came about to accommodate a baptistery at the west end. NOTE: The plans below do not show the buildings in their current states. The plan of the Abbey of St Gall, Switzerland Germany, Speyer Cathedral France, Autun Cathedral France, Angoulême Cathedral The Abbey Church of St. Gall, Switzerland, shows the plan that was to become common throughout Germanic Europe. It is a Latin Cross with a comparatively long nave and short transepts and eastern end, which is apsidal. The nave is aisled, but the chancel and transepts are not. It has an apsidal west end, which was to become a feature of Churches of Germany, such as Worms Cathedral. Speyer Cathedral, Germany, also has aisless transept and chancel. It has a markedly modular look. A typical Germanic characteristic is the presence of towers framing the chancel and the west end. There is marked emphasis on the western entrance, called Westwerk, which is seen in several other churches. Each vault compartment covers two narrow bays of the nave At Autun Cathedral, France, the pattern of the nave bays and aisles extends beyond the crossing and into the chancel, each aisle terminating in an apse. Each nave bay is separated at the vault by a transverse rib. Each transept projects to the width of two nave bays. The entrance has a narthex which screens the main portal. This type of entrance was to be elaborated in the Gothic period on the transepts at Chartres. Angoulême Cathedral, France, is one of several instances in which the Byzantine churches of Constantinople seem to have been influential in the design in which the main spaces are roofed by domes. This structure has necessitated the use of very thick walls, and massive piers from which the domes spring. There are radiating chapels around the apse, which is a typically French feature and was to evolve into the chevet. As was typically the case in England, Ely Cathedral was a Benedictine monastery, serving both monastic and secular function. To facilitate this, the chancel or "presbytery" is longer than usually found in Europe, as are the aisled transepts which contained chapels. In England, emphasis was placed on the orientation of the chapels to the east. The very large piers at the crossing signify that there was once a tower. The western end having two round towers flanking a tall central tower was unique in Britain. Ely Cathedral was never vaulted and retains a wooden ceiling over the nave. The cathedral of Santiago de Compostela shares many features with Ely, but is typically Spanish in its expansive appearance. Santiago held the body of St. James and was the most significant pilgrimage site in Europe. The narthex, the aisles, the large aisled transepts and numerous projecting chapels reflect this. The chancel is short, compared to that of Ely, and the altar set so as to provide clear view to a vast congregation simultaneously. In section, the typical aisled church or cathedral has a nave with a single aisle on either side. The nave and aisles are separated by an arcade carried on piers or on columns. The roof of the aisle and the outer walls help to buttress the upper walls and vault of the nave, if present. Above the aisle roof are a row of windows known as the clerestory, which give light to the nave. During the Romanesque period there was a development from this two-stage elevation to a three-stage elevation in which there is a gallery, known as a triforium, between the arcade and the clerestory. This varies from a simple blind arcade decorating the walls, to a narrow arcaded passage, to a fully developed second story with a row of windows lighting the gallery. This drawing is a reconstruction by Dehio of the appearance of the Romanesque Konstanz Cathedral before its alterations in the Gothic style. It has a typical elevation of nave and aisles with wooden panelled ceilings and an apsidal east end. Exterior elevation, Peterborough Cathedral Church and cathedral east ends The eastern end of a Romanesque church is almost always semi-circular, with either a high chancel surrounded by an ambulatory as in France, or a square end from which an apse projects as in Germany and Italy. Where square ends exist in English churches, they are probably influenced by Anglo Saxon churches. Peterborough and Norwich Cathedrals have retained round east ends in the French style. However, in France, simple churches without apses and with no decorative features were built by the Cistercians who also founded many houses in England, frequently in remote areas. The Cathedral of Santa Maria d'Urgell, Spain, has an apsidal east end projecting at a lower level to the choir and decorated with an arcade below the roofline. This form is usual in Italy and Germany. The Abbey of Sant'Antimo has a high apsidal end surrounded by an ambulatory and with small projecting apses Church and cathedral façades and external decoration Romanesque church facades, generally to the west end of the building, are usually symmetrical, have a large central portal made significant by its mouldings or porch, and an arrangement of arched-topped windows. In Italy there is often a single central ocular or wheel window. The common decorative feature is arcading. Smaller churches often have a single tower that is usually placed to the western end in France or England, either centrally or to one side, while larger churches and cathedrals often have two. In France, Saint-Étienne, Caen, presents the model of a large French Romanesque facade. It is a symmetrical arrangement of nave flanked by two tall towers each with two buttresses of low flat profile that divide the facade into three vertical units. The lowest stage is marked by large doors, each set within an arch in each of the three vertical sections. The wider central section has two tiers of three identical windows, while in the outer sections there are two tiers of single windows, giving emphasis to the mass of the towers. The towers rise above the facade through three further tiers, the lowest of tall blind arcading, the next of arcading pierced by two narrow windows and the third of two large windows, divided into two lights by a colonnette. This facade can be seen as the foundation for many other buildings, including both French and English Gothic churches. While the form is typical of northern France, its various components were common to many Romanesque churches of the period across Europe. Similar facades are found in Portugal. In England, Southwell Cathedral has maintained this form, despite the insertion of a huge Gothic window between the towers. Lincoln and Durham must once have looked like this. In Germany, Limburg Cathedral has a rich variety of openings and arcades in horizontal storeys of varying heights. The churches of San Zeno Maggiore, Verona, and San Michele, Pavia, present two types of facade that are typical of Italian Romanesque, that which reveals the architectural form of the building, and that which screens it. At San Zeno, the components of nave and aisles are made clear by the vertical shafts that rise to the level of the central gable and by the varying roof levels. At San Miniato al Monte the definition of the architectural parts is made even clearer by the polychrome marble, a feature of many Italian Medieval facades, particularly in Tuscany. At San Michele the vertical definition is present as at San Zeno, but the rooflines are screened behind a single large gable decorated with stepped arcading. At Santa Maria della Pieve, Arezzo, this screening is carried even further, as the roofline is horizontal and the arcading rises in many different levels while the colonettes that support them have a great diversity of decoration. In the Rhineland and Netherlands the Carolingian form of west end known as the westwerk prevailed. Towers and apse of the western end are often incorporated into a multi-storey structure that bears little structural or visual relationship to the building behind it. These westwerks take a great variety of forms as may be seen at Maria Laach Abbey, St Gertrude, Nivelles, and St Serviatius, Maastricht. The Old Cathedral of Coimbra, Portugal, is fortress-like and battlemented. The two central openings are deeply recessed. Pisa Cathedral, Italy. The entire building is faced with marble striped in white and grey. On the facade this pattern is overlaid with architectonic decoration of blind arcading below tiers of dwarf galleries. The three portals became increasingly common. Angoulême Cathedral, France. The facade here, richly decorated with architectonic and sculptural forms, has much in common with that at Empoli in that it screens the form of the building behind it. Saint-Étienne, Abbaye aux Hommes, Caen, France, 11th century, with its tall towers, three portals and neat definition of architectural forms became a model for the facades of many later cathedrals across Europe. 14th-century spires Southwell Cathedral, England, 1120, follows the Norman model with pyramidal spires as were probably at Saint-Étienne. The Perpendicular window and battlement are late Gothic. Lisbon Cathedral, Portugal, 1147, has a similar form to the Old Cathedral of Coimbra above with the addition of two sturdy bell towers in the Norman manner and a wheel window. Parma Cathedral, Italy, 1178, has a screen facade ornamented with galleries. At the centre is an open porch surmounted by a ceremonial balcony. The tower, (Gothic 1284) is a separate structure as usual in Italy. Towers were an important feature of Romanesque churches and a great number of them are still standing. They take a variety of forms: square, circular and octagonal, and are positioned differently in relation to the church building in different countries. In northern France, two large towers, such as those at Caen, were to become an integral part of the facade of any large abbey or cathedral. In central and southern France this is more variable and large churches may have one tower or a central tower. Large churches of Spain and Portugal usually have two towers. Many abbeys of France, such as that at Cluny, had many towers of varied forms. This is also common in Germany, where the apses were sometimes framed with circular towers and the crossing surmounted by an octagonal tower as at Worms Cathedral. Large paired towers of square plan could also occur on the transept ends, such as those at Tournai Cathedral in Belgium. In Germany, where four towers frequently occur, they often have spires that may be four or eight sided, or the distinctive Rhenish helm shape seen on the cathedrals of Limburg or Speyer. It is also common to see bell or onion-shaped spires of the Baroque period surmounting Romanesque towers in central and Eastern Europe. In England, for large abbeys and cathedral buildings, three towers were favoured, with the central tower being the tallest. This was often not achieved, through the slow process of the building stages, and in many cases the upper parts of the tower were not completed until centuries later as at Durham and Lincoln. Large Norman towers exist at the cathedrals of Durham, Exeter, Southwell, Norwich and Tewkesbury Abbey. Such towers were often topped during the late Medieval period with a Gothic spire of wooden construction covered with lead, copper or shingles. In the case of Norwich Cathedral, the huge, ornate, 12th-century crossing-tower received a 15th-century masonry spire rising to a height of 320 feet and remaining to this day. In Italy towers are almost always free standing and the position is often dictated by the landform of the site, rather than aesthetics. This is the case in nearly all Italian churches both large and small, except in Sicily where a number of churches were founded by the Norman rulers and are more French in appearance. As a general rule, large Romanesque towers are square with corner buttresses of low profile, rising without diminishing through the various stages. Towers are usually marked into clearly defined stages by horizontal courses. As the towers rise, the number and size of openings increases as can be seen on the right tower of the transept of Tournai Cathedral where two narrow slits in the fourth level from the top becomes a single window, then two windows, then three windows at the uppermost level. This sort of arrangement is particularly noticeable on the towers of Italian churches, which are usually built of brick and may have no other ornament. Two fine examples occur at Lucca, at the church of San Frediano and at the Duomo. It is also seen in Spain. In Italy there are a number of large free-standing towers that are circular, the most famous of these being the Leaning Tower of Pisa. In other countries where circular towers occur, such as Germany, they are usually paired and often flank an apse. Circular towers are uncommon in England, but occur throughout the Early Medieval period in Ireland. The most massive Romanesque crossing tower is that at Tewkesbury Abbey, in England, where large crossing towers are characteristic. (See pic. St Alban's Cathedral, above) The Leaning Tower of Pisa with its encircling arcades is the best known (and most richly decorated) of the many circular towers found in Italy. Romanesque churches generally have a single portal centrally placed on the west front, the focus of decoration for the facade of the building. Some churches such as Saint-Étienne, Caen, (11th century) and Pisa Cathedral (late 12th century) had three western portals, in the manner of Early Christian basilicas. Many churches, both large and small, had lateral entrances that were commonly used by worshippers. Romanesque doorways have a character form, with the jambs having a series of receding planes, into each of which is set a circular shaft, all surmounted by a continuous abacus. The semi-circular arch which rises from the abacus has the same seried planes and circular mouldings as the jambs. There are typically four planes containing three shafts, but there may be as many as twelve shafts, symbolic of the apostles. The opening of the portal may be arched, or may be set with a lintel supporting a tympanum, generally carved, but in Italy sometimes decorated with mosaic or fresco. A carved tympanum generally constitutes the major sculptural work of a Romanesque church. The subject of the carving on a major portal may be Christ in Majesty or the Last Judgement. Lateral doors may include other subjects such as the Birth of Christ. The portal may be protected by a porch, with simple open porches being typical of Italy, and more elaborate structures typical of France and Spain. The mouldings of the arched central west door of Lincoln Cathedral are decorated by chevrons and other formal and figurative ornament typical of English Norman. The "Gallery of Kings" above the portal is Gothic The Porta Platerias, Cathedral of Santiago de Compostela, by Master Esteban, has two wide openings with tympanums supported on brackets. The sculptured frieze above is protected by an eave on corbels. The structure of large churches differed regionally and developed across the centuries. The use of piers of rectangular plan to support arcades was common, as at Mainz Cathedral and St Gertrude Nivelle, and remained usual in smaller churches across Europe, with the arcades often taking the form of openings through the surface of a wall. In Italy, where there was a strong tradition of using marble columns, complete with capital, base and abacus, this remained prevalent, often reusing existent ancient columns, as at San Miniato al Monte. A number of 11th-century churches have naves distinguished by huge circular columns with no clerestory, or a very small one as at St Philibert, Tournus. In England stout columns of large diameter supported decorated arches, gallery and clerestory, as at the nave of Malmesbury Abbey (see "Piers and columns", above). By the early 12th century composite piers had evolved, in which the attached shafts swept upward to a ribbed vault or were continued into the mouldings of the arcade, as at Vézelay Abbey, Saint-Étienne, Caen, and Peterborough Cathedral. The nature of the internal roofing varied greatly, from open timber roofs, and wooden ceilings of different types, which remained common in smaller churches, to simple barrel vaults and groin vaults and increasingly to the use of ribbed vaults in the late 11th and 12th centuries, which were to become a common feature of larger abbey churches and cathedrals. A number of Romanesque churches are roofed with a series of Domes. At Fontevrault Abbey the nave is covered by four domes, while at the Church of Saint Front, Périgueux, the church is of Greek cross plan, with a central dome surrounded by four smaller domes over the nave, chancel and transepts. Internal decoration varied across Europe. Where wide expanses of wall existed, they were often plastered and painted. Wooden ceilings and timber beams were decorated. In Italy walls were sometimes faced with polychrome marble. Where buildings were constructed of stone that was suitable for carving, many decorative details occur, including ornate capitals and mouldings. The apsidal east end was often a focus of decoration, with both architectonic forms such as arcading and pictorial features such as carved figures, murals and occasionally mosaics. Stained glass came into increasing use from the 11th century. In many churches the eastern end has been rebuilt in a later style. Of England's Norman cathedrals, no eastern end remains unchanged. In France the eastern terminals of the important abbeys of Caen, Vézelay and, most significantly, the Basilica of St Denis were completely rebuilt in the Gothic style. In Germany, major reconstructions of the 19th century sought to return many Romanesque buildings to their original form. Examples of simple Romanesque apses can be seen in the images of St Gertrude, Nivelles; St Philibert, Tournus, and San Miniato al Monte. The Church of St Philibert, Tournus, (990-1019) has tall circular piers supporting the arcade and is roofed with a series of barrel vaults supported on arches. Small clerestory windows light the vault. Abbey of St Mary Magdalene, Vézelay, (consecrated 1104) has clusters of vertical shafts rising to support transverse arches and a groin vault. The dressed polychrome stonework has exquisitely detailed mouldings. East end is Gothic. The nave of Peterborough Cathedral (1118–93) in three stages of arcade, gallery & clerestory, typical of Norman abbey churches. The rare wooden ceiling retains its original decoration (c. 1230). Gothic arches beneath tower (c. 1350). Among the structures associated with church buildings are crypts, porches, chapter houses, cloisters and baptisteries. Crypts are often present as an underlying structure to a substantial church, and are generally a completely discrete space, but occasionally, as in some Italian churches, may be a sunken space under a raised chancel and open, via steps, to the body of the nave. Romanesque crypts have survived in many instances, such as Canterbury Cathedral, when the church itself has been rebuilt. The usual construction of a Romanesque crypt is with many short stout columns carrying groin vaults, as at Worcester Cathedral. Porches sometimes occur as part of the original design of a facade. This is very much the case in Italy, where they are usually only one bay deep and are supported on two columns, often resting on couchant lions, as at St Zeno, Verona.See above. Elsewhere, porches of various dates have been added to the facade or side entrance of existent churches and may be quite a substantial structure, with several bays of vaulting supported on an open or partially open arcade, and forming a sort of narthex as at the Church of St Maria, Laach.See above In Spain, Romanesque churches often have large lateral porches, like loggias. Chapter houses often occur adjacent to monastic or cathedral churches. Few have survived intact from the Romanesque period. Early chapter houses were rectangular in shape, with the larger ones sometimes having groin or ribbed vaults supported on columns. Later Romanesque chapter houses sometimes had an apsidal eastern end. The chapter house at Durham Cathedral is a wide space with a ribbed vault, restored as originally constructed in 1130. The circular chapter house at Worcester Cathedral, built by Bishop Wulfstan (1062–95), was the first circular chapter house in Europe and was much imitated in England. Cloisters are generally part of any monastic complex and also occur at cathedral and collegiate churches. They were essential to the communal way of life, a place for both working during daylight hours and relaxing during inclement weather. They usually abut the church building and are enclosed with windowless walls on the outside and an open arcade on the inside, looking over a courtyard or "cloister garth". They may be vaulted or have timber roofs. The arcades are often richly decorated and are home to some of the most fanciful carved capitals of the Romanesque period with those of Santo Domingo de Silos in Spain and the Abbey of St Pierre Moissac, being examples. Many Romanesque cloisters have survived in Spain, France, Italy and Germany, along with some of their associated buildings. Baptisteries often occur in Italy as a free standing structure, associated with a cathedral. They are generally octagonal or circular and domed. The interior may be arcaded on several levels as at Pisa Cathedral. Other notable Romanesque baptisteries are that at Parma Cathedral remarkable for its galleried exterior, and the polychrome Baptistery of San Giovanni of Florence Cathedral, with vault mosaics of the 13th century including Christ in Majesty, possibly the work of the almost legendary Coppo di Marcovaldo. The groin-vaulted crypt of Worcester Cathedral The cloister of Lavaudieu Abbey The Baptistery of Parma Cathedral Arcading is the single most significant decorative feature of Romanesque architecture. It occurs in a variety of forms, from the Lombard band, which is a row of small arches that appear to support a roofline or course, to shallow blind arcading that is often a feature of English architecture and is seen in great variety at Ely Cathedral, to the open dwarf gallery, first used at Speyer Cathedral and widely adopted in Italy as seen on both Pisa Cathedral and its famous Leaning Tower. Arcades could be used to great effect, both externally and internally, as exemplified by the church of Santa Maria della Pieve, in Arezzo. Overlapping arches form a blind arcade at St Lawrence's church Castle Rising, England. (1150) The semi-circular arches form pointed arches where they overlap, a motif which may have influenced Gothic. Dwarf galleries are a major decorative feature on the exterior of Speyer Cathedral, Germany (1090–1106), surrounding the walls and encircling the towers. This was to become a feature of Rhenish Romanesque. The eastern apse of Parma Cathedral, Italy (early 12th century) combines a diversity of decorative features: blind arcading, galleries, courses and sculptured motifs. The arcading on the facade of Lucca Cathedral, Tuscany (1204) has many variations in its decorative details, both sculptural and in the inlaid polychrome marble. Polychrome blind arcading of the apse of Monreale Cathedral, Sicily (1174–82) The decoration indicates Islamic influence in both the motifs and the fact that all the arches, including those of the windows, are pointed. The Romanesque period produced a profusion of sculptural ornamentation. This most frequently took a purely geometric form and was particularly applied to mouldings, both straight courses and the curved moldings of arches. In La Madeleine, Vezelay, for example, the polychrome ribs of the vault are all edged with narrow filets of pierced stone. Similar decoration occurs around the arches of the nave and along the horizontal course separating arcade and clerestory. Combined with the pierced carving of the capitals, this gives a delicacy and refinement to the interior. In England, such decoration could be discrete, as at Hereford and Peterborough cathedrals, or have a sense of massive energy as at Durham where the diagonal ribs of the vaults are all outlined with chevrons, the mouldings of the nave arcade are carved with several layers of the same and the huge columns are deeply incised with a variety of geometric patterns creating an impression of directional movement. These features combine to create one of the richest and most dynamic interiors of the Romanesque period. Although much sculptural ornament was sometimes applied to the interiors of churches, the focus of such decoration was generally the west front, and in particular, the portals. Chevrons and other geometric ornaments, referred to by 19th-century writers as "barbaric ornament", are most frequently found on the mouldings of the central door. Stylized foliage often appears, sometimes deeply carved and curling outward after the manner of the acanthus leaves on Corinthian capitals, but also carved in shallow relief and spiral patterns, imitating the intricacies of manuscript illuminations. In general, the style of ornament was more classical in Italy, such as that seen around the door of San Giusto in Lucca, and more "barbaric" in England, Germany and Scandinavia, such as that seen at Lincoln and Speyer Cathedrals. France produced a great range of ornament, with particularly fine interwoven and spiralling vines in the "manuscript" style occurring at Saint-Sernin, Toulouse. The portal of the Hermitage of St Segundo, Avila, has paired creatures. and decorative bands of floral and interlacing. The pairing of creatures could draw on Byzantine and Celtic models. On these mouldings around the portal of Lincoln Cathedral are formal chevron ornament, tongue-poking monsters, vines and figures, and symmetrical motifs. St Martin's Church, Gensac-la-Pallue has capitals with elaborate interlacing. With the fall of the Roman Empire, the tradition of carving large works in stone and sculpting figures in bronze died out. The best-known surviving large sculptural work of Proto-Romanesque Europe is the life-size wooden Crucifix commissioned by Archbishop Gero of Cologne in about 960–65. During the 11th and 12th centuries, figurative sculpture flourished in a distinctly Romanesque style that can be recognised across Europe, although the most spectacular sculptural projects are concentrated in South-Western France, Northern Spain and Italy. Major figurative decoration occurs particularly around the portals of cathedrals and churches, ornamenting the tympanum, lintels, jambs and central posts. The tympanum is typically decorated with the imagery of Christ in Majesty with the symbols of the Four Evangelists, drawn directly from the gilt covers of medieval Gospel Books. This style of doorway occurs in many places and continued into the Gothic period. A rare survival in England is that of the "Prior's Door" at Ely Cathedral. In France, many have survived, with impressive examples at the Abbey of Saint-Pierre, Moissac, the Abbey of Sainte-Marie, Souillac, and Abbey of la Madaleine, Vézelay – all daughter houses of Cluny, with extensive other sculpture remaining in cloisters and other buildings. Nearby, Autun Cathedral has a Last Judgement of great rarity in that it has uniquely been signed by its creator Giselbertus (who was perhaps the patron rather than the sculptor). The same artist is thought to have worked at la Madaleine Vezelay which uniquely has two elaborately carved tympanum, the early inner one representing the Last Judgement and that on the outer portal of the narthex representing Jesus sending forth the Apostles to preach to the nations. It is a feature of Romanesque art, both in manuscript illumination and sculptural decoration, that figures are contorted to fit the space that they occupy. Among the many examples that exist, one of the finest is the figure of the Prophet Jeremiah from the pillar of the portal of the Abbey of Saint-Pierre, Moissac, France, from about 1130. A significant motif of Romanesque design is the spiral, a form applied to both plant motifs and drapery in Romanesque sculpture. An outstanding example of its use in drapery is that of the central figure of Christ on the outer portal at La Madaleine, Vezelay. Many of the smaller sculptural works, particularly capitals, are Biblical in subject and include scenes of Creation and the Fall of Man, episodes from the life of Christ and those Old Testament scenes that prefigure his Death and Resurrection, such as Jonah and the Whale and Daniel in the lions' den. Many Nativity scenes occur, the theme of the Three Kings being particularly popular. The cloisters of Santo Domingo de Silos Abbey in Northern Spain, and Moissac are fine examples surviving complete. Details of the portal of Oloron Cathedral show a demon, a lion swallowing a man and kings with musical instruments. A relief from St Trophime, Arles, showing King Herod and the Three Kings, follows the conventions in that the seated Herod is much larger than the standing figures. Notre-Dame-en-Vaux, Châlons-en-Champagne. This paired capital representing Christ washing the feet of the disciples is lively and naturalistic. The large wall surfaces and plain curving vaults of the Romanesque period lent themselves to mural decoration. Unfortunately, many of these early wall paintings have been destroyed by damp or the walls have been replastered and painted over. In most of Northern Europe such pictures were systematically destroyed in bouts of Reformation iconoclasm. In other countries they have suffered from war, neglect and changing fashion. A classic scheme for the full painted decoration of a church, derived from earlier examples often in mosaic, had, as its focal point in the semi-dome of the apse, Christ in Majesty or Christ the Redeemer enthroned within a mandorla and framed by the four winged beasts, symbols of the Four Evangelists, comparing directly with examples from the gilt covers or the illuminations of Gospel Books of the period. If the Virgin Mary was the dedicatee of the church, she might replace Christ here. On the apse walls below would be saints and apostles, perhaps including narrative scenes, for example of the saint to whom the church was dedicated. On the sanctuary arch were figures of apostles, prophets or the twenty-four "elders of the Apocalypse", looking in towards a bust of Christ, or his symbol the Lamb, at the top of the arch. The north wall of the nave would contain narrative scenes from the Old Testament, and the south wall from the New Testament. On the rear west wall would be a Doom painting or Last Judgement, with an enthroned and judging Christ at the top. One of the most intact schemes to exist is that at Saint-Savin-sur-Gartempe in France. (See picture above under "Vault") The long barrel vault of the nave provides an excellent surface for fresco, and is decorated with scenes of the Old Testament, showing the Creation, the Fall of Man and other stories including a lively depiction of Noah's Ark complete with a fearsome figurehead and numerous windows through with can be seen the Noah and his family on the upper deck, birds on the middle deck, while on the lower are the pairs of animals. Another scene shows with great vigour the swamping of Pharaoh's army by the Red Sea. The scheme extends to other parts of the church, with the martyrdom of the local saints shown in the crypt, and Apocalypse in the narthex and Christ in Majesty. The range of colours employed is limited to light blue-green, yellow ochre, reddish brown and black. Similar paintings exist in Serbia, Spain, Germany, Italy and elsewhere in France. A frieze of figures occupies the zone below the semi-dome in the apse. Abbey of St Pere of Burgal, Catalonia, Spain In England the major pictorial theme occurs above the chancel arch in parish churches. St John the Baptist, Clayton, Sussex The oldest-known fragments of medieval pictorial stained glass appear to date from the 10th century. The earliest intact figures are five prophet windows at Augsburg, dating from the late 11th century. The figures, though stiff and formalised, demonstrate considerable proficiency in design, both pictorially and in the functional use of the glass, indicating that their maker was well accustomed to the medium. At Canterbury and Chartres Cathedrals, a number of panels of the 12th century have survived, including, at Canterbury, a figure of Adam digging, and another of his son Seth from a series of Ancestors of Christ. Adam represents a highly naturalistic and lively portrayal, while in the figure of Seth, the robes have been used to great decorative effect, similar to the best stone carving of the period. Many of the magnificent stained glass windows of France, including the famous windows of Chartres, date from the 13th century. Far fewer large windows remain intact from the 12th century. One such is the Crucifixion of Poitiers, a remarkable composition that rises through three stages, the lowest with a quatrefoil depicting the Martyrdom of St Peter, the largest central stage dominated by the crucifixion and the upper stage showing the Ascension of Christ in a mandorla. The figure of the crucified Christ is already showing the Gothic curve. The window is described by George Seddon as being of "unforgettable beauty". King David from Augsburg Cathedral, late 11th century. One of a series of prophets that are the oldest stained glass windows in situ. Two panels of lively figures, Seth and Adam from the 12th-century Ancestors of Christ, Canterbury Cathedral, now set into a Perpendicular Gothic window with panels of many different dates. King Otto II from a series of Emperors (12th and 13th centuries) The panels are now set into Gothic windows, Strasbourg Cathedral Detail of a small panel showing Kings David and Solomon set in an architectonic frame from a large window at Strasbourg. Late 12th century. The alternation of red and blue is a typical device of simpler window designs. It is approximately 1/3 the height, and is much less complex in execution than the Emperor series of which Otto II is a part.See left Transitional style and the continued use of Romanesque forms During the 12th century, features that were to become typical of Gothic architecture began to appear. It is not uncommon, for example, for a part of building that has been constructed over a lengthy period extending into the 12th century, to have very similar arcading of both semi-circular and pointed shape, or windows that are identical in height and width, but in which the later ones are pointed. This can be seen on the towers of Tournai Cathedral and on the western towers and facade at Ely Cathedral. Other variations that appear to hover between Romanesque and Gothic occur, such as the facade designed by Abbot Suger at the Abbey of Saint-Denis, which retains much that is Romanesque in its appearance, and the Facade of Laon Cathedral, which, despite its Gothic form, has round arches. Abbot Suger's innovative choir of the Abbey of Saint-Denis, 1140–44, led to the adoption of the Gothic style by Paris and its surrounding area, but other parts of France were slower to take it up, and provincial churches continued to be built in the heavy manner and rubble stone of the Romanesque, even when the openings were treated with the fashionable pointed arch. In England, the Romanesque groundplan, which in that country commonly had a very long nave, continued to affect the style of building of cathedrals and those large abbey churches which were also to become cathedrals at the dissolution of the monasteries in the 16th century. Despite the fact that English cathedrals were built or rebuilt in many stages, substantial areas of Norman building can be seen in many of them, particularly in the nave arcades. In the case of Winchester Cathedral, the Gothic arches were literally carved out of the existent Norman piers. Other cathedrals have sections of their building which are clearly an intermediate stage between Norman and Gothic, such as the western towers of Ely Cathedral and part of the nave at Worcester Cathedral. The first truly Gothic building in England is the long eastern end of Canterbury Cathedral commenced in 1175. In Italy, although many churches such as Florence Cathedral and Santa Maria Novella were built in the Gothic style, or utilising the pointed arch and window tracery, Romanesque features derived from the Roman architectural heritage, such as sturdy columns with capitals of a modified Corinthian form, continued to be used. The pointed vault was utilised where convenient, but it is commonly interspersed with semicircular arches and vaults wherever they conveniently fit. The facades of Gothic churches in Italy are not always easily distinguishable from the Romanesque. Germany was not quick to adopt the Gothic style, and when it did so in the 1230s, the buildings were often modelled very directly upon French cathedrals, as Cologne Cathedral was modelled on Amiens. The smaller churches and abbeys continued to be constructed in a more provincial Romanesque manner, the date only being registered by the pointed window openings. The facade of Laon Cathedral, 1225, a Gothic cathedral, maintains rounded arches and arcading in the Romanesque manner. Ely Cathedral, England, the central western tower and framing smaller towers all had transitional features, 1180s. The tower to the left fell. Gothic porch, 1250s; lantern, 1390s. The facade of the Cathedral of Genoa has both round and pointed arches, and paired windows, a continuing Romanesque feature of Italian Gothic architecture. Romanesque castles, houses and other buildings The Romanesque period was a time of great development in the design and construction of defensive architecture. After churches and the monastic buildings with which they are often associated, castles are the most numerous type of building of the period. While most are in ruins through the action of war and politics, others, like William the Conqueror's White Tower within the Tower of London have remained almost intact. In some regions, particularly Germany, large palaces were built for rulers and bishops. Local lords built great halls in the countryside, while rich merchants built grand town houses. In Italy, city councils constructed town halls, while wealthy cities of Northern Europe protected their trading interests with warehouses and commercial premises. All over Europe, dwellers of the town and country built houses to live in, some of which, sturdily constructed in stone, have remained to this day with sufficient of their form and details intact to give a picture of the style of domestic architecture that was in fashion at the time. Examples of all these types of buildings can be found scattered across Europe, sometimes as isolated survivals like the two merchants' houses on opposite sides of Steep Hill in Lincoln, England, and sometimes giving form to a whole medieval city like San Gimignano in Tuscany, Italy. These buildings are the subject of a separate article. During the 19th century, when Gothic Revival architecture was fashionable, buildings were occasionally designed in the Romanesque style. There are a number of Romanesque Revival churches, dating from as early as the 1830s and continuing into the 20th century where the massive and "brutal" quality of the Romanesque style was appreciated and designed in brick. The Natural History Museum, London, designed by Alfred Waterhouse, 1879, on the other hand, is a Romanesque revival building that makes full use of the decorative potential of Romanesque arcading and architectural sculpture. The Romanesque appearance has been achieved while freely adapting an overall style to suit the function of the building. The columns of the foyer, for example, give an impression of incised geometric design similar to those of Durham Cathedral. However, the sources of the incised patterns are the trunks of palms, cycads and tropical tree ferns. The animal motifs, of which there are many, include rare and exotic species. The type of modern buildings for which the Romanesque style was most frequently adapted was the warehouse, where a lack of large windows and an appearance of great strength and stability were desirable features. These buildings, generally of brick, frequently have flattened buttresses rising to wide arches at the upper levels after the manner of some Italian Romanesque facades. This style was adapted to suit commercial buildings by opening the spaces between the arches into large windows, the brick walls becoming a shell to a building that was essentially of modern steel-frame construction, the architect Henry Hobson Richardson giving his name to the style, Richardsonian Romanesque. Good examples of the style are Marshall Field's Wholesale Store, Chicago, by H.H. Richardson, 1885, and the Chadwick Lead Works in Boston, USA, by William Preston, 1887. The style also lent itself to the building of cloth mills, steelworks and powerstations. The 19th-century reconstruction of the westwerk of the Romanesque Speyer Cathedral. see above - The traceried window to the left of the building indicates that the steeply gabled vestry dates from the Gothic period. - Gerville (1818): Je vous ai quelquefois parlé d'architecture romane. C’est un mot de ma façon qui me paraît heureusement inventé pour remplacer les mots insignifiants de saxone et de normande. Tout le monde convient que cette architecture, lourde et grossière, est l'opus romanum dénaturé ou successivement dégradé par nos rudes ancêtres. Alors aussi, de la langue latine, également estropiée, se faisait cette langue romane dont l'origine et la dégradation ont tant d'analogie avec l'origine et les progrès de l'architecture. Dites-moi donc, je vous prie, que mon nom romane est heureusement trouvé. English: I have sometimes spoken to you about Romanesque architecture. It is a word of my own which I invented (I think successfully) to replace the insignificant words of Saxon and Norman. Everyone agrees that this architecture, heavy and rough, is the opus romanum successively denatured or degraded by our rude ancestors. So too, out of the crippled Latin language, was made this Romance language whose origin and degradation have so much analogy with the origin and progress of architecture. Tell me, please, that my name Roman (esque) was invented with success. - de Caumont (1824): Le nom romane que nous donnons à cette architecture, qui ne doit avoir qu'un puisqu'elle est partout la même sauf de légères differences de localité, a d'ailleurs le mérite d'en indiquer l'origine et il n'est pas nouveau puisqu'on s'en sert déjà pour désigner la langue du même temps La langue romane est la langue latine dégénérée. L'architecture romane est l'architecture romaine abâtardie. ( English: The name Roman (esque) we give to this architecture, which should be universal as it is the same everywhere with slight local differences, also has the merit of indicating its origin and is not new since it is used already to describe the language of the same period. Romance language is degenerated Latin language. Romanesque architecture is debased Roman architecture) |Wikimedia Commons has media related to Romanesque architecture.| - List of Romanesque buildings - List of regional characteristics of Romanesque churches - Romanesque secular and domestic architecture - Portuguese Romanesque architecture - Romanesque art - Romanesque sculpture - Spanish Romanesque - Romanesque Revival Architecture in the United Kingdom - Bannister Fletcher, A History of Architecture on the Comparative Method. - Gidon 1934, p. 285-286 - Gidon, Ferdinand (1934). "L’invention de l’expression architecture romane par Gerville (1818) d’après quelques lettres de Gerville à Le Prévost". Bulletin de la Société des antiquaires de Normandie (in French) 42: 268–288. - de Caumont, Arcisse (8 May 1824). "Essai sur l'architecture religieuse du moyen-âge, particulièrement en Normandie". Mémoires de la Société des antiquaires de Normandie (in French) (Mancel): 535–677. Retrieved 2012-06-24. - Williams, Elizabeth (1 January 1985). "The perception of romanesque art in the romantic period: archaeological attitudes in france in the 1820s and1830s". Forum for Modern Language Studies XXI (4): 303–321. doi:10.1093/fmls/XXI.4.303. - Jean Hubert, Romanesque Art. - Date from Hartmann-Virnich, as below - de Caumont 1824, p. 550 - Gunn, William (1819). An inquiry into the origin and influence of Gothic architecture. R. and A. Taylor. p. 6. Retrieved 2012-07-06. - Andreas Hartmann-Virnich: Was ist Romanik, Darmstadt 2004, p. 28-30 - Rolf Toman, Romanesque: Architecture, Sculpture, Painting - Helen Gardner, Art through the Ages. - George Holmes, ed. The Oxford History of Medieval Europe. - Rolf Toman, pp. 114-117 - Copplestone, pp.188-89 - Rolf Toman, pp. 70-73 - Rolf Toman, pp. 18, 177, 188 - "In the years that followed the year 1000, we witnessed the rebuilding of churches all over the universe, but especially in Italy and Gaul." Chronicle of Raoul Glaber, quoted by Jean Hubert, Romanesque Art. - famous for the ancient Roman "Mouth of Truth" set into the wall of its narthex - famous for the 15th-century Ghiberti Doors - traditionally the marriage place of Romeo and Juliet - John Harvey, English Cathedrals - Alec Clifton-Taylor, The Cathedrals of England - Rolf Toman, Romanesque. - "Architecture". National Tourism Organisation of Serbia. Retrieved 2007-09-28. - Rene Hyughe, Larousse Encyclopedia of Byzantine and Medieval Art - This technique was also used in the Classical world, notably at the Parthenon. - Nikolaus Pevsner, An Outline of European Architecture - Banister Fletcher, p.307 - Stephenson, Hammond & Davi 2005, p. 172. - Jones, Murray & Murray 2013, p. 512. - Porter 1928, p. 48. - Kimball, F., & Edgell, G. H. (1918). A History of Architecture. New York. Harper & Brothers. 621 pages (page 252). - With the exception of the Plan of St. Gall, which is from an ancient manuscript (and probably does not reflect an actual construction), they are all hypothetical reconstructions of groundplans as they existed in the 12th or 13th centuries. The Abbey Church of St. Gall has been replaced by a Baroque Church. Speyer has had its west front rebuilt twice, Ely Cathedral has lost the eastern arm, being replaced in the Gothic style, the central tower being replaced with the unique octagon and the northwest tower, never rebuilt. It has also gained a west porch. Santiago has had some substantial changes including a Baroque west front. - Crossley, Frederick H. (1962). The English Abbey. - Banister Fletcher p. 309 - "Romànic de la Vall de Camprodon". Elripolles.com. 2010-03-09. Retrieved 2011-06-11. - Alec Clifton-Taylor says "With the Cathedral of Durham we reach the incomparable masterpiece of Romanesque architecture not only in England but anywhere." - See details at Cologne Cathedral. - Howe, Jeffery. "Romanesque Architecture (slides)". A digital archive of architecture. Boston College. Retrieved 2007-09-28. - James Hall, A History of Ideas and Images in Italian Art, p154, 1983, John Murray, London, ISBN 0-7195-3971-4 - George Seddon in Lee, Seddon and Stephens, Stained Glass - Wim Swaan, Gothic Cathedrals - Conant, Kenneth J., Carolingian and Romanesque Architecture: 800 to 1200 (4th, illustrated, reprint ed.). Yale University Press. 1993. ISBN 978-0-300-05298-5. - V.I. Atroshenko and Judith Collins, The Origins of the Romanesque, Lund Humphries, London, 1985, ISBN 0-85331-487-X - Rolf Toman, Romanesque: Architecture, Sculpture, Painting, Könemann, (1997), ISBN 3-89508-447-6 - Banister Fletcher, A History of Architecture on the Comparative method (2001). Elsevier Science & Technology. ISBN 0-7506-2267-9. - Alfred Clapham, Romanesque Architecture in England British Council (1950) - Helen Gardner; Fred S. Kleiner, Christin J. Mamiya, Gardner's Art through the Ages. Thomson Wadsworth, (2004) ISBN 0-15-505090-7. - George Holmes, editor, The Oxford Illustrated History of Medieval Europe, Oxford University Press, (1992) ISBN 0-19-820073-0 - René Huyghe, Larousse Encyclopedia of Byzantine and Medieval Art, Paul Hamlyn, (1958) - François Ischer, Building the Great Cathedrals. Harry N. Abrams, (1998). ISBN 0-8109-4017-5. - Jones, Tom Devonshire; Murray, Linda; Murray, Peter, eds. (2013). The Oxford Dictionary of Christian Art and Architecture (illustrated ed.). Oxford University Press. ISBN 978-0-199-68027-6. - Nikolaus Pevsner, An Outline of European Architecture. Pelican Books (1964) - Porter, Arthur Kingsley (1928). Spanish Romanesque Sculpture, Volume 1 (illustrated ed.). Hacker Art Books. - John Beckwith, Early Medieval Art, Thames and Hudson, (1964) - Peter Kidson, The Medieval World, Paul Hamlyn, (1967) - T. Francis Bumpus,, The Cathedrals and Churches of Belgium, T. Werner Laurie. (1928) - Alec Clifton-Taylor, The Cathedrals of England, Thames and Hudson (1967) - John Harvey, English Cathedrals, Batsford (1961). - Stephenson, Davis; Hammond, Victoria; Davi, Keith F. (2005). Visions of Heaven: the Dome in European Architecture (illustrated ed.). Princeton Architectural Press. p. 174. ISBN 978-1-56898-549-7. - Trewin Copplestone, World Architecture, and Illustrated History, Paul Hamlyn, (1963) - Tadhg O'Keefe, Archeology and the Pan-European Romanesque , Duckworth Publishers, (2007), ISBN 0715634348 |Look up Romanesque in Wiktionary, the free dictionary.| - Corpus of Romanesque Sculpture in Britain and Ireland - Overview of French Romanesque art - French Romanesque art through 300 places (Italian) (French) (Spanish) (English) - Romanesque Churches in Southern Burgundy - Spanish and Zamora´s Romanesque art, easy navigation (Spanish) - Spanish Romanesque art (Spanish) - Círculo Románico - Visigothic, Mozarabic and Romanesque art in Europe - Romanesque Churches in Portugal - The Nine Romanesque Churches of the Vall de Boi - Pyrenees (English) - Satan in the Groin - exhibitionist carvings on medieval churches - An illustrated article by Peter Hubert on the cusped arch - Corrèze Illustrated history (French) - The Encyclopedia of Romanesque Art in Spain: a work in progress (Spanish) - Saint-Trophime Digital Media Archive (creative commons-licensed HD documentation) on the Romanesque Church of St. Trophime, using data from a World Monuments Fund/CyArk research partnership - Cerisy-la-Forêt abbey, a masterwork of French Norman architecture

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Calculating the inverse of a linear function is easy: just make x the subject of the equation, and replace y with x in the resulting expression. Finding the inverse of a quadratic function is considerably trickier, not least because Quadratic functions are not, unless limited by a suitable domain, one-one functions. 1Make y or f(x) the subject of the formula if it isn't already. During your algebraic manipulation, make sure that you do not change the function in any way and perform the same operations to both "sides" of the equation. 2Rearrange the function so that it is in the form y=a(x-h)2+k. This is not only essential for you to find the inverse of the function, but also for you to determine whether the function even has an inverse. You can do this by two methods: - By completing the square - "Take common" from the whole equation the value of a (the coefficient of x2). Do this by writing the value of a, starting a bracket, and writing the whole equation, then dividing each term by the value of a, as shown in the diagram on the right. Leave the left hand side of the equation untouched, as there has been no net change to the right hand side. - Complete the square. The coefficient of x is (b/a). Halve it, to give (b/2a), and square it, to give (b/2a)2. Add and subtract it from the equation. This will have no net effect on the equation. If you look closely, you will see that the first three terms inside the bracket are in the form a2+2ab+b2, where a is x, and b is (b/2a). Of course these two values will be numerical, rather than algebraic for a real equation. This is a completed square. - Because the first three terms are now a perfect square, you can write them in the form (a-b)2 or (a+b)2. The sign between the two terms will be the same as the sign of the coefficient of x in the equation. - Take the term which is outside the perfect square out of the square bracket. This brings the equation into the form y=a(x-h)2+k, as intended. - By comparing coefficients - Form an identity in x. On the left, put the function as it is expressed in terms of x, and on the right put the function in the form that you want it to be, in this case a(x-h)2+k. This will enable you to find out the values of a, h, and k that are true for all values of x. - Open and expand the bracket on the right hand side of the identity. We shall not be touching the left hand side of the equation, and may omit it from our working. Note that all working on the right hand side is algebraic as shown and not numerical. - Identify the coefficients of each power of x. Then group them and place them in brackets, as shown on the right. - Compare the coefficients of each power of x. The coefficient of x2 on the right hand side must equal that on the left hand side. This gives the value of a. The coefficient of x on the right hand side also must equal that on the left hand side. This leads to the formation of an equation in a and h, which can be solved by substituting the value of a, which has already been found. The coefficient of x0, or 1, on the left hand side must also equal that on the right hand side. Comparing them yields an equation that will help us find the value of k. - Using the values of a,h, and k found above, we can write the equation in the desired form. - By completing the square 3Ensure that the value of h is either on the boundary of the domain, or outside it. The value of h gives the x-coordinate of the turning point of the equation. A turning point within the domain would mean that the function is not one-one, and hence does not have an inverse. Note that the equation is a(x-h)2+k. Thus if there is (x+3) inside the bracket, the value of h is negative 3. 4Make (x-h)2 the subject of your formula. Do this by subtracting the value of k from both sides of the equation, and then dividing both sides of the equation by a. By now you will have numerical values for a,h, and k, so use those, not the symbols. 5Square-Root both sides of the equation. This will remove the power of two from (x-h). Do not forget to put the "+/-" sign on the other side of the equation. 6Decide between the + and the - sign, as you can not have both (having both would make it a one to many "function", which would make it invalid as the same). For this, look at the domain. If the domain lies to the left of the stationary point i.e. x < a certain value, use the - sign. If the domain lies to the right of the stationary point i.e. x > a certain value, use the + sign. Then, make x the subject of the formula. 7Replace y with x, and x with f-1(x), and congratulate yourself on having successfully found the inverse of a quadratic function. Questions and Answers Give us 3 minutes of knowledge! - Check your inverse by calculating the value of f(x) for a certain value of x, and then substituting that value of f(x) in the inverse to see if it returns the original value of x. For example, if the function of 3 [f(3)] is 4, then substituting 4 in the inverse should return 3. - If it is not too much trouble you can also check the inverse by inspecting its graph. It should look like the original function reflected across the line y=x. In other languages: Español: encontrar la inversa de una función cuadrática, Italiano: Trovare l'Inversa di una Funzione Quadratica, Русский: найти функцию, обратную квадратичной функции, Português: Encontrar o Inverso de uma Função Quadrática Thanks to all authors for creating a page that has been read 182,769 times.

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Gardening Level 2: Let's Get Growing This activity book is Level 2 (B) of the 4-H Gardening Curriculum series, written by university experts. This level focuses on the following gardening skills: using transplants in a garden, developing a planting calendar to start seeds indoors, understanding plant responses, growing plants from plant parts, making a worm box, making compost, judging vegetables, growing vegetables for cash. Target Age: Grades 7-8. Unit 1: Let's Plan! Unit 2: Dig In (Planting) Unit 3: While You Wait Unit 4: Watch Out! (Garden Care) Unit 5: Now What? (Harvesting and Storage) Unit 6: Imagine That! (Careers) You might also be interested in...

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What is osteosarcoma? Cancer starts when cells in the body begin to grow out of control. Cells in nearly any part of the body can become cancer, and can spread to other areas of the body. To learn more about how cancers start and spread, see What Is Cancer? Osteosarcoma (also called osteogenic sarcoma) is a type of cancer that starts in the bones. To understand osteosarcoma, it helps to know about bones and what they do. About normal bones Many people think of bones as just being part of the skeleton, like the steel girders that support a building. But bones actually do a number of different things. - Some bones help support and protect our vital organs. Examples include the skull bones, breast bone (sternum), and ribs. These types of bones are often referred to as flat bones. - Other bones, such as those in the arms and legs, make a framework for our muscles that helps us move. These are called long bones. - Bones also make new blood cells. This is done in the soft, inner part of some bones called the bone marrow, which contains blood-forming cells. New red blood cells, white blood cells, and platelets are made in bone marrow. - Bones also provide the body with a place to store minerals such as calcium. Because bones are very hard and don’t change shape − at least once we reach adulthood − we might not think of bones as being alive, but they are. Like all other tissues of the body, bones have many kinds of living cells. Two main types of cells in our bones help them stay strong and keep their shape. - Osteoblasts help build up bones by forming the bone matrix (the connective tissue and minerals that give bone its strength). - Osteoclasts break down bone matrix to prevent too much of it from building up, and they help bones keep their proper shape. By depositing or removing minerals from the bones, osteoblasts and osteoclasts also help control the amount of these minerals in the blood. Osteosarcoma is the most common type of cancer that develops in bone. Like the osteoblasts in normal bone, the cells that form this cancer make bone matrix. But the bone matrix of an osteosarcoma is not as strong as that of normal bones. Most osteosarcomas occur in children and young adults. Teens are the most commonly affected age group, but osteosarcoma can occur at any age. In children and young adults, osteosarcoma usually develops in areas where the bone is growing quickly, such as near the ends of the long bones. Most tumors develop in the bones around the knee, either in the distal femur (the lower part of the thigh bone) or the proximal tibia (the upper part of the shinbone). The proximal humerus (the part of the upper arm bone close to the shoulder) is the next most common site. However, osteosarcoma can develop in any bone, including the bones of the pelvis (hips), shoulder, and jaw. This is especially true in older adults. Subtypes of osteosarcoma Several subtypes of osteosarcoma can be identified by how they look on x-rays and under the microscope. Some of these subtypes have a better prognosis (outlook) than others. Based on how they look under the microscope, osteosarcomas can be classified as high grade, intermediate grade, or low grade. The grade of the tumor tells doctors how likely it is that the cancer will grow and spread to other parts of the body. High-grade osteosarcomas: These are the fastest growing types of osteosarcoma. When seen under a microscope, they do not look like normal bone and have many cells in the process of dividing into new cells. Most osteosarcomas that occur in children and teens are high grade. There are many types of high-grade osteosarcomas (although the first 3 are the most common). - Small cell - High-grade surface (juxtacortical high grade) Other high-grade osteosarcomas include: - Pagetoid: a tumor that develops in someone with Paget disease of the bone - Extra-skeletal: a tumor that starts in a part of the body other than a bone - Post-radiation: a tumor that starts in a bone that had once received radiation therapy Intermediate-grade osteosarcomas: These uncommon tumors fall in between high-grade and low-grade osteosarcomas. (They are usually treated as if they are low-grade osteosarcomas.) - Periosteal (juxtacortical intermediate grade) Low-grade osteosarcomas: These are the slowest growing osteosarcomas. The tumors look more like normal bone and have few dividing cells when seen under a microscope. - Parosteal (juxtacortical low grade) - Intramedullary or intraosseous well differentiated (low-grade central) The grade of the tumor plays a role in determining its stage and the type of treatment used. For more on staging, see the section “How is osteosarcoma staged?” Other types of bone tumors Several other types of tumors can start in the bones. Malignant (cancerous) bone tumors Ewing tumors are the second most common bone cancer in children. They are described in our document Ewing Family of Tumors. Most other types of bone cancers are usually found in adults and are rare in children. These include: - Chondrosarcoma (cancer that develops from cartilage) - Malignant fibrous histiocytoma - Malignant giant cell tumor of bone For more information on these cancers, see our document Bone Cancer. Many types of cancer that start in other organs of the body can spread to the bones. These are sometimes referred to as metastatic bone cancers, but they are not true bone cancers. For example, prostate cancer that spreads to the bones is still prostate cancer and is treated like prostate cancer. For more information, see the document Bone Metastasis. Benign (non-cancerous) bone tumors Not all bone tumors are cancer. Benign bone tumors do not spread to other parts of the body. They are usually not life threatening and can often be cured by surgery. There are many types of benign bone tumors. - Osteomas are benign tumors formed by bone cells. - Chondromas are benign tumors formed by cartilage cells. - Osteochondromas are benign tumors with both bone and cartilage cells. Other benign bone tumors include eosinophilic granuloma of bone, non-ossifying fibroma, enchondroma, xanthoma, benign giant cell tumor of bone, and lymphangioma. The rest of this document covers only osteosarcoma. Last Medical Review: 04/18/2014 Last Revised: 01/27/2016

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Watching this resources will notify you when proposed changes or new versions are created so you can keep track of improvements that have been made. Favoriting this resource allows you to save it in the “My Resources” tab of your account. There, you can easily access this resource later when you’re ready to customize it or assign it to your students. The act of bending a joint, especially a bone joint. The counteraction of extension. Not only are a variety of movements possible with synovial joints, but in order to maintain flexibility, these joints need to be moved daily. Failure to maintain flexibility of joints makes movement more difficult and increases the probability of falls and injuries. A synovial joint, also known as a diarthrosis, is the most common and most movable type of joint in the body of a mammal. As with most other joints, synovial joints achieve movement at the point of contact of the articulating bones. Structural and functional differences distinguish synovial joints from cartilaginous joints (synchondroses and symphyses) and fibrous joints (sutures, gomphoses, and syndesmoses). The main structural differences between synovial and fibrous joints are the existence of capsules surrounding the articulating surfaces of a synovial joint and the presence of lubricating synovial fluid within those capsules (synovial cavities). Several movements may be performed by synovial joints. Abduction is the movement away from the mid-line of the body. Adduction is the movement toward the middle line of the body . Extension is the straightening of limbs (increase in angle) at a joint. Flexion is bending the limbs (reduction of angle) at a joint . Rotation is a circular movement around a fixed point. There are six types of synovial joints. Some are relatively immobile, but are more stable. Others have multiple degrees of freedom, but at the expense of greater risk of injury. The six types of joints include: Gliding joints, which only allow sliding movement Hinge joints, which allow flexion and extension in one plane Pivot joints, which allow bone rotation about another Condyloid joints, which allow flexion, extension, abduction, and adduction movements Saddle joints, which permit the same movement as condyloid joints (and condylar joints and saddle joints combine to form compound joints) Ball and socket joints, which allow all movements except gliding synovial joints, unlike cartilaginous and fibrous joints, achieve movement at bony points of contact, synovial joints are the most moveable of the body joints, synovial joints have lubricating capsules surrounding the articular surfaces of the joints, or synovial joints are also known as diarthroses

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A concretion is a hard, compact mass of matter formed by the precipitation of mineral cement within the spaces between particles, and is found in sedimentary rock or soil. Concretions are often ovoid or spherical in shape, although irregular shapes also occur. The word 'concretion' is derived from the Latin con meaning 'together' and crescere meaning 'to grow'. Concretions form within layers of sedimentary strata that have already been deposited. They usually form early in the burial history of the sediment, before the rest of the sediment is hardened into rock. This concretionary cement often makes the concretion harder and more resistant to weathering than the host stratum. There is an important distinction to draw between concretions and nodules. Concretions are formed from mineral precipitation around some kind of nucleus while a nodule is a replacement body. Descriptions dating from the 18th century attest to the fact that concretions have long been regarded as geological curiosities. Because of the variety of unusual shapes, sizes and compositions, concretions have been interpreted to be dinosaur eggs, animal and plant fossils (called pseudofossils), extraterrestrial debris or human artifacts. - 1 Origins - 2 Appearance - 3 Composition - 4 Occurrence - 5 Types of concretion - 6 Gallery - 7 See also - 8 Citations - 9 References - 10 External links Detailed studies (i.e., Boles et al., 1985; Thyne and Boles, 1989; Scotchman, 1991; Mozley and Burns, 1993; McBride et al., 2003; Chan et al., 2005; Mozley and Davis, 2005) published in peer-reviewed journals have demonstrated that concretions form subsequent to burial during diagenesis. They quite often form by the precipitation of a considerable amount of cementing material around a nucleus, often organic, such as a leaf, tooth, piece of shell or fossil. For this reason, fossil collectors commonly break open concretions in their search for fossil animal and plant specimens. Some of the most unusual concretion nuclei, as documented by Al-Agha et al. (1995), are World War II military shells, bombs, and shrapnel, which are found inside siderite concretions found in an English coastal salt marsh. Depending on the environmental conditions present at the time of their formation, concretions can be created by either concentric or pervasive growth (Mozley, 1996; Raiswell and Fisher, 2000). In concentric growth, the concretion grows as successive layers of mineral accrete to its surface. This process results in the radius of the concretion growing with time. In case of pervasive growth, cementation of the host sediments, by infilling of its pore space by precipitated minerals, occurs simultaneously throughout the volume of the area, which in time becomes a concretion. Concretions vary in shape, hardness and size, ranging from objects that require a magnifying lens to be clearly visible to huge bodies three meters in diameter and weighing several thousand pounds. The giant, red concretions occurring in Theodore Roosevelt National Park, in North Dakota, are almost 3 m (9.8 ft) in diameter. Spheroidal concretions, as large as 9 m (30 ft) in diameter, have been found eroding out of the Qasr El Sagha Formation within the Faiyum depression of Egypt. Concretions are usually similar in color to the rock in which they are found. Concretions occur in a wide variety of shapes, including spheres, disks, tubes, and grape-like or soap bubble-like aggregates. They are commonly composed of a carbonate mineral such as calcite; an amorphous or microcrystalline form of silica such as chert, flint, or jasper; or an iron oxide or hydroxide such as goethite and hematite. They can also be composed of other minerals that include dolomite, ankerite, siderite, pyrite, marcasite, barite and gypsum. Although concretions often consist of a single dominant mineral, other minerals can be present depending on the environmental conditions which created them. For example, carbonate concretions, which form in response to the reduction of sulfates by bacteria, often contain minor percentages of pyrite. Other concretions, which formed as a result of microbial sulfate reduction, consist of a mixture of calcite, barite, and pyrite. Concretions are found in a variety of rocks, but are particularly common in shales, siltstones, and sandstones. They often outwardly resemble fossils or rocks that look as if they do not belong to the stratum in which they were found. Occasionally, concretions contain a fossil, either as its nucleus or as a component that was incorporated during its growth but concretions are not fossils themselves. They appear in nodular patches, concentrated along bedding planes, protruding from weathered cliffsides, randomly distributed over mudhills or perched on soft pedestals. Types of concretion Concretions vary considerably in their compositions, shapes, sizes and modes of origin. Septarian concretions or septarian nodules, are concretions containing angular cavities or cracks, called "septaria". The word comes from the Latin word septum; "partition", and refers to the cracks/separations in this kind of rock. There is an incorrect explanation that it comes from the Latin word for "seven", septem, referring to the number of cracks that commonly occur. Cracks are highly variable in shape and volume, as well as the degree of shrinkage they indicate. Although it has commonly been assumed that concretions grew incrementally from the inside outwards, the fact that radially oriented cracks taper towards the margins of septarian concretions is taken as evidence that in these cases the periphery was stiffer while the inside was softer, presumably due to a gradient in the amount of cement precipitated. The process that created the septaria, which characterize septarian concretions, remains a mystery. A number of mechanisms, e.g. the dehydration of clay-rich, gel-rich, or organic-rich cores; shrinkage of the concretion's center; expansion of gases produced by the decay of organic matter; brittle fracturing or shrinkage of the concretion interior by either earthquakes or compaction; and others, have been proposed for the formation of septaria (Pratt 2001). At this time, it is uncertain, which, if any, of these and other proposed mechanisms is responsible for the formation of septaria in septarian concretions (McBride et al. 2003). Septaria usually contain crystals precipitated from circulating solutions, usually of calcite. Siderite or pyrite coatings are also occasionally observed on the wall of the cavities present in the septaria, giving rise respectively to a panoply of bright reddish and golden colors. Some septaria may also contain small calcite stalactites and well-shaped millimetric pyrite single crystals. A spectacular example of septarian concretions, which are as much as 3 meters (9.8 feet) in diameter, are the Moeraki Boulders. These concretions are found eroding out of Paleocene mudstone of the Moeraki Formation exposed along the coast near Moeraki, South Island, New Zealand. They are composed of calcite-cemented mud with septarian veins of calcite and rare late-stage quartz and ferrous dolomite (Boles et al. 1985, Thyne and Boles 1989). Very similar concretions, which are as much as 3 meters (9.8 feet) in diameter and called "Koutu Boulders", litter the beach between Koutu and Kauwhare points along the south shore of the Hokianga Harbour of Hokianga, North Island, New Zealand. The much smaller septarian concretions found in the Kimmeridge Clay exposed in cliffs along the Wessex Coast of England are more typical examples of septarian concretions (Scotchman 1991). Cannonball concretions are large spherical concretions, which resemble cannonballs. These are found along the Cannonball River within Morton and Sioux Counties, North Dakota, and can reach 3 m (9.8 ft) in diameter. They were created by early cementation of sand and silt by calcite. Similar cannonball concretions, which are as much as 4 to 6 m (13 to 20 ft) in diameter, are found associated with sandstone outcrops of the Frontier Formation in northeast Utah and central Wyoming. They formed by the early cementation of sand by calcite (McBride et al. 2003). Somewhat weathered and eroded giant cannonball concretions, as large as 6 meters (20 feet) in diameter, occur in abundance at "Rock City" in Ottawa County, Kansas. Large and spherical boulders are also found along Koekohe beach near Moeraki on the east coast of the South Island of New Zealand. The Moeraki Boulders and Koutu Boulders of New Zealand are examples of septarian concretions, which are also cannonball concretions. Large spherical rocks, which are found on the shore of Lake Huron near Kettle Point, Ontario, and locally known as "kettles", are typical cannonball concretions. Cannonball concretions have also been reported from Van Mijenfjorden, Spitsbergen; near Haines Junction, Yukon Territory, Canada; Jameson Land, East Greenland; near Mecevici, Ozimici, and Zavidovici in Bosnia-Herzegovina; in Alaska in the Kenai Peninsula Captain Cook State Park on north of Cook Inlet beach. Reports of cannonball concretions have also come from Bandeng and Zhanlong hills near Gongxi Town, Hunan Province, China. Hiatus concretions are distinguished by their stratigraphic history of exhumation, exposure and reburial. They are found where submarine erosion has concentrated early diagenetic concretions as lag surfaces by washing away surrounding fine-grained sediments (Zaton 2010). Their significance for stratigraphy, sedimentology and paleontology was first noted by Voigt (1968) who referred to them as Hiatus-Konkretionen. "Hiatus" refers to the break in sedimentation that allowed this erosion and exposure. They are found throughout the fossil record but are most common during periods in which calcite sea conditions prevailed, such as the Ordovician, Jurassic and Cretaceous (Zaton 2010). Most are formed from the cemented infillings of burrow systems in siliciclastic or carbonate sediments. A distinctive feature of hiatus concretions separating them from other types is that they were often encrusted by marine organisms including bryozoans, echinoderms and tube worms in the Paleozoic (e.g., Wilson 1985) and bryozoans, oysters and tube worms in the Mesozoic and Cenozoic (e.g., Taylor and Wilson 2001). Hiatus concretions are also often significantly bored by worms and bivalves (Taylor and Wilson 2001). Elongate concretions form parallel to sedimentary strata and have been studied extensively due to the inferred influence of phreatic (saturated) zone groundwater flow direction on the orientation of the axis of elongation (e.g., Johnson, 1989; McBride et al., 1994; Mozley and Goodwin, 1995; Mozley and Davis, 2005). In addition to providing information about the orientation of past fluid flow in the host rock, elongate concretions can provide insight into local permeability trends (i.e., permeability correlation structure; Mozley and Davis, 1996), variation in groundwater velocity (Davis, 1999), and the types of geological features that influence flow. Elongate concretions are well known in the Kimmeridge Clay formation of northwest Europe. In outcrops, where they have acquired the name "doggers", they are typically only a few metres across, but in the subsurface they can be seen to penetrate up to tens of metres of along-hole dimension. Unlike limestone beds, however, it is impossible to consistently correlate them between even closely spaced wells. Moqui Marbles, also called Moqui balls, and "Moki marbles", are iron oxide concretions which can be found eroding in great abundance out of outcrops of the Navajo Sandstone within south-central and southeastern Utah. These concretions range in shape from spheres to discs, buttons, spiked balls, cylindrical forms, and other odd shapes. They range from pea-size to baseball-size. They were created by the precipitation of iron, which was dissolved in groundwater. They are further described by (Chan and Parry 2002, Chan et al. 2005).(Loope et al. 2010,2011) Kansas pop rocks Kansas pop rocks are concretions of either iron sulfide, i.e. pyrite and marcasite, or in some cases jarosite, which are found in outcrops of the Smoky Hill Chalk Member of the Niobrara Formation within Gove County, Kansas. They are typically associated with thin layers of altered volcanic ash, called bentonite, that occur within the chalk comprising the Smoky Hill Chalk Member. A few of these concretions enclose, at least in part, large flattened valves of inoceramid bivalves. These concretions range in size from a few millimeters to as much as 0.7 m (2.3 ft) in length and 12 cm (0.39 ft) in thickness. Most of these concretions are oblate spheroids shape. Other "pop rocks" are small polycuboidal pyrite concretions, which are as much as 7 cm (0.23 ft) in diameter (Hattin 1982). These concretions are called "pop rocks" because they explode if thrown in a fire. Also, when they are either cut or hammered, they produce sparks and a burning sulfur smell. Contrary to what has been published on the Internet, none of the iron sulfide concretions, which are found in the Smoky Hill Chalk Member were created by either the replacement of fossils or by metamorphic processes. In fact, metamorphic rocks are completely absent from the Smoky Hill Chalk Member (Hattin 1982). Instead, all of these the iron sulfide concretions were created by the precipitation of iron sulfides within anoxic marine calcareous ooze after it had accumulated and before it had lithified into chalk. Iron sulfide concretions, such as the Kansas Pop rocks, consisting of either pyrite and marcasite, are nonmagnetic (Hobbs and Hafner 1999). On the other hand, iron sulfide concretions, which either are composed of or contain either pyrrhotite or smythite, will be magnetic to varying degrees (Hoffmann, 1993). Prolonged heating of either a pyrite or marcasite concretion will convert portions of either mineral into pyrrhotite causing the concretion to become slightly magnetic. Calcium carbonate disc concretions These so-called fairy stones consist of single or multiple discs, usually 6–10 cm in diameter and often with concentric grooves on their surfaces. They form in Quaternary clay as calcium carbonate migrates to some small fossil or pebble. Fairy stones are particularly common in the Harricana River valley in the Abitibi-Témiscamingue administrative region of Quebec, and in Östergötland county, Sweden. - Bowling Ball Beach - Calcrete, CaCO3 concretions in arid and semi-arid soils - Caliche (mineral), synonym of calcrete - Dinocochlea in the Natural History Museum, London - Clay dogs - Gypcrust, CaSO4 concretions in arid and semi-arid soils - Klerksdorp sphere - Martian spherules - Moeraki Boulders (New Zealand) - Mushroom Rock State Park, Kansas - Nodule (geology), a replacement body, not to be confused with a concretion - Rock City, Kansas - Speleothems, CaCO3 formations in caves - Glossary of terms in soil science (PDF). Ottawa: Agriculture Canada. 1976. p. 13. ISBN 0662015339. - "septarian". dictionary.reference.com. Retrieved March 20, 2014. - "SEPTARIAN NODULES". Archived from the original on 5 September 2013. - Dann, C., and Peat, N. (1989) Dunedin, North and South Otago. Wellington: GP Books. ISBN 0-477-01438-0 - "The Epoch Times - Mysterious Huge Stone Eggs Discovered in Hunan Province". - Al-Agha, M.R., S.D. Burley, C.D. Curtis, and J. Esson, 1995, Complex cementation textures and authigenic mineral assemblages in Recent concretions from the Lincolnshire Wash (east coast, UK) driven by Fe(0) Fe(II) oxidation: Journal of the Geological Society, London, v. 152, pp. 157–171. - Boles, J.R., C.A. Landis, and P. Dale, 1985, The Moeraki Boulders; anatomy of some septarian concretions:, Journal of Sedimentary Petrology. v. 55, n. 3, pp. 398–406. - Chan, M.A. and W.T. Parry, 2002, 'Mysteries of Sandstone Colors and Concretions in Colorado Plateau Canyon Country PDF version, 468 KB : Utah Geological Survey Public Information Series. n. 77, pp. 1–19. - Chan, M.A., B.B. Beitler, W.T. Parry, J. Ormo, and G. Komatsu, 2005. Red Rock and Red Planet Diagenesis: Comparison of Earth and Mars Concretions PDF version, 3.4 MB : GSA Today, v. 15, n. 8, pp. 4–10. - Davis, J.M., 1999, Oriented carbonate concretions in a paleoaquifer: Insights into geologic controls on fluid flow: Water Resources Research, v. 35, p. 1705-1712. - Hattin, D.E., 1982, Stratigraphy and depositional environment of the Smoky Hill Chalk Member, Niobrara Chalk (Upper Cretaceous) of the type area, western Kansas: Kansas Geological Survey Bulletin 225:1-108. - Hobbs, D., and J. Hafnaer, 1999, Magnetism and magneto-structural effects in transition-metal sulphides: Journal of Physics: Condensed Matter, v. 11, pp. 8197–8222. - Hoffmann, V., H. Stanjek, and E. Murad, 1993, Mineralogical, magnetic and mössbauer data of symthite (Fe9S11) : Studia Geophysica et Geodaetica, v. 37, pp. 366–381. - Johnson, M.R., 1989, Paleogeographic significance of oriented calcareous concretions in the Triassic Katberg Formation, South Africa: Journal of Sedimentary Petrology, v. 59, p. 1008-1010. - Loope D.B., Kettler R.M., Weber K.A., 2011, Morphologic Clues to the origin of Iron Oxide-Cemented Sphereoids, Boxworks, and Pipelike Concretions, Navajo Sandstone of South-Central Utah, U.S.A, The Journal of Geology, Vol. 119, No. 5 (September 2011), pp. 505–520 - Loope D.B., Kettler R.M., Weber K.A., 2011, Follow the water: Connecting a CO2 reservoir and bleached sandstone to iron-rich concretions in the Navajo Sandstone of south-central Utah, USA, GEOLOGY FORUM, November 2011, Geological Society of America doi:10.1130/G32550Y.1 - McBride, E.F., M.D. Picard, and R.L. Folk, 1994, Oriented concretions, Ionian Coast, Italy: evidence of groundwater flow direction: Journal of Sedimentary Research, v. 64, p. 535-540. - McBride, E.F., M.D. Picard, and K.L. Milliken, 2003, Calcite-Cemented Concretions in Cretaceous Sandstone, Wyoming and Utah, U.S.A.: Journal of Sedimentary Research. v. 73, n. 3, p. 462-483. - Mozley, P.S., 1996, The internal structure of carbonate concretions: A critical evaluation of the concentric model of concretion growth: Sedimentary Geology: v. 103, p. 85-91. - Mozley, P.S., and Goodwin, L., 1995, Patterns of cementation along a Cenozoic normal fault: A record of paleoflow orientations: Geology: v. 23, p 539-542. - Mozley, P.S., and Burns, S.J., 1993, Oxygen and carbon isotopic composition of marine carbonate concretions: an overview: Journal of Sedimentary Petrology, v. 63, p. 73-83. - Mozley, P.S., and Davis, J.M., 2005, Internal structure and mode of growth of elongate calcite concretions: Evidence for small-scale microbially induced, chemical heterogeneity in groundwater: Geological Society of America Bulletin, v. 117, 1400-1412. - Pratt, B.R., 2001, "Septarian concretions: internal cracking caused by synsedimentary earthquakes": Sedimentology, v. 48, p. 189-213. - Raiswell, R., and Q.J. Fisher, 2000, Mudrock-hosted carbonate concretions: a review of growth mechanisms and their influence on chemical and isotopic composition: Journal of Geological Society of London. v. 157, p. 239-251 - Scotchman, I.C., 1991, The geochemistry of concretions from the Kimmeridge Clay Formation of southern and eastern England: Sedimentology. v. 38, pp. 79-106. - Thyne, G.D., and J.R. Boles, 1989, Isotopic evidence for origin of the Moeraki septarian concretions, New Zealand: Journal of Sedimentary Petrology. v. 59, n. 2, pp. 272-279. - Voigt, E., 1968, Uber-Hiatus-Konkretion (dargestellt an Beispielen aus dem Lias): Geologische Rundschau. v. 58, pp. 281–296. - Wilson, M.A., 1985, Disturbance and ecologic succession in an Upper Ordovician cobble-dwelling hardground fauna: Science. v. 228, pp. 575-577. - Wilson, M.A., and Taylor, P.D., 2001, Palaeoecology of hard substrate faunas from the Cretaceous Qahlah Formation of the Oman Mountains: Palaeontology. v. 44, pp. 21-41. - Zaton, M., 2010, Hiatus concretions: Geology Today. v. 26, pp. 186–189. |Wikimedia Commons has media related to Concretion.| - Dietrich, R.V., 2002, Carbonate Concretions--A Bibliography, The Wayback Machine. and PDF file of Carbonate Concretions--A Bibliography, CMU Online Digital Object Repository, Central Michigan University, Mount Pleasant, Michigan. - Biek, B., 2002, Concretions and Nodules in North Dakota North Dakota Geological Survey, Bismarck, North Dakota. - Epoch Times Staff, 2007, Mysterious Huge Stone Eggs Discovered in Hunan Province Epoch Times International. Photographs of large cannonball concretions recently found in Hunan Province, China. - Everhart, M., 2004, A Field Guide to Fossils of the Smoky Hill ChalkPart 5: Coprolites, Pearls, Fossilized Wood and other Remains Part of the Oceans of Kansas web site. - Hansen, M.C., 1994, Ohio Shale Concretions PDF version, 270 KB Ohio Division of Geological Survey GeoFacts n. 4, pp. 1–2. - Hanson, W.D., and J.M. Howard, 2005, Spherical Boulders in North-Central Arkansas PDF version, 2.8 MB Arkansas Geological Commission Miscellaneous Publication n. 22, pp. 1–23. - Heinrich, P.V., 2007, The Giant Concretions of Rock City Kansas PDF version, 836 KB BackBender's Gazette. vol. 38, no. 8, pp. 6–12. - Hokianga Tourism Association, nd, Koutu Boulders ANY ONE FOR A GAME OF BOWLS? and Koutu Boulders, Hokianga Harbour, Northland, New Zealand High-quality pictures of cannonball concretions. - Irna, 2006, All that nature can never do, part IV : stone spheres - Irna, 2007a, Stone balls : in France too! - Irna, 2007b, Stone balls in Slovakia, Czech Republic and Poland - Katz, B., 1998, Concretions Digital West Media, Inc. - Kuban, Glen J., 2006-2008. Nevada Shoe Prints? - McCollum, A., nd, Sand Concretions from Imperial Valley, a collection of articles maintained by an American artist. - Mozley, P.S., Concretions, bombs, and groundwater, on-line version of an overview paper originally published by the New Mexico Bureau of Geology and Mineral Resources. - United States Geological Survey, nd, Cannonball concretion - University of Utah, 2004, Earth Has 'Blueberries' Like Mars 'Moqui Marbles' Formed in Groundwater in Utah's National Parks press release about iron oxide and Martian concretions

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1 Answer | Add Yours Temperature is the measure of the flow of energy, particularly of heat. When you put two things of different temperature in contact with each other, they will eventually have the same temperature (thermal equilibrium). What happens is that heat flows from the hotter object to the coller object. For example, in your first beaker containing hot water, it will eventually cool to room temperature as heat flows from the hot water and beaker to the surroundings. At this point, you may notice that the beaker is also hot (it is probably in thermal equilibrium with the hot water). The cooling rate depends on the heat capacity of the substances. The lower the heat capacity, the quicker heat is gained or lost (depending on the situation). Now, if we want to know the difference in rate of cooling of a single beaker and a lot of beaker packed closely together, we simply have to note their heat capacity. Note that heat capacity depends on the mass of the substance. Assuming that the entire system (beaker + water) is a single entity of a certain heat capacity, the single beaker will cool faster than the group of beakers. [Of course this will change and will be different if the beakers are far apart from each other]. We’ve answered 302,729 questions. We can answer yours, too.Ask a question

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As hurricane Sandy made its way to the Eastern coast of the United States in October 2012, meteorologists called the storm unprecedented in terms of its potential for damage and fatalities. Few events on Earth rival the sheer power of a hurricane. Also known as tropical cyclones and typhoons, these fierce storms can churn the seas into a violent topography of 50-foot (15-meter) peaks and valleys, redefine coastlines and reduce whole cities to watery ruin. Some researchers even theorize that the dinosaurs were wiped out by prehistoric hypercanes, a kind of super-hurricane stirred to life by the heat of an asteroid strike [source: National Geographic]. Every year, the world experiences hurricane season. During this period, hundreds of storm systems spiral out from the tropical regions surrounding the equator, and between 40 and 50 of these storms intensify to hurricane levels. In the Northern Hemisphere, the season runs from June 1 to Nov. 30, while the Southern Hemisphere generally experiences hurricane activity from January to March. So 75 percent of the year, it's safe to say that someone somewhere is probably worrying about an impending hurricane. A hurricane builds energy as it moves across the ocean, sucking up warm, moist tropical air from the surface and dispensing cooler air aloft. Think of this as the storm breathing in and out. The hurricane escalates until this "breathing" is disrupted, like when the storm makes landfall. At this point, the storm quickly loses its momentum and power, but not without unleashing wind speeds as high as 185 mph (300 kph) on coastal areas. In this article, we'll explore the lifecycle and anatomy of a hurricane, as well as the methods we use to classify and track these ultimate storm systems as they hurtle across the globe. Defining a Hurricane To understand how a hurricane works, you have to understand the basic principles of atmospheric pressure. The gases that make up Earth's atmosphere are subject to the planet's gravity. In fact, the atmosphere weighs in at a combined 5.5 quadrillion tons (4.99 quadrillion metric tons). The gas molecules at the bottom, or those closest to the Earth's surface where we all live, are compressed by the weight of the air above them. The air closest to us is also the warmest, as the atmosphere is mostly heated by the land and the sea, not by the sun. To understand this principle, think of a person frying an egg on the sidewalk on a hot, sunny day. The heat absorbed by the pavement actually fries the egg, not the heat coming down from the sun. When air heats up, its molecules move farther apart, making it less dense. This air then rises to higher altitudes where air molecules are less compressed by gravity. When warm, low-pressure air rises, cool, high-pressure air seizes the opportunity to move in underneath it. This movement is called a pressure gradient force. These are some of the basic forces at work when a low-pressure center forms in the atmosphere -- a center that may turn into what people in the North Atlantic, North Pacific and Caribbean regions call a hurricane. What else is happening? Well, as we know, warm, moist air from the ocean's surface begins to rise rapidly. As it rises, its water vapor condenses to form storm clouds and droplets of rain. The condensation releases heat called latent heat of condensation. This latent heat warms the cool air, causing it to rise. This rising air is replaced by more warm, humid air from the ocean below. And the cycle continues, drawing more warm, moist air into the developing storm and moving heat from the surface to the atmosphere. This exchange of heat creates a pattern of wind that circulates around a center, like water going down a drain. But what about those signature ferocious winds? Converging winds at the surface are colliding and pushing warm, moist air upward. This rising air reinforces the air that's already ascending from the surface, so the circulation and wind speeds of the storm increase. In the meantime, strong winds blowing the same speed at higher altitudes (up to 30,000 feet or 9,000 meters) help to remove the rising hot air from the storm's center, maintaining a continual movement of warm air from the surface and keeping the storm organized. If the high-altitude winds don't blow at the same speed at all levels -- if wind shears are present -- the storm becomes disorganized and weakens. Even higher in the atmosphere (above 30,000 feet or 9,000 meters) high-pressure air over the storm's center also removes heat from the rising air, further driving the air cycle and the hurricane's growth. As high-pressure air is sucked into the low-pressure center of the storm, wind speeds increase. Then you have a hurricane to contend with. How a Hurricane Forms This content is not compatible on this device. Source: NASA Observatorium You never hear about hurricanes hitting Alaska. That's because hurricanes develop in warm, tropical regions where the water is at least 80 degrees Fahrenheit (27 degrees Celsius). The storms also require moist air and converging equatorial winds. Most Atlantic hurricanes begin off the west coast of Africa, starting as thunderstorms that move out over the warm, tropical ocean waters. A hurricane's low-pressure center of relative calm is called the eye. The area surrounding the eye is called the eye wall, where the storm's most violent winds occur. The bands of thunderstorms that circulate outward from the eye are called rain bands. These storms play a key role in the evaporation/condensation cycle that feeds the hurricane. The rotation of a hurricane is a product of the Coriolis force, a natural phenomenon that causes fluids and free-moving objects to veer to the right of their destination in the Northern Hemisphere and to the left in the Southern Hemisphere. Imagine flying a small plane directly south. While you're moving southward, the planet is rotating. If you plotted a flight from the North Pole to the equator on a map, the path will appear to curve to the right. So in the Northern Hemisphere, winds deflect to the right. In the Southern Hemisphere, they deflect to the left. This wind deflection gets storms spinning. As a result, hurricanes in the Northern Hemisphere rotate counterclockwise and clockwise in the Southern Hemisphere. The force also affects the actual path of the hurricane, bending them to the right (clockwise) in the Northern Hemisphere and to the left (counterclockwise) if you're south of the equator. If you can't remember, just move within five degrees of the equator; the Coriolis force is too weak there to help form hurricanes. Hurricanes often begin their lives as clusters of clouds and thunderstorms called tropical disturbances. These low-pressure areas feature weak pressure gradients and little or no rotation. Most of these disturbances die out, but a few persevere down the path to hurricane status. In these cases, the thunderstorms in the disturbance release latent heat, which warms areas in the disturbance. This causes the air density inside the disturbance to lower, dropping the surface pressure. Wind speeds increase as cooler air rushes underneath the rising warm air. As this wind is subject to the Coriolis force, the disturbance begins to rotate. The incoming winds bring in more moisture, which condenses to form more cloud activity and releases latent heat in the process. On the next page, we'll explore the brief, violent life of a hurricane. Lifecycle of a Hurricane Given the destruction the storm unleashes, it's easy to think of a hurricane as a kind of monster. It may not be a living organism, but it does require sustenance in the form of warm, moist air. And if a tropical disturbance continues to find enough of this "food" and to encounter optimal wind and pressure conditions, it will just keep growing. It can take anywhere from hours to days for a tropical disturbance to develop into a hurricane. But if the cycle of cyclonic activity continues and wind speeds increase, the tropical disturbance advances through three stages: - Tropical depression: wind speeds of less than 38 mph - Tropical storm: wind speeds of 39 to 73 mph - Hurricane: wind speeds greater than 74 mph Between 80 and 100 tropical storms develop each year around the world. Many of them die out before they can grow too strong, but around half of them eventually achieve hurricane status. Hurricanes vary widely in physical size. Some storms are compact, with only a few bands of wind and rain trailing behind them. Other storms are looser -- the bands of wind and rain spread out over hundreds or thousands of miles. Hurricane Floyd, which hit the eastern United States in September 1999, was felt from the Caribbean islands to New England. Once a hurricane has formed and intensified, the only remaining path for the atmospheric juggernaut is dissipation. Eventually, the storm will encounter conditions that deny it the warm, moist air it requires. When a hurricane moves onto cooler waters at a higher latitude, gradient pressure decreases, winds slow, and the entire storm is tamed, from a tropical cyclone to a weaker extratropical cyclone that peters out in days. That important supply of warm, moist air also vanishes when the hurricane makes landfall. Condensation and the release of latent heat diminishes, and the friction of an uneven landscape decreases wind speeds. This causes winds to move more directly into the eye of the storm, eliminating the large pressure difference that fuels the storm's awesome power. Hurricanes can unleash incredible damage when they hit. With enough advance warning though, cities and coastal areas can give residents the time they need to fortify the area and even evacuate. To better classify each hurricane and prepare those affected for the intensity of the storm, meteorologists rely on rating systems. Australian meteorologists use a slightly different scale to classify hurricanes. While the Australian scale of cyclone intensity also ranks storms by wind speed and damage on a scale of 1 to 5, it covers both hurricanes and tropical storms. On the next page, we'll look at the tremendous damage hurricanes can inflict when they collide with coastal areas. Over the course of millennia, hurricanes have cemented their reputation as destroyers. Many people even frame them as the embodiment of nature's power or acts of divine wrath. The word "hurricane" itself actually derives from "Hurakan," a destructive Mayan god. No matter how you choose to sum up or personify these powerful acts of nature, the damage they inflict stems from several different aspects of the storm. Hurricanes deliver massive downpours of rain. A particularly large storm can dump dozens of inches of rain in just a day or two, much of it inland. That amount of rain can create flooding, potentially devastating large areas in the path of the hurricane's fierce center. In addition, high sustained winds within the storm can cause widespread structural damage to both man-made and natural structures. These winds can roll over vehicles, collapse walls and blow over trees. The prevailing winds of a hurricane push a wall of water, called a storm surge, in front of it. If the storm surge happens to coincide with high tide, it causes beach erosion and significant inland flooding. The hurricane itself is often just the beginning. The storm's winds often spawn tornadoes, which are smaller, more intense cyclonic storms that cause additional damage. You can read more about them in How Tornadoes Work. The extent of hurricane damage doesn't just depend on the strength of the storm, but also the way it makes contact with the land. In many cases, the storm merely grazes the coastline, sparing the shores its full power. Hurricane damage also greatly depends on whether the left or right side of a hurricane strikes a given area. The right side of a hurricane packs more punch because the wind speed and the hurricane's speed of motion complement one another there. On the left side, the hurricane's speed of motion subtracts from the wind speed. This combination of winds, rain and flooding can level a coastal town and cause significant damage to cities far from the coast. In 1996, Hurricane Fran swept 150 miles (241 km) inland to hit Raleigh, N.C. Tens of thousands of homes were damaged or destroyed, millions of trees fell, power was out for weeks in some areas and the total damage was measured in the billions of dollars. Tracking a Hurricane To monitor and track the development and movement of a hurricane, meteorologists rely on remote sensing by satellites, as well as data gathered by specially equipped aircraft. On the ground, Regional Specialized Meteorological Centers, a network of global centers designated by the World Meteorological Organization, are charged with tracking and notifying the public about extreme weather. Weather satellites use different sensors to gather different types of information about hurricanes. They track visible clouds and air circulation patterns, while radar measures rain, wind speeds and precipitation. Infrared sensors also detect vital temperature differences within the storm, as well as cloud heights. The Hurricane Hunters are members of the 53rd Weather Reconnaissance Squadron/403rd Wing, based at Keesler Air Force Base in Biloxi, Miss. Since 1965, the Hurricane Hunters team has used the C-130 Hercules, a very sturdy turboprop plane to fly into tropical storms and hurricanes. The only difference between this plane and the cargo version is the specialized, highly sensitive weather equipment installed on the WC-130. The team can cover up to five storm missions per day, anywhere from the mid-Atlantic to Hawaii. The Hurricane Hunters gather information about wind speeds, rainfall and barometric pressures within the storm. They then relay this information back to the National Hurricane Center in Miami, Fla. If you're curious about these foolhardy pilots, read Why would someone fly an airplane into a hurricane? Meteorologists take all the storm data they receive and use it to create computer forecast models. Based on a great deal of current and past statistical data, these virtual storms allow scientists to forecast a hurricane's path and changes in intensity well in advance of landfall. With this data, governments and news agencies ideally can warn residents of coastal areas and greatly reduce the loss of life during a hurricane. Long-term forecasting now allows meteorologists to predict how many hurricanes will take place in an upcoming season and to study trends and patterns in global climate. While personifying a massive, destructive force certainly makes for a jazzier headline, the practice of naming hurricanes originated with meteorologists, not media outlets. Often more than one tropical storm is active at the same time, so what better way to tell them apart than by naming them? For several hundred years, residents of the West Indies often named hurricanes after the Catholic saint's day on which the storm made landfall. If a storm arrived on the anniversary of a previous storm, a number was assigned. For example, Hurricane San Felipe struck Puerto Rico on Sept. 13, 1876. Another storm struck Puerto Rico on the same day in 1928, so this storm was named Hurricane San Felipe the Second. During World War II, weather officials only gave hurricanes masculine names. These names closely followed radio code names for letters of the alphabet. This system, like the West Indian saints system, drew from a limited naming pool. In the early 1950s, weather services began naming storms alphabetically and with only feminine names. By the late 1970s, this practice was replaced with the equal opportunity system of alternating masculine and feminine names. The World Meteorological Organization (WMO) continues this practice to this day. The first hurricane of the season is given a name starting with the letter A, the second with the letter B and so on. As the storms affect varying portions of the globe, the naming lists draw from different cultures and nationalities. Hurricanes in the Pacific Ocean are assigned a different set of names than Atlantic storms. For example, the first hurricane of the 2001 hurricane season was a Pacific Ocean storm near Acapulco, Mexico, named Adolf. The first Atlantic storm of the 2001 season was named Allison. A list of names through 2011 is available from the National Hurricane Center. If a hurricane inflicts significant damage, a country affected by the storm can request that the name of the hurricane be "retired" by the WMO. A retired name can't be reissued to a tropical storm for at least 10 years. This helps to avoid public confusion and to simplify both historical and legal record keeping. Our modern understanding of hurricanes depends largely on a mere century's worth of scientific study and record keeping, but the storms have been dictating the course of human history for millennia. After all, they're a part of an atmospheric system that predates the human race by billions of years. While scientists are largely left to speculate about the strength of Mesozoic Era storms, geologists have discovered evidence of Iron Age hurricanes in layers of ground sediment. When storm surges wash over land and into lakes, they leave fans of sand behind. Scientists can carbon date organic materials above and below the sand to determine an approximate storm date. A team from Louisiana State University studied thousands of years worth of lake bed evidence and discovered that, over the past 3,400 years, a dozen Category 4 or higher hurricanes hit the area -- yet most of them occurred 1,000 years or more ago [source: Young]. Findings such as these allow scientists to better study long-term weather patterns and possibly make better sense of current climate trends. As far as human records go, the ancient Mayans of South America made some of the earliest mentions of hurricanes in their hieroglyphics. The centuries that follow are littered with accounts of hurricanes affecting the outcomes of wars, colonization efforts and an untold number of personal lives. Just to name a few, hurricane activity thwarted the following sea ventures through the destruction and scattering of ocean fleets: - The 1274 Mongol invasion of Japan - A 1559 attempt by the Spanish to recapture Florida - The French defense of a Floridian fort, subsequently lost to the Spanish in 1565 - The Spanish Armada's attack on England in 1588 - A 1640 Dutch attack on Cuba - British dominance over the French in the Caribbean Islands in 1780 Today, modern meteorology prevents most hurricanes from arriving unannounced, greatly decreasing the massive hurricane fatality rates of previous centuries. But even with advance warning, governments and the residents of coastal areas still have to properly prepare for the coming storms. Meanwhile, some experts look to the future with concern. Some point to periods of intense hurricane activity in Earth's past and worry that such trends may return. Others argue that global warming brought on by the increased production of greenhouse gasses will lead to larger hurricane zones and more powerful storms. After all, hurricanes thrive on warm, moist waters, and a warmer Earth could provide more sustenance for tropical storms. Explore the links on the next page to learn more about hurricanes and the Earth's weather, including a story about those crazy pilots who fly their planes into hurricanes. Should you get in your bathtub during a tornado? Read on to find out why — and why not. More Great Links - "Atmosphere." Britannica Student Encyclopædia. 2008. (Aug. 5, 2008)http://student.britannica.com/comptons/article-196868/atmosphere - Drye, Willie. "Hurricanes of History -- From Dinosaur Times to Today." National Geographic News. Jan. 28, 2005. (Aug. 19, 2008) http://news.nationalgeographic.com/news/2005/01/0128_050128_tv_hurricane.html - "Evolution of the atmosphere." Britannica Online Encyclopædia. 2008. (Aug. 8, 2008)http://www.britannica.com/EBchecked/topic/1424734/evolution-of-the-atmosphere - "The History of Hurricanes." Federal Emergency Management Agency. (Aug. 21, 2008)http://www.fema.gov/kids/hurr_hist.htm - "Hurricane Timeline: 1495 to 1800." South Florida Sun-Sentinel. 2008. (Aug. 21, 2008)http://www.sun-sentinel.com/news/weather/hurricane/sfl-hc-history-1495to1800,0,3354030.htmlstory - "Jet stream." Britannica Online Encyclopædia. 2008. (Aug. 8, 2008)http://www.britannica.com/EBchecked/topic/303269/jet-stream - "Lightning." Britannica Online Encyclopædia. 2008. (Aug. 8, 2008)http://www.britannica.com/EBchecked/topic/340767/lightning#default - Reynolds, Ross. "Cambridge Guide to Weather." Cambridge University Press. 2000. - Tarbuck, Edward and Frederick Lutgens. "Earth Science: Eleventh Edition." Pearson Prentice Hall. 2006. - Toothman, Jessika. "How Clouds Work." HowStuffWorks.com. May 5, 2008. (Aug. 8, 2008)http://science.howstuffworks.com/cloud.htm - "Tropical Cyclone." Britannica Online Encyclopædia. 2008. (Aug. 20, 2008)http://www.britannica.com/EBchecked/topic/606551/tropical-cyclone - Vogt, Gregory L. "The Atmosphere: Planetary Heat Engine." Twenty-First Century Books. 2007. - Wilson, Tracy V. "How the Earth Works." HowStuffWorks.com. April 21, 2006. (Aug. 8, 2008)http://science.howstuffworks.com/Earth.htm - Young, Emma. "Raiders of the lost storms." New Scientist. June 10, 2006. (Aug. 21, 2008)http://environment.newscientist.com/channel/earth/hurricane-season/mg19025551.300-raiders-of-the-lost-storms.html

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There are many different types of particles, with different particle sizes and properties. Atoms and molecules are called microscopic particles. Subatomic particles are particles that are smaller than atoms. The proton, the neutron, and the electron are subatomic particles. These are the particles which make atoms. The proton has a positive charge (a + charge). The neutron has a neutral charge. The electron has a negative charge (a - charge), and it is the smallest of these three particles. In atoms, there is a small nucleus in the center, which is where the protons and neutrons are, and electrons orbit the nucleus. Protons and neutrons are made up of quarks. Quarks are subatomic particles, but they are also elementary particles because we do not know if they are made up of even smaller particles. There are six different types of quarks. These are the up quark, the down quark, the strange quark, the charm quark, the bottom quark, and the top quark. A neutron is made of two down quarks and one up quark. The proton is made up of two up quarks and one down quark.

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| || | | "If we look we'll find 'em... the microbes are there. They're these little packages of secrets that are waiting to be opened." | - Anna-Louise Reysenbach Introduction Microbes flourish. Inside your gut, in the mucky soil of a marsh, in Antarctic ice, in the hot springs of Yellowstone, in habitats seemingly incompatible with life, microbes flourish. They were present on Earth 3.5 to 4 billion years ago, and they've been evolving and expanding into new environments ever since. Replicating quickly, exchanging genetic material with each other and with other organisms, bacteria and archaea have become ubiquitous. Not only are they everywhere, but these tiny organisms also manipulate the environments in which they live. Their presence has driven the development of new ecosystems - some of which allowed for the evolution of more complex organisms. Without microbes, the recycling of essential nutrients on Earth would halt. Microbes communicate; some generate the signals for the formation of metabolically diverse communities. Some use sophisticated signaling to establish complex relationships with higher organisms. In this unit we will examine examples of the broad diversity of microorganisms and consider their roles in various ecosystems, both natural and man-made. We will also discuss some of the practical applications that derive from the wealth of metabolic diversity that microorganisms possess. Let's start at the beginning ... three or four billion years ago. © Annenberg Foundation 2016. All rights reserved. Legal Policy

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After reading this tutorial you might want to check out some of our other Mathematics Quizzes as well. In the first two tutorials, we have considered a random experiment that has only one characteristic and hence its outcome is a random variable X that assumes a single value. However, in the following tutorial, we will deal with random experiments having 2 (or more) characteristics and hence random variables X, Y (or more). Such random variables are called jointly distributed rvs. x1: height of a person x2: weight of a person x3: blood pressure of a person x4: sugar count of a personHence, x1 , x2, x3, x4 are jointly distributed. However, here we will consider a two dimensional random variable (X,Y). We will study the following cases: Also we shall study some other characteristics of jointly distributed random variables and transformation of random vectors. (For Discrete Variables X and Y) : Joint Probability, Probability mass function, Marginal Probability mass function, Conditional Probability Mass Function, Independence of events: (For Continuous Variables X and Y) : Probability density function, Marginal Probability Distribution Function, Conditional Probability Distribution Function, Independence of events: Properties common to both cases: Properties of CDF, Product Moments, Central moments, Non Central moments. Q:Find the relation between Geometric and Pascal distributions, Exponential and Gamma distributions. Q:Suppose a shopkeeper has 10 pens of a brand out of which 5 are good(G), 2 have defective inks(DI) and 3 have defective caps(DC). If 2 pens are selected at random, find the probability i. Not more than one is DI and not more than one is DC. ii. P(DI<2) Q:The joint probability mass function of (X, Y) is given by p(x, y) = k (4x + 4y), x = 1, 2, 3; y = 0, 1, 2. Find the (i) marginal distributions of Y (ii)P(X ≤ 2 | Y ≤ 1). Q:Check whether X and Y are independent: P(X=1, Y=1) = 1/4, P(X=1, Y=0) = 1/4, P(X=0, Y=1) = 1/4, P(X=0, Y=0) = 1/4 Q:A shopping mall has parking facility for both 2-wheelers and 4-wheelers. On a randomly selected day, let X and Y be the proportion of 2 and 4 wheelers respectively. The Joint pdf of X and Y are: f ( x, y ) = ( x + 2y ) * 2/3 ; 0 ≤ x ≤ 1; 0 ≤ y ≤ 1 =0 elsewhere i. Find the marginal densities of X and Y. ii. Find the probability that the proportion of two wheelers is less than half. Q:Prove the additive property of: Binomial distribution, Poisson distribution. Q:The amount of rainfall recorded in Jalna in June is a rv X and the amount in July is a rv Y. X and Y have a bivariate normal distribution. (X, Y) ~ (6, 4, 1, 0.25, 0.1) Find: (i) P(X ≤ 5) (ii) P(Y ≤ 5| X = 5) Q:Let X1,.....Xn be i.i.d with cdf F(x) and pdf f(x). Find the distribution of min and max of X. Complete Tutorial with Problems and Solutions :

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- 1. In your own words, define inventing. Then do the following: - a. Explain to your merit badge counselor the role of inventors and their inventions in the economic development of the United States. - b. List three inventions and state how they have helped humankind. - 2. Do ONE of the following: - a. Identify and interview with a buddy (and with your parent’s permission and merit badge counselor’s approval) an individual in your community who has invented a useful item. Report what you learned to your counselor. - b. Read about three inventors. Select the one you find most interesting and tell your counselor what you learned. - 3. Do EACH of the following: - a. Define the term intellectual property. Explain which government agencies oversee the protection of intellectual property, the types of intellectual property that can be protected, how such property is protected, and why protection is necessary. - b. Explain the components of a patent and the different types of patents available. - c. Examine your Scouting gear and find a patent number on a camping item you have used. With your parent’s permission, use the Internet to find out more about that patent. Compare the finished item with the claims and drawings in the patent. Report what you learned to your counselor. - d. Explain to your counselor the term patent infringement. - 4. Discuss with your counselor the types of inventions that are appropriate to share with others, and explain why. Tell your counselor about one nonpatented or noncopyrighted invention and its impact on society. - 5. Choose a commercially available product that you have used on an overnight camping trip with your troop. Make recommendations for improving the product, and make a sketch that shows your recommendations. Discuss your recommendations with your counselor. - 6. Think of an item you would like to invent that would solve a problem for your family, troop, chartered organization, community, or a special-interest group. Then do EACH of the following, while keeping a notebook to record your progress: - a. Talk to potential users of your invention and determine their needs. Then, based on what you have learned, write a statement describing the invention and how it would help solve a problem. This statement should include detailed sketch of the invention. - b. Create a model of the invention using clay, cardboard, or any other readily available material. List the materials necessary to build a prototype of the invention. - c. Share the idea and the model with your counselor and potential users of your invention. Record their feedback in your notebook. - 7. Build a working prototype of the item you invented for requirement 6*. Test and evaluate the invention. Among the aspects to consider in your evaluation are cost, usefulness, marketability, appearance, and function. Describe how your original vision and expectations for your invention are similar or dissimilar to the prototype you built. Have your counselor evaluate and critique your prototype. - Before you begin building the prototype, you must have your counselor’s approval, based on the design and building plans you have already shared. - 8. Do ONE of the following: - a. Participate with a club or team (robotics team, science club, or engineering club) that builds a useful item. Share your experience with your counselor. - b. Visit a museum or exhibit dedicated to an inventor or invention, and create a presentation of your visit to share with a group such as your troop or patrol. - 9. Discuss with your counselor the diverse skills, education, training, and experience it takes to be an inventor. Discuss how you can prepare yourself to be creative and inventive to solve problems at home, in school, and in your community. Discuss three career fields that might utilize the skills of an inventor. | The official source for the information shown in this article or section is:| scoutingmagazine.org, 2015 Edition (BSA Supply SKU #620714)

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In programming, a series of objects all of which are the same size and type. Each object in an array is called an array element. For example, you could have an array of integers or an array of characters or an array of anything that has a defined data type. The important characteristics of an array are: Stay up to date on the latest developments in Internet terminology with a free weekly newsletter from Webopedia. Join to subscribe now. Webopedia's student apps roundup will help you to better organize your class schedule and stay on top of assignments and homework. Read More »20 Ways to Shorten a URL If you need to shorten a long URL try this list of 20 free online redirection services. Read More »Top 10 Tech Terms of 2015 The most popular Webopedia definitions of 2015. Read More » This Webopedia study guide describes the different parts of a computer system and their relations. Read More »Network Fundamentals Study Guide Networking fundamentals teaches the building blocks of modern network design. Learn different types of networks, concepts, architecture and... Read More »The Five Generations of Computers Learn about each of the five generations of computers and major technology developments that have led to the current devices that we use today. Read More »

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Perform an Indian folk song with the correct nuance (degrees of differences). Work cooperatively in a Describe the characteristics of Indian vocal music. Music is an essential ingredient in the lives of people in India. Folk and classical music accompany various styles of Indian dances. For Indians, music is a means to get touch with the Supreme Being. A song is an apparatus for the communication and an interaction process between the worshipper and deity. Through the clear expression of the rasa or dynamics, the singer feels the presence of Brahma, the creator. Vocal music in India is a way to express deep devotion to God. It is manifested through the art of vocalization which becomes more then just a vocal warm up but an act of Most of the Indian classical songs of Northern India are devotional but few are The most notable is the dhun or kirtan for the Hindus, bhajan, the shabad for Sikho and the Kawali (qawali) for the Muslims. Not all Indian music are serious. Gangal is one style which is known for its rich romantic and poetic content. The lakshan geet is a style which is oriented towards musical Swarmalika is used for pedagogic purposes. In style, sargam is used instead of words. Hymn to Shiva It is an example of an Indian song. There is an English translation of the The notations are easy to sing. The flat sign is used before some of the notes in order to sing with the correct It has four measures to the base clef. It has four beats. Sing the song and add the drone accompaniments vocally or on the Percussion instruments can also be used as an accompaniment to add color to the singing. Some Indian songs are used to describe the scenic beauty of a particular region in the country. This song Tamil Nad speaks of the beauty of the Land of the Tamils. Indian songs are also used to bid a person farewell. This song Vijaya means Deva – Dasi Dance This is another Indian composition which is highly rhythmical. The flow to the tone and rhythm is Clapping can be used in case there are no drums available. The notations of the songs from India are classified into two forms: The text of the songs is all about A particular slide catching your eye? Clipping is a handy way to collect important slides you want to go back to later.

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Missouri Compromise 1820<br />1819 Missouri applied to join the Union as a slave state.<br />This gave the South a majority in the Senate.<br />Henry Clay proposed admitting Maine as a Free state to maintain the balance in Congress.<br />Future states from the LouisianaPurchase would be Free above the 3630’ N. Latitude and states below that line would be Slave.<br /> Slavery in the West<br />David Wilmot, Congressman from PA, submitted a bill, the WilmotProviso that would ban slavery in any of the territories gained from Mexico.<br />The bill passed in the House, but NOT in the Senate, leaving the question of slavery in the west unresolved.<br /> Opposing Views <br />Believed in extending the Missouri Compromise line or Popular Sovereignty<br />Believed slavery should not be restricted and slaves should be returned to their owners<br />Wanted to ban slavery throughout the entire country<br />Abolitionists<br />Southerners<br />Moderates<br /> Opposing Views<br />Wanted to ban slavery in all parts of the country<br />Abolitionists<br />Believed in extending the Missouri Compromise or Popular Sovereignty<br />Moderates<br />Believed that slavery should be allowed everywhere and runaways should be returned to their owners<br />Southerners<br /> Free-Soil Party 1848<br /><ul><li>Main goal was to keep slavery from spreading to the western territories. They did not look to ban slavery where it already existed. This was the first election where slavery was an important issue.</li></li></ul><li> Compromise of 1850<br />Chapter 16 Section 2<br /> Slavery Debate Erupts Again<br />California applies for statehood as a Free state in 1850.<br />Southerners feared that they would be out voted in the Senate and it was suggested that they should secede from the Union.<br />Like many northerners, Webster viewed slavery as evil. The breakup of the United States, however, he believed was worse. To save the Union, Webster was willing to compromise. He would support southern demands that northerners be forced to return fugitive slaves.<br /> John C. Calhoun refused to compromise insisting that fugitive slaves be returned to their owners.<br />Henry Clay feared that if a compromise was not reached the country would break apart.<br /> Compromise of 1850<br />Calhoun died and Clay became ill as Congress still debated the slavery issue.<br />Stephen Douglas of IL took up Clay’s fight for compromise to pass legislation to satisfy both North and South.<br />The compromise consisted of 5 separate components.<br /> Compromise of 1850<br />California is admitted as a free state.<br />Territories of New Mexico and Utah would uphold popular sovereignty.<br />Bans slave trade in Washington, D.C.<br />Settled the dispute over the Texas/New Mexico Border.<br />Passed the Fugitive Slave Act. <br />What does this mean for the Missouri Compromise?<br /> Fugitive Slave Act<br />Required all citizens to return fugitive, runaway, slaves to their owners.<br />Anyone who helped or allowed fugitives to escape would be fined $1,000 (equal to $25,480 today)<br />African Americans suspected to be a runaway was not allowed a trial by jury.<br />Judges were paid $10.00 ($250.00) for charging blacks as runaways and returned to the south. $5.00 ($125.00) for deciding they were free. <br /> Uncle Tom’s Cabin<br />A book written by Harriet Beecher Stowe in 1852.<br />Harriet Beecher Stowe lived along the Ohio River where many slaves crossed to get to freedom.<br />Book was fictitious, but based on the stories she heard from escaped slaves.<br />Gave people in the North a better understanding of what it meant to be a slave and saw slavery as a moral problem<br /> The Crisis Deepens<br />Chapter 16 Section 3<br /> Kansas Nebraska Act 1854<br />Compromise of 1850 nullified the Missouri Compromise, but only clarified how the slavery issue would be handled in the Mexican Cession.<br />So what about the Kansas and Nebraska territories?<br />Stephen Douglas proposed that both territories’ settlers decide whether slavery would be allowed in those territories upon applying for statehood. This is called…<br />Popular Sovereignty<br /> Predictions<br />What was the reaction to the Kansas Nebraska Act in the North?<br />“Opponents of slavery called the act a ‘criminal betrayal of precious rights.’ Slavery could now spread to areas that had been free for more than 30 years. Some northerners protested by openly challenging the Fugitive Slave Act.”<br />Do you think popular sovereignty will solve the issue of slavery?<br /> Crisis Turns Violent…<br />Initial settlers in the Kansas territory were from the neighboring states for the purpose of acquiring cheap land.<br />Few of these settlers owned slaves.<br />Under the Kansas Nebraska Act was the territory of Kansas going to enter the Union as a Free state or a Slave state?<br />It would be decided using Popular Sovereignty<br /> Crisis Turns Violent <br />To increase the number of slave owners in the Kansas territory Border Ruffians, proslavery settlers from Missouri rode across the border into Kansas to increase the number of slave owners in the territory.<br />These slave owners voted in the government elections and illegally voted in a proslavery government.<br />The original non-slave owners who initially settled the territory refused to obey the proslavery government and elected their own legislature.<br /> Bleeding Kansas<br />A proslavery band of men decided to attack a antislavery town of Lawrence, KS destroying homes and a Free-Soil Newspaper.<br />In retaliation, John Brown, an abolitionist, with his four sons, attacked the proslavery town Pottawatomie(paht uh waht uh mee) Creek.<br />In the middle of the night he dragged five proslavery settlers from their beds and murdered them.<br />This created both sides to use hit-and-run tactics, guerilla warfare, on the other, killing over 200 people.<br /> Violence in the Senate<br />Senator, Charles Sumner of MA, criticized the proslavery government of Kansas and verbally attacked proslavery southerners, specifically Andrew Butler.<br />Due to his age, his nephew, Congressmen Preston Brooks felt he couldn’t defend himself, so he marched on the Senate floor and beat Sumner over the head with a cane till he was bleeding and unconscious.<br /> Dred Scott v. Sanford <br />1857 antislavery lawyers submitted a lawsuit, alegalcase to settle a dispute, on behalf of a slave Dred Scott who’s owner had died.<br />His lawyers argued that, because his lawyer moved him to reside in IL and WI, both free states, he should be set free.<br /> Dred Scott Decision 1857<br />The Supreme Court decided: <br />1) that DredScott was property, therefore not a citizen he was incapable to filing a lawsuit to begin with. <br />2) According to the constitution, no citizen can be deprived of property thus, Congress did not have the power to outlaw slavery in any territory.<br /> “That the history of the nation during the last four years has fully established the propriety and necessity of the organization and perpetuation of the Republican Party and that the causes which called it into existence are permanent in their nature and now, more than ever before, demand its peaceful and constitutional triumph.”<br />The Republican Party Emerges<br />Chapter 16 Section 4<br /> Republican Party<br />Neither of the major political parties, Whigs or Democrats, would take a stand on the issue of slavery.<br />In 1854, Free Soilers, Northern Democrats, antislavery Whigs formed the Republican Party.<br />Their major goal was to keep the spread of slavery in the West. <br /> Abraham Lincoln, Republican<br />Lincoln entered the national political scene during his debates with Stephen Douglas in 1858 for the Illinois Senate.<br />During the series of debates Douglas supported popular sovereignty, Lincoln argued that slavery should not be allowed in the territories because the “House divided against itself could not stand.”<br /> John Brown Raid<br /><ul><li>Brown led a group of men to a federal arsenal, gun warehouse, at Harper’s Ferry, VA. He believed that once weapons were available slaves would join him and revolt against their owners. No revolt took place and he was arrested by troops commanded by Robert E. Lee.</li></li></ul><li>Hero or Villain?<br />Brown was found guilty of murder and treason, actions against one’s country.<br />Because he acted with such dignity through out his trial and his head held high many northerners considered him a martyr, someone who is willing to give their life for a cause.<br />Many southerners became convinced that the North wanted to destroy slavery. <br />Why?<br /> The Nation Divides<br />Chapter 16 Section 5<br /> Election of 1860<br />Setting the Scene: <br />Republican Convention Chicago, IL<br />“Fire the salute,” ordered the delegate. “Old Abe is nominated!” Amid the celebration, though, a delegate from Kentucky struck a somber note. “Gentlemen, we are on the brink of a great civil war.”<br /> Election of 1860<br />Abraham Lincoln – Republican<br />Prevent the spread of slavery in the western territories<br />Stephen Douglas – Northern Democrat<br />Refused to support slavery<br />John Breckinridge – Southern Democrat<br />Supported spread of slavery<br />John Bell – Constitutional Union <br />Moderate Southerner who wanted to keep the Union together<br /> Lincoln was able to gain a majority vote without even being listed on 10 of the Southern ballots.<br /> Southern Reaction<br />South believed that when Lincoln took office he would abolish slavery.<br />The South no longer had a voice in the federal government and congress, as well as the President, was against their interests – slavery.<br />Governor of SC, Francis W. Pickens, wrote to other southern states that it was their duty to secede from the Union.<br />SC seceded from the Union December 20, 1860.<br /> Confederate States of America<br />By February 1861 the following states made up the Confederacy:<br />South Carolina (first to secede)<br />Alabama, Florida, Georgia, Louisiana, Mississippi, and Texas .<br />At a convention held in Montgomery, AL, Jefferson Davis was appointed their president.<br />Davis served in Mexican War, as senator from MS, supporter of state’s rights, Secretary of War under President Pierce.<br /> The Right to Secede?<br />Most southerners believed that they had every right to secede. After all, the Declaration of Independence said that “it is the right of the people to alter or to abolish” a government that denies the rights of its citizens. <br />Lincoln, they believed, would deny white southerners the right to own slaves.<br /> Civil War Begins<br />Lincoln took the oath of office on March 4, 1861.<br />Lincoln’s First Inaugural Address.<br />April 1861, Fort Sumter <br />Federal fort off the coast of South Carolina is in need of food and supplies.<br />Lincoln informs the governor that he is shipping food and not weapons or ammunition<br />As part of the Confederacy, SC could not allow the Union to have control within its borders.<br />Confederate soldiers bombarded the fort with shells forcing the Union to surrender on April 13, 1861.<br /> A particular slide catching your eye? Clipping is a handy way to collect important slides you want to go back to later.

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Christopher ColumbusResearch Paper Christopher Columbus and over other 25,000+ free term papers, essays and research papers examples are available on the website! Autor: Sretieff12345 • October 11, 2012 • 943 Words (4 Pages) • 636 Views Christopher Columbus', voyage during 1492-1502 altered the course of European history as we know it today. During his voyages, throughout the West Indies, "Christopher Columbus paved the way for others to conquer and settle the new land in the name of the Spanish crown". Although at least "ten different powers would eventually play a role in settling the Caribbean", the British, French, and the Dutch would compete mostly with the Spanish. The Dutch were interested in trading, and would eventually introduce the British and French to the plantation systems of sugarcane production. The British, and French were interested in colonization and would later have their colonies forced to trade only with their mother countries. To portray Columbus' influence on European history I will concentrate on the influences he had on the Dutch's wants and needs compared to that of the British and the French. Columbus represented a culture that was expanding its power. European countries were exploring to gain more access to natural resources in different parts of world in order to increase their authority in comparison to the competitors. The riches of this new world attracted other European powers. The British, Dutch and French challenged Spain's monopoly in the 17th century. Columbus and his people used "piracy, smuggling, and outright war to take over lands and set up their own colonies". It was the Dutch, for example, that captured "Guiana, and the British captured St. Kitts, Barbados and Jamaica from Spain". First, the Dutch were accredited with the cultivation of sugarcane in their early Brazilian colonies. As a result of producing sugarcane, they set up trading centers on the few small islands they had settled in. During the Dutch period 1570 to 1678 "the Dutch shipping industry became the main provider of supplies and slaves for the other Caribbean colonies" and became the primary resource for sugar. Today, the "only reminders of Dutch activities in the Caribbean are Suriname and 6 small island possessions in the Lesser Antilles". In 1640, Close to the end of the Dutch period, the Dutch introduced British and French colonialists to the production of sugarcane. By the middle of the "17th century, some British pirates settled among logwood forests on the coast of the Bay of Honduras, which later became the Settlement of Belize. The French were also settling in North America and the Caribbean". Within Europe the "increasing market for sugar ensured the colonies an early success", and opened a new way of conducting business with their mother countries. During the 18th century, the French and British fought for domination over the "New World". "The British took control of more and more territories in the Caribbean and by the 19th century were the major power in the Caribbean. Through this era, the "British

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|Languages:||English, Munsee, and Unami| |Religions:||Christianity, Native American Church,| traditional tribal religion |Related:||Other Algonquian peoples| The Lenape are a Native American tribe and First Nations band government. They are also called Delaware Indians and their historical territory was along the Delaware River watershed, western Long Island and the Lower Hudson Valley. Most Lenape were pushed out of their Delaware homeland during the 18th century by expanding European colonies, exacerbated by losses from intertribal conflicts. Lenape communities were weakened by newly introduced diseases, mainly smallpox, and violent conflict with Europeans. Iroquois people occasionally fought the Lenape. Surviving Lenape moved west into the upper Ohio River basin. The American Revolutionary War and United States' independence pushed them further west. In the 1860s, the United States government sent most Lenape remaining in the eastern United States to the Indian Territory (present-day Oklahoma and surrounding territory) under the Indian removal policy. In the 21st century, most Lenape now reside in the US state of Oklahoma, with some communities living also in Wisconsin, Ontario (Canada) and in their traditional homelands. Lenape kinship system has matrilineal clans, that is, children belong to their mother's clan, from which they gain social status and identity. The mother's eldest brother was more significant as a mentor to the male children than was their father, who was of another clan. Hereditary leadership passed through the maternal line, and women elders could remove leaders of whom they disapproved. Agricultural land was managed by women and allotted according to the subsistence needs of their extended families. Families were matrilocal; newlywed couples would live with the bride's family, where her mother and sisters could also assist her with her growing family. Lenni-Lenape (or Lenni-Lenapi) comes from their autonym, Lenni, which may mean "genuine, pure, real, original," and Lenape, meaning "Indian" or "man". (cf. Anishinaabe.) Alternately, lënu may be translated as "man." The Lenape, when first encountered by whites, were a loose association of related peoples who spoke similar languages and shared familial bonds in an area known as Lenapehoking, the Lenape traditional territory, which spanned what is now eastern Pennsylvania, New Jersey, southern New York, and eastern Delaware. The tribe's other name, "Delaware," is not of Native American origin. English colonists named the Delaware River for the first governor of Virginia, Thomas West, 3rd Baron De La Warr, whose title was ultimately derived from French. (For etymology of the surname, see Earl De La Warr§Etymology.) The English then began to call the Lenape the Delaware Indians because of where they lived. Swedes also settled in the area, and early Swedish sources listed the Lenape as the Renappi. See main article: Lenapehoking. Traditional Lenape lands, the Lenapehoking, encompassed the Delaware Valley of eastern Pennsylvania and western New Jersey from the Lehigh River south into eastern Delaware and the Delaware Bay, western Long Island, New York Bay, and the Lower Hudson Valley in New York. The Lenape lived in numerous small towns along the rivers and streams that fed the waterways. The Unami and Munsee languages belong to the Eastern Algonquian language group. Although the Unami and Munsee speakers people are related, they consider themselves as distinct, as they used different words and lived on opposite sides of the Kitatinny Mountains of modern Pennsylvania. Today, only elders speak the language - although some young Lenape youth and adults learn the ancient language. The German and English-speaking Moravian missionary John Heckewelder wrote: "The Monsey tong [sic] is quite different even though [it and Lenape] came out of one parent language." William Penn, who first met the Lenape in 1682, stated that the Unami used the following words: "mother" was anna, "brother" was "isseemus," "friend" was netap. Penn instructed his fellow Englishmen: “If one asks them for anything they have not, they will answer, mattá ne hattá,” which to translate is, not I have, instead of I have not." According to the Moravian missionary David Zeisberger, the Unami word for "food" is May-hoe-me-chink; in Munsee it is Wool-as-gat. The Unami word for "hill" is Ah-choo; in Munsee it is Watts Unk. Sometimes the languages shared words, such as "corn," which is Xash-queem, or "wolf," which is too-may. In contemporary Unami orthography, food is michëwakàn; hill is ahchu; corn is xàskwim; and wolf is tëme. Zeisberger and Heckewelder lived among the Unami and Munsee people in Pennsylvania and Ohio during the late-18th and early-19th centuries and interviewed them. David Zeisberger wrote A Lenâpé-English Dictionary: From An Anonymous [Manuscript] In The Archives Of The Moravian Church At Bethlehem, [Pennsylvania], David Zeisberger's History of Northern American Indians, The Diary of David Zeisberger: A Moravian Missionary Among the Ohio Indians, Grammar of the Language of the Lenni Lenape or Delaware Indians, and Zeisberger’s Indian Dictionary: English, German, Iroquois—The Onondaga and Algonquin—The Delaware. The "Delaware" that Zeisberger translated is Munsee, and not Unami. John Heckewelder wrote extensively on the Lenape in his History, Manners, and Customs of the Indian Nations Who Once Inhabited Pennsylvania and Neighboring States, as well as The Names Which the Lenni Lenape or Delaware Indians Gave to Rivers, Streams, and Localities. At the time of first European contact, a Lenape individual would have identified primarily with his or her immediate family and clan, friends, and/or village unit; then with surrounding and familiar village units; next with more distant neighbors who spoke the same dialect; and ultimately, with all those in the surrounding area who spoke mutually comprehensible languages, including the Nanticoke people, who lived to their south and west in present western Delaware and eastern Maryland, and the Munsee, who lived to their north. Among many Algonquian peoples along the East Coast, the Lenape were considered the "grandfathers" from whom other Algonquian-speaking peoples originated. Consequently, in inter-tribal councils, the Lenape were given respect as one would to elders. Lenape has three phratries, which in turn had twelve clans. These are: By 1682, when William Penn arrived to his American commonwealth, the Lenape had been so reduced by disease, famine, and war that the sub-clan mothers had reluctantly resolved to consolidate their families into the main clan family. This is why William Penn and all those after him believed that the Lenape clans had always only had three divisions ('Turtle, Turkey, and Wolf) when, in fact, they had over thirty on the eve of European contact. For example, some time between 1650 and 1680, the Bear, Deer, etc. families, with few members left, absorbed into the leading Wolf Family. Members of each clan were found throughout Lenape territory and clan lineage was traced through the mother. While clan mothers controlled the land, the houses, and the families, the clan fathers provided the meat, cleared the fields, built the houses, and protected the clan. Upon reaching adulthood, a Lenape male would marry outside of his clan, a practice known by ethnographers as, "exogamy". The practice effectively prevented inbreeding, even among individuals whose kinship was obscure or unknown. This means that a male from the Turkey Clan was expected to marry a female from either the Turtle or Wolf clans. His children, however, would not belong to the Turkey Clan, but to the mother's clan. As such, a person's mother's brothers (the person's matrilineal uncles) played a large role in his or her life as they shared the same clan lineage. To add clarity to the clan system, all males, as a part of their passage rites into adulthood, were tattooed with their clan symbol on their chests. This is why many English, Dutch, and Swedish traders believed that the Lenape had three or more tribes, when in fact, they were one nation of kindred people. Those of a different language stock, such as the Iroquois (or, in the Unami language, the Maax-waas Len [Bear People] or Minquas), were regarded as foreign. As in the case of the Iroquois, the animosity of difference and competition spanned many generations, and different language tribes became traditional enemies. Ethnicity seems to have mattered little to the Lenape and many other "tribes". Archaeological excavations have found Lenape burials that included identifiably ethnic Iroquois remains interred along with those of Lenape. The two groups were bitter enemies since before recorded history, but intermarriage occurred. In addition, both tribes practiced adopting young captives from warfare into their tribes and assimilating them as full tribal members. Early Europeans who first wrote about Indians found matrilineal social organization to be unfamiliar and perplexing. Because of this, Europeans often tried to interpret Lenape society through more familiar European arrangements. As a result, the early records are full of clues about early Lenape society, but were usually written by observers who did not fully understand what they were seeing. For example, a man's maternal uncle (his mother's brother), and not his father, was usually considered to be his closest male ancestor, since his uncle belonged to his mother's clan and his father belonged to a different one. The maternal uncle played a more prominent role in the lives of his sister's children than did the father. Early European chroniclers did not understand this concept. The band assigned land of their common territory to a particular clan for hunting, fishing, and cultivation. Individual private ownership of land was unknown, as the land belonged to the clan collectively while they inhabited it, but women often had rights to traditional areas for cultivation. Clans lived in fixed settlements, using the surrounding areas for communal hunting and planting until the land was exhausted. In a common practice known as "agricultural shifting", the group then moved to found a new settlement within their territory. The Lenape practiced large-scale agriculture to augment a mobile hunter-gatherer society in the regions around the Delaware River. The Lenape were largely a sedentary people who occupied campsites seasonally, which gave them relatively easy access to the small game that inhabited the region: fish, birds, shellfish and deer. They developed sophisticated techniques of hunting and managing their resources. By the time of the arrival of Europeans, the Lenape were cultivating fields of vegetation through the slash and burn technique. This extended the productive life of planted fields. They also harvested vast quantities of fish and shellfish from the bays of the area, and, in southern New Jersey, harvested clams year-round. The success of these methods allowed the tribe to maintain a larger population than nomadic hunter-gatherers could support. Scholars have estimated that at the time of European settlement, there may have been about 15,000 Lenape total in approximately 80 settlement sites around much of the New York City area, alone. In 1524 Lenape in canoes met Giovanni da Verrazzano, the first European explorer to enter New York Harbor. At the time of European contact, the Lenape practiced agriculture, mostly companion planting. The women cultivated many varieties of the "Three Sisters:" corn, beans, and squash. The men also practiced hunting and the harvesting of seafood. The people were primarily sedentary rather than nomadic; they moved to seasonal campsites for particular purposes such as fishing and hunting. European settlers and traders from the seventeenth-century colonies of New Netherland and New Sweden traded with the Lenape for agricultural products, mainly maize, in exchange for iron tools. The Lenape also arranged contacts between the Minquas or Susquehannocks and the Dutch and Swedish West India companies to promote the fur trade. The Lenape were major producers of wampum or shell beads, which they traditionally used for ritual purposes and as ornaments. After the Dutch arrival, they began to exchange wampum for beaver furs provided by Iroquoian-speaking Susquehannock and other Minquas. They exchanged these furs for Dutch and, from the late 1630s, also Swedish imports. Relations between some Lenape and Minqua polities briefly turned sore in the late 1620s and early 1630s, but were relatively peaceful most of the time. The early European settlers, especially the Dutch and Swedes, were surprised at the Lenape's skill in fashioning clothing from natural materials. In hot weather both men and women wore only loin cloth and skirt respectively, while they used beaver pelts or bear skins to serve as winter mantles. Additionally, both sexes might wear buckskin leggings and moccasins in cold weather. Deer hair, dyed a deep scarlet, was a favorite component of headdresses and breast ornaments for males. The Lenape also adorned themselves with various ornaments made of stone, shell, animal teeth, and claws. The women often wore headbands of dyed deer hair or wampum. They painted their skin skirts or decorated them with porcupine quills. These skirts were so elaborately appointed that, when seen from a distance, they reminded Dutch settlers of fine European lace. The winter cloaks of the women were striking, fashioned entirely from the iridescent body feathers of wild turkeys. The first recorded contact with Europeans and people presumed to have been the Lenape was in 1524. The explorer Giovanni da Verrazzano was greeted by local Lenape who came by canoe, after his ship entered what is now called Lower New York Bay. The early interaction between the Lenape and Dutch traders in the 17th century was primarily through the fur trade; specifically, the Lenape trapped and traded beaver pelts for European-made goods. According to Dutch settler Isaac de Rasieres, who observed the Lenape in 1628, the Lenape's primary crop was maize, which they planted in March. They quickly adopted European metal tools for this task. In May, the Lenape planted kidney beans near the maize plants; the latter served as props for the climbing bean vines. They also planted squash, whose broad leaves cut down on weeds and conserved moisture in the soil. The women devoted their summers to field work and harvested the crops in August. Women cultivated varieties of maize, squash and beans, and did most of the fieldwork, processing and cooking of food. The men limited their agricultural labor to clearing the field and breaking the soil. They primarily hunted and fished during the rest of the year. Dutch settler David de Vries, who stayed in the area from 1634 to 1644, described a Lenape hunt in the valley of the Achinigeu-hach (or "Ackingsah-sack," the Hackensack River), in which one hundred or more men stood in a line many paces from each other, beating thigh bones on their palms to drive animals to the river, where they could be killed easily. Other methods of hunting included lassoing and drowning deer, as well as forming a circle around prey and setting the brush on fire. At the time of sustained European contact in the 16th centuries and 17th centuries, the Lenape were a powerful Native American nation who inhabited a region on the mid-Atlantic coast spanning the latitudes of southern Massachusetts to the southern extent of Delaware in what anthropologists call the Northeastern Woodlands. Although never politically unified, the confederation of the Delaware roughly encompassed the area around and between the Delaware and lower Hudson rivers, and included the western part of Long Island in present-day New York. Some of their place names, such as Manhattan, Raritan, and Tappan were adopted by Dutch and English colonists to identify the Lenape people that lived there. Based on the historical record of the mid-seventeenth century, it has been estimated that most Lenape polities consisted of several hundred people but it is conceivable that some had been considerably larger prior to close contact, given the wars between the Susquehannocks and the Iroquois, both of whom were armed by the Dutch fur traders, while the Lenape were at odds with the Dutch and so lost that particular arms race. Smallpox devastated native communities even located far from European settlements by the 1640s. The Lenape and Susquehannocks fought a war in the middle of the 17th century that left the Delaware a tributary state even as the Susquehannocks had defeated the Province of Maryland between 1642-50s, New Amsterdam was founded in 1624 by the Dutch in what would later become New York City. Dutch settlers also founded a colony at present-day Lewes, Delaware on June 3, 1631 and named it Zwaanendael (Swan Valley). The colony had a short life, as in 1632 a local band of Lenape killed the 32 Dutch settlers after a misunderstanding escalated over Lenape defacement of the insignia of the Dutch West India Company. In 1634, the Iroquoian-speaking Susquehannock went to war with the Lenape over access to trade with the Dutch at New Amsterdam. They defeated the Lenape, and some scholars believe that the Lenape may have become tributaries to the Susquehannock. After the warfare, the Lenape referred to the Susquehannock as "uncles." The Iroquois added the Lenape to the Covenant Chain in 1676; the Lenape were tributary to the Five Nations (later Six) until 1753, shortly before the outbreak of the French and Indian War (a part of the Seven Years' War in Europe). The Lenape's quick adoption of trade goods, and their need to trap furs to meet high European demand, resulted in their disastrous over-harvesting of the beaver population in the lower Hudson Valley. With the fur sources exhausted, the Dutch shifted their operations to present-day upstate New York. The Lenape who produced wampum in the vicinity of Manhattan Island temporarily forestalled the negative effects of the decline in trade. Lenape population fell sharply during this period, due to high fatalities from epidemics of infectious diseases carried by Europeans, such as measles and smallpox, to which they had no natural immunity, as the diseases had arisen on the Asian continent and moved west into Europe, where they had become endemic in the cities. The Lenape had a culture in which the clan and family controlled property. Europeans often tried to contract for land with the tribal chiefs, confusing their culture with that of neighboring tribes such as the Iroquois. The Lenape would petition for grievances on the basis that not all their families had been recognized in the transaction (not that they wanted to "share" the land). After the Dutch arrival in the 1620s, the Lenape were successful in restricting Dutch settlement until the 1660s to Pavonia in present-day Jersey City along the Hudson. The Dutch finally established a garrison at Bergen, which allowed settlement west of the Hudson within the province of New Netherland. This land was purchased from the Lenape after the fact. In 1682, William Penn and Quaker colonists created the English colony of Pennsylvania beginning at the lower Delaware River. A peace treaty was negotiated between the newly arriving English and Lenape at what is now known as Penn Treaty Park. In the decades immediately following, some 20,000 new colonists arrived in the region, putting pressure on Lenape settlements and hunting grounds. Although Penn endeavored to live peaceably with the Lenape and to create a colony that would do the same, he also expected his authority and that of the colonial government to take precedence. His new colony effectively displaced many Lenape and forced others to adapt to new cultural demands. Penn gained a reputation for benevolence and tolerance, but his efforts resulted in more effective colonization of the ancestral Lenape homeland than previous ones. William Penn died in 1718. His heirs, John and Thomas Penn, and their agents were running the colony, and had abandoned many of the elder Penn's practices. Trying to raise money, they contemplated ways to sell Lenape land to colonial settlers. The resulting scheme culminated in the so-called Walking Purchase. In the mid-1730s, colonial administrators produced a draft of a land deed dating to the 1680s. William Penn had approached several leaders of Lenape polities in the lower Delaware to discuss land sales further north. Since the land in question did not belong to their polities, the talks came to nothing. But colonial administrators had prepared the draft that resurfaced in the 1730s. The Penns and their supporters tried to present this draft as a legitimate deed. Lenape leaders in the lower Delaware refused to accept it. According to historian Steven Harper, what followed was a "convoluted sequence of deception, fraud, and extortion orchestrated by the Pennsylvania government that is commonly known as the Walking Purchase." In the end, all Lenape who still lived on the Delaware were driven off the remnants of their homeland under threats of violence. Some Lenape polities eventually retaliated by attacking Pennsylvania settlements. When they fought British colonial expansion to a standstill at the height of the Seven Years' War, the British government investigated the causes of Lenape resentment. The British asked William Johnson, Superintendent of Indian Affairs, to lead the investigation. Johnson had become wealthy as a trader and acquired thousands of acres of land in the Mohawk River Valley from the Iroquois Mohawk of New York. Beginning in the 18th century, the Moravian Church established missions among the Lenape. The Moravians required the Christian converts to share their pacifism, as well as to live in a structured and European-style mission village. Moravian pacifism and unwillingness to take loyalty oaths caused conflicts with British authorities, who were seeking aid against the French and their Native American allies during the French and Indian War (Seven Years' War). The Moravians' insistence on Christian Lenapes' abandoning traditional warfare practices alienated mission populations from other Lenape and Native American groups, who revered warriors. The Moravians accompanied Lenape relocations to Ohio and Canada, continuing their missionary work. The Moravian Lenape who settled permanently in Ontario after the American Revolutionary War were sometimes referred to as "Christian Munsee", as they mostly spoke the Munsee branch of the Delaware language. During the French and Indian War, the Lenape initially sided with the French, as they hoped to prevent further British colonial encroachment in their territory. But, such leaders as Teedyuscung in the east and Tamaqua in the vicinity of modern Pittsburgh shifted to building alliances with the English. After the end of the war, however, Anglo-American settlers continued to kill Lenape, often to such an extent that the historian Amy Schutt writes the dead since the wars outnumbered those killed during the war. The Treaty of Easton, signed in 1758 between the Lenape and the Anglo-American colonists, required the Lenape to move westward, out of present-day New York and New Jersey and into Pennsylvania, then Ohio and beyond. Sporadically they continued to raid European-American settlers from far outside the area. In 1763 Bill Hickman, Lenape, warned English colonists in the Juniata River region of an impending attack. Many Lenape joined in Pontiac's War, and were numerous among those Native Americans who besieged Pittsburgh. In April 1763 Teedyuscung was killed when his home was burned. His son Captain Bull responded by attacking settlers from New England who had migrated to the Wyoming Valley of Pennsylvania. The settlers had been sponsored by the Susquehanna Company. The Lenape were the first Indian tribe to enter into a treaty with the new United States government, with the Treaty of Fort Pitt signed in 1778 during the American Revolutionary War. By then living mostly in the Ohio Country, the Lenape supplied the Continental Army with warriors and scouts in exchange for food supplies and security. During the American Revolution, the Munsee-speaking Lenape (then called Delaware) bands of the Ohio Country were deeply divided over which side, if any, to take in the conflict. Their bands lived in numerous villages around their main village of Coshocton. At the time of the Revolutionary War, the Lenape villages lay between the western frontier strongholds of the British and the Patriots: the American colonists had Fort Pitt (present-day Pittsburgh) and the British with Indian allies controlled the area of Fort Detroit (in present-day Michigan). Some Lenape decided to take up arms against the American colonials and moved to the west, closer to Detroit, where they settled on the Scioto and Sandusky rivers. Those Lenape sympathetic to the United States remained at Coshocton, and leaders signed the Treaty of Fort Pitt (1778) with the Americans. Through this, the Lenape hoped to establish the Ohio Country as a state inhabited exclusively by Native Americans, as part of the new United States. A third group of Lenape, many of them converted Christian Munsees, lived in several mission villages run by Moravians. (They spoke the Munsee branch of Delaware, an Algonquian language.) White Eyes, the Lenape chief who had negotiated the treaty, died in 1778. Many Lenape at Coshocton eventually joined the war against the Americans. In response, Colonel Daniel Brodhead led an expedition out of Fort Pitt and on 19 April 1781 destroyed Coshocton. Surviving residents fled to the north. Colonel Brodhead convinced the militia to leave the Lenape at the Moravian mission villages unmolested, since they were unarmed non-combatants. Brodhead's having to restrain the militia from attacking the Moravian villages was a reflection of the brutal nature of frontier warfare. Violence had escalated on both sides. Relations between regular Continental Army officers from the East (such as Brodhead) and western militia were frequently strained. The tensions were worsened by the American government's policy of recruiting some Indian tribes as allies in the war. Western militiamen, many of whom had lost friends and family in Indian raids against settlers' encroachment, blamed all Indians for the acts of some. During the early 1770s, missionaries, including David Zeisberger and John Heckewelder, arrived in the Ohio Country near the Delaware villages. The Moravian Church sent these men to convert the natives to Christianity. The missionaries established several missions, including Gnadenhutten, Lichtenau, and Schoenbrunn. The missionaries asked that the natives forsake all of their traditional customs and ways of life. Many Delaware did adopt Christianity, but others refused to do so. The Delaware became a divided people during the 1770s, including in Killbuck's family. Killbuck resented his grandfather for allowing the Moravians to remain in the Ohio Country. The Moravians believed in pacifism, and Killbuck believed that every convert to the Moravians deprived the Delaware of a warrior to stop further white settlement of their land. During the French and Indian War, Killbuck assisted the English against their French enemy. In 1761, Killbuck led an English supply train from Fort Pitt to Fort Sandusky. The British paid him one dollar per day. Later Killbuck became a leader in a very dangerous time for the Delaware. The American Revolution had just begun, and Killbuck found his people caught between the English in the West and the Americans in the East. At the war's beginning, Killbuck and many Delaware claimed to be neutral. In 1778, Killbuck permitted American soldiers to traverse Delaware territory so that the soldiers could attack Fort Detroit. In return, Killbuck requested that the Americans build a fort near the natives' major village of Coshocton to provide the Delaware with protection from English attacks. The Americans agreed and built Fort Laurens, which they garrisoned. Other Indian groups, especially the Wyandot, the Mingo, the Munsee, the Shawnee, and the Wolf Clan of the Delaware, favored the British. They believed that by their proclamation of 1763, restricting Anglo-American settlement to east of the Appalachian Mountains, that the British would help them preserve a Native American territory. The British planned to attack Fort Laurens in early 1779 and demanded that the neutral Delawares formally side with the British. Killbuck warned the Americans of the planned attack. His actions helped save the fort, but the Americans abandoned it in August 1779. The Delaware had lost their protectors and, in theory, faced attacks from the British, their native allies, and the American settlers who flooded into the area in the late 1770s and early 1780s after the war. Most Delaware formally joined the British after the American withdrawal from Fort Laurens. Facing pressure from the British, the Americans, and even his fellow natives, Killbuck hoped a policy of neutrality would save his people from destruction. It did not. The amateur anthropologist Silas Wood published a book claiming that there were several American Indian tribes that were distinct to Long Island, New York. He collectively called them the Metoac. Modern scientific scholarship has shown that two linguistic groups representing two Algonquian cultural identities lived on the island, not "13 individual tribes" as asserted by Wood. The bands to the west were Lenape. Those to the east were more related culturally to the Algonquian tribes of New England across Long Island Sound, such as the Pequot. Wood (and earlier settlers) often misinterpreted the Indian use of place names for identity as indicating their names for "tribes." Over a period of 176 years, European settlers progressively crowded the Lenape out of the East Coast and Ohio, and pressed them to move further west. Most members of the Munsee-language branch of the Lenape left the United States after the British were defeated in the American Revolutionary War. Their descendants live on three Indian reserves in Western Ontario, Canada. They are descendants of those Lenape of Ohio Country who sided with the British during the Revolutionary War. The largest reserve is at Moraviantown, Ontario, where the Turtle Phratry settled in 1792 following the war. Two groups migrated to Oneida County, New York by 1802, the Brotherton Indians of New Jersey and the Stockbridge-Munsee. After 1819, they removed to Wisconsin, under pressure from state and local governments. By the Treaty of St. Mary’s, signed October 3, 1818 in St. Mary's, Ohio, the Delaware ceded their lands in Indiana for lands west of the Mississippi and an annuity of $4,000. Over the next few years, the Delaware settled on the James River in Missouri near its confluence with Wilsons Creek, occupying eventually about 40000acres of the approximately 2000000acres allotted to them. Anderson, Indiana is named after Chief William Anderson, whose father was Swedish. The Delaware Village in Indiana was called Anderson's Town, while the Delaware Village in Missouri on the James River was often called Anderson’s Village. The tribes' cabins and cornfields were spread out along the James River and Wilsons Creek. Many Delaware participated in exploration of the western United States, working as trappers with the mountain men, and as guides and hunters for wagon trains. They served as army guides and scouts in events such as the Second Seminole War, Frémont's expeditions, and the conquest of California during the Mexican-American War. Occasionally, they played surprising roles as Indian allies. Sagundai accompanied one of Frémont's expeditions as one of his Delaware guides. From California, Fremont needed to communicate with Senator Benton. Sagundai volunteered to carry the message through some 2,200 kilometres of hostile territory. He took many scalps in this adventure, including that of a Comanche with a particularly fine horse, who had outrun both Sagundai and the other Comanche. Sagundai was thrown when his horse stepped into a prairie-dog hole, but avoided the Comanche's lance, shot the warrior dead, and caught his horse and escaped the other Comanche. When Sagundai returned to his own people in present-day Kansas, they celebrated his exploits with the last war and scalp dances of their history. These were held at Edwardsville, Kansas. By the terms of the "Treaty of the James Fork" made September 24, 1829 and ratified by the US Senate in 1830, the Delaware were forced to move further west. They were granted lands in Indian Territory in exchange for lands on the James Fork of the White River in Missouri. These lands, in what is now Kansas, were west of the Missouri and north of the Kansas River. The main reserve consisted of about 1000000acres with an additional "outlet" strip 10miles wide extending to the west. In 1854 Congress passed the Kansas–Nebraska Act, which created the Territory of Kansas and opened the area for white settlement. It also authorized negotiation with Indian tribes regarding removal. The Delaware were reluctant to negotiate for yet another relocation, but they feared serious trouble with white settlers, and conflict developed. As the Delaware were not considered United States citizens, they had no access to the courts, and no way to enforce their property rights. The United States Army was to enforce their rights to reservation land after the Indian Agent had both posted a public notice warning trespassers and served written notice on them, a process generally considered onerous. Major B.F. Robinson, the Indian Agent appointed in 1855, did his best, but could not control the hundreds of white trespassers who stole stock, cut timber, and built houses and squatted on Delaware lands. By 1860 the Delaware had reached consensus to leave Kansas, which was in accord with the government's Indian removal policy. The main body of Lenape arrived in Indian Territory in the 1860s. As a result of the multiple removals, each leaving some Lenape who chose to stay in place, Lenape people and descendents are located today in New Jersey, Wisconsin and southwest Oklahoma. The two largest groups are the Delaware Nation (Anadarko, Oklahoma) and the Delaware Tribe of Indians (Bartlesville, Oklahoma). The Delaware Tribe of Indians were required to purchase land from the reservation of the Cherokee Nation; they made two payments totaling $438,000. A court dispute followed over whether the sale included rights for the Delaware as citizens within the Cherokee Nation. While the dispute was unsettled, the Curtis Act of 1898 dissolved tribal governments and ordered the allotment of communal tribal lands to individual households of members of tribes. After the lands were allotted in 160-acre (650,000 m²) lots to tribal members in 1907, the government sold "surplus" land to non-Indians. The Delaware migrated into Texas in the late 18th and early 19th centuries. Elements of the Delaware migrated from Missouri into Texas around 1820, settling around the Red River and Sabine River. The Delaware were peaceful and shared their territory in Spanish Texas with the Caddo and other immigrating bands, as well as with the Spanish and ever-increasing American population. This peaceful trend continued after Mexico won their independence from Spain in 1821. In 1828, Mexican General Manuel de Mier y Terán made an inspection of eastern Mexican Texas and estimated that the region housed between 150 to 200 Delaware families. The Delaware requested Mier y Terán to issue them land grants and send teachers, so they might learn to read and write the Spanish language. The General, impressed with how well they had adapted to the Mexican culture, sent their request to Mexico City, but the authorities never granted the Delaware any legal titles. The situation changed when the Texas Revolution began in 1835. Texas officials were eager to gain the support of the Texas tribes to their side and offered to recognize their land claims by sending three commissioners to negotiate a treaty. A treaty was agreed upon in February 1836 which mapped the boundaries of Indian lands; but, this agreement was never officially ratified by the Texas government. The Delaware remained friendly after Texas won its independence. Republic of Texas President, Sam Houston favored a policy of peaceful relations with all tribes. He sought the services of the friendly Delaware and in 1837 enlisted several Delaware to protect the frontier from hostile western tribes. Delaware scouts joined with Texas Rangers as they patrolled the western frontier. Houston also tried to get the Delaware land claims recognized but his efforts were only met by opposition. The next Texan President, Mirabeau B. Lamar, completely opposed all Indians. He considered them as illegal intruders who threatened the settlers safety and lands and issued an order for their removal from Texas. The Delaware were sent north of the Red River into Indian Territory, however, a few scattered Delawares remained in Texas. In 1841, Houston was reelected to a second term as president and his peaceful Indian policy was then reinstated. A treaty with the remaining Delaware and a few other tribes was negotiated in 1843 at Fort Bird and the Delaware were enlisted to help him make peace with the Comanche. Delaware scouts and their families were allowed to settle along the Brazos and Bosque rivers in order to influence the Comanche to come to the Texas government for a peace conference. The plan was successful and the Delaware helped bring the Comanches to a treaty council in 1844. In 1845, the Republic of Texas agreed to annexation by the US to become an American State. The Delaware continued their peaceful policy with the Americans and served as interpreters, scouts and diplomats for the US Army and the Indian Bureau. In 1847, John Meusebach was assisted by Jim Shaw (Delaware), in settling the German communities in the Texas Hill Country. For the remainder of his life, Shaw worked as a military scout in West Texas. In 1848, John Conner (Delaware) guided the Chihuahua-El Paso Expedition and was granted a league of land by a special act of the Texas legislature in 1853. The expeditions of the map maker Randolph B. Marcy through West Texas in 1849, 1852, and 1854 were guided by Black Beaver (Delaware). In 1854, despite the history of peaceful relations, the last of the Texas Delaware were moved by the American government to the Brazos Indian Reservation near Graham, Texas. In 1859 the US forced the remaining Delaware to remove from Texas to a location on the Washita River in the vicinity of present Anadarko, Oklahoma. In 1979, the United States Bureau of Indian Affairs revoked the tribal status of the Delaware living among Cherokee in Oklahoma. They began to count the Delaware as Cherokee. The Delaware had this decision overturned in 1996, when they were recognized by the federal government as a separate tribal nation. The Cherokee Nation filed suit to overturn the independent federal recognition of the Delaware. The tribe lost federal recognition in a 2004 court ruling in favor of the Cherokee Nation, but regained it on 28 July 2009. After recognition, the tribe reorganized under the Oklahoma Indian Welfare Act. Members approved a constitution and by laws in a May 26, 2009 vote. Jerry Douglas was elected as tribal chief. In 2004, the Delaware Nation filed suit against Pennsylvania in the United States District Court for the Eastern District of Pennsylvania, seeking to reclaim 315acres included in the 1737 Walking Purchase to build a casino. In the suit titled "The Delaware Nation v. Commonwealth of Pennsylvania" the plaintiffs acting as the successor in interest and political continuation of the Lenni Lenape and of Lenape Chief “Moses” Tundy Tatamy, claimed aboriginal and fee title to the 315 acres of land located in Forks Township in Northampton County, near the town of Tatamy, Pennsylvania. After the Walking Purchase, Chief Tatamy was granted legal permission for him and his family to remain on this parcel of land, known as “Tatamy's Place". In addition to suing the state, the tribe also sued the township, the county and elected officials, including Gov. Ed Rendell. The court held that the justness of the extinguishment of aboriginal title is nonjusticiable, including in the case of fraud. Because the extinguishment occurred prior to the passage of the first Indian Nonintercourse Act in 1790, that Act did not avail the Delaware. As a result the court granted the Commonwealth's motion to dismiss. In its conclusion the court stated: ... we find that the Delaware Nation's aboriginal rights to Tatamy's Place were extinguished in 1737 and that, later, fee title to the land was granted to Chief Tatamy-not to the tribe as a collectivity. Three Lenape tribes are federally recognized in the United States. They are as follows: The Canadian Lenape left the United States in the late 1700s following the American Revolutionary War and settled in what is now Ontario. Consequently, Canada recognizes three Lenape First Nations (with four Indian reserves); they are located in Southwestern Ontario: New Jersey has two state recognized tribes, who are in part Lenape: the Nanticoke Lenni-Lenape Indians of New Jersey and Ramapough Lenape Nation. In Delaware, the Lenape are organized and state-recognized as the Lenape Indian Tribe of Delaware. Some Lenape or Delaware live in communities known as Urban Indians in their historic homeland in a number of states such as Delaware, Maryland, New Jersey, and Virginia. New York City and Philadelphia are known to have some Lenape residents. Some Lenape live within the city of Tulsa, Oklahoma, and others live in diaspora across the country. Large communities of Lenape people live in the vicinities of Bartlesville, Oklahoma and Anadarko, Oklahoma. Additionally, over a dozen unrecognized tribes claim Lenape descent. Unrecognized Lenape organizations in Colorado, Idaho, and Kansas have petitioned the United States federal government for recognition. The Walam Olum, which purported to be an account of the Delaware's migration to the lands around the Delaware River, emerged through the works of Constantine Samuel Rafinesque in the 19th century. For many decades, scholars believed it was genuine. In the 1980s and 1990s, newer textual analysis suggested it was a hoax. In Cormac McCarthy's Blood Meridian, the group of American scalphunters are aided by an unspecified number of Delaware, who serve as scouts and guides through the western deserts. In The Light in the Forest, True Son is adopted by a band of Lenape. In Mark Raymond Harrington's 1938 novel, The Indians of New Jersey: Dickon among the Lenapes, a group of Lenape find a shipwrecked English boy. His gradual integration into the tribe provides a study of Lenape life, society, weaponry, and beliefs. The book includes a glossary for Lenape terms. Trouble's Daughter: The Story of Susanna Hutchinson, Indian Captive is a young adult novel of a fictional kidnapping by the Lenape Turtle Clan of a daughter of Anne Hutchinson, the religious reformer and founder of the Rhode Island colony. Moon of Two Dark Horses is a novel of the friendship between a white settler and a Lenape boy at the time of the Revolutionary War. Standing in the Light, The Captive Diary of Catherine Carey Logan, part of the Dear America series of fictional diaries, is a novel by Mary Pope Osborne. It tells the story of the capture of a teenage girl and her brother by a band of Lenape, and the youths' assimilation into Lenape culture. Peter (Per) Lindeström's Geographia Americae with an Account of the Delaware Indians is one of the few sympathetic contemporary accounts of Lenape life in the lower Delaware River valley during the 17th century. Moravian missionary John Heckewelder published a sympathetic account of the Lenape in exile in the Ohio Valley. His account, published in 1818, provides some alternate Lenape tribal history disputing the tributary relationship with the Susquehannock. "Scouts of '76: a tale of the revolutionary war", a 1924 book by Charles E. Willis, contains an account of the contributions of the Lenni Lenape to the American Revolution when they lived in the area of Lake Wawayanda. The Ramapough Lenape Nation is central to the Sundance Channel series The Red Road, in a newly (Federally) recognized reservation straddling the border between New York and New Jersey.

http://everything.explained.today/Lenape/

Current Access Through Your IP Address Current Access Through Your Registered Email Address A subscription to JoVE is required to view this article. You will only be able to see the first 20 seconds. Ligation can be defined as the act of joining, and in biology the term refers to an enzymatic reaction that joins two biomolecules with a covalent bond. This video describes the application of DNA ligation in molecular biology research. In the cell, DNA ligases are enzymes that identify and seal breaks in DNA by catalyzing the formation of phosphodiester bonds between the 3’-hydroxyl and 5’-phosphate groups of the DNA backbone. Ligation occurs as part of normal cellular processes, such as DNA replication, to repair single and double strand DNA breaks. In the laboratory, DNA ligases is routinely used in molecular cloning - a process that joins endonuclease-digested DNA fragments, or inserts, with an endonuclease-digested vector, such as a plasmid, so that the fragment can be introduced into host cells and then replicated. Endonuclease digestions involve the use of restriction endonucleases, or restriction enzymes, which create nicks at specific stretches of DNA. These nicks can resemble single strand breaks producing 3’ and 5’ overhangs, called sticky ends or double strand breaks with no overhangs, called blunt-ends. Ligating sticky ends is advantageous, because the complimentary overhanging base pairs stabilize the reaction. Because blunt end ligations don’t have any complimentary base pairing, the ligation is less efficient and more difficult for the enzyme to join the ends. Sticky and blunt ends cannot, under normal circumstances, be ligated together. However, the Klenow fragment, the product of DNA polymerase 1, digested with subtilisin can convert sticky ends to blunt ends. Klenow possesses 3’ to 5’ exonuclease activity that chews up 3’ overhangs and polymerase activity that blunts 5’overhangs by extending the 3’ end of the complementary strand. When the goal is to insert a gene into a plasmid, resealing of vector DNA, called self-ligation, is a common undesirable outcome for a ligation reaction. Alkaline phosphatase treatment of vector DNA post-digestion removes 5’phosphates on both ends and prevents this undesirable outcome. As we mentioned previously, vector and insert DNAs are digested with endonucleases prior to beginning a ligation. Following gel-purification of digested vector and insert, DNA concentrations are measured a spectrophotometer to determine the concentration of the purified vector and insert. From this concentration, the number of molecules of insert or vector in 1 µl can be determined based on the average molecular weight for a DNA base pair and the number of base pairs in each fragment. Based on the calculated molecular concentration of vector and insert, a 3 to 1 ratio of insert to vector is calculated, to determine the volume of vector and insert used in the reaction. This 3 to 1 ratio of DNA insert to vector is desirable, because it ups the probability of the insert being ligated into vector versus vector ligating itself. Now that we have determined the amount of vector and insert DNA to use in the reaction, we proceed to set up the ligation reaction on ice. The order of adding in which reaction components should be added to your microfuge tube is as follows: sterile water enough to a make a 10 µl final volume, in our case we’ll use 4 µl, 1 µl 10X of ligation buffer, 1 µl 10mM of ATP, 1 µl of vector and 3 µl insert DNA, as calculated, and finally 1 µl DNA Ligase. The reaction is mixed thoroughly, centrifuged and incubated at the appropriate temperature. Whether you are doing a sticky or blunt end ligation impacts the temperature and duration of the ligation reaction. For example, a sticky end ligation with a six base pair overhang can be carried out near room temperature for about 1 hr, because the complementary ends stabilize the joining of fragments. Short overhangs or blunt end ligations should be carried out between 14-20˚C overnight. Now that we learned how to set up a ligation reaction, let’s have a look at some of the applications of this procedure. Ligations can be used to directly insert PCR-amplified fragments into linearized plasmids. Here you see a researcher taking a sample of frozen mouse brain, isolating genomic DNA from it, and then subjecting it to bisulfite PCR, which is a PCR-based method to detect methylated DNA. PCR products are then directly ligated into the plasmid to create a library of genes that are methylated in that particular brain region. Ligations can be used to attach oligonucleotide linkers, which contain binding sites for PCR primers, to purify DNA fragments. When working with tumor samples, scientists can use this approach to sequence tumor genomic DNA, with the hope of identifying tumor-causing mutations. In this video, ligation is performed on DNA isolated from formaldehyde fixed cells and subsequently treated with a restriction enzyme and klenow in presence of biotin, which is then used to pull down ligated DNA. This DNA is then amplified using PCR and the products sequenced to identify chromatin interactions at various scales as shown. You have now learned about DNA ligase, various principles involved in setting up ligation in the laboratory, potential problems and fixes and various applications of ligation in molecular biology research. Thanks for watching. In molecular biology, ligation refers to the joining of two DNA fragments through the formation of a phosphodiester bond. An enzyme known as a ligase catalyzes the ligation reaction. In the cell, ligases repair single and double strand breaks that occur during DNA replication. In the laboratory, DNA ligase is used during molecular cloning to join DNA fragments of inserts with vectors – carrier DNA molecules that will replicate target fragments in host organisms. This video provides an introduction to DNA ligation. The basic principle of ligation is described as well as a step-by-step procedure for setting up a generalized ligation reaction. Critical aspects of ligation reactions are discussed, such as how the length of a sticky end overhang affects the reaction temperature and how the ratio of DNA insert to vector should be tailored to prevent self-ligation. Molecular tools that assist with ligations like the Klenow Fragment and shrimp alkaline phosphatase (SAP) are mentioned, and applications , such as proximity ligations and the addition of linkers to fragments for sequencing are also presented. JoVE Science Education Database. Basic Methods in Cellular and Molecular Biology. DNA Ligation Reactions. JoVE, Cambridge, MA, doi: 10.3791/5069 (2016). In this video, PCR is used to amplify regions in bacterial genomic DNA called clustered regularly interspaced short palindromic repeat (CRISPR) sequences. Ligations are used to introduce endonuclease-digested and gel-purified fragments into a vector, which is then transformed into E. coli. CRISPR sequences are of interest to scientists, because they are important components of bacterial defense against viral infections. The telomere is a repetitive nucleotide sequence at the end of a chromosome, which protects an organism’s genetic material from degradation. In this video, the ligation of adaptors containing PCR primers is used to determine the G-overhang structure of the telomeres found in Trypanosoma brucei - the causative agent of African sleeping sickness. A BioBrick part is a DNA sequence of defined structure and function that has standardized upstream and downstream sequences. Ligation is used in this video to introduce Biobrick parts into a plasmid that enables E. coli to metabolize hydrocarbons. This proof of concept study shows the possibility of a sustainable approach to oil-remediation through synthetic biology. In this article, linearized ?-phage DNA is modified through the annealing and ligation of modified oligonucleotides to form a replication fork. The ligation of biotin and digoxygenin labeled probes on opposite ends of the ?-phage replication fork allows for the real time observation of DNA replication through microscopy. Ligation is used to add biotin and digoxigenin to ? DNA. The conjugated DNA is adhered to a flow cell and magnetic beads were attached for uses with magnetic tweezers. This procedure allows for the measurement of forces exerted by individual proteins.

http://www.jove.com/science-education/5069/dna-ligation-reactions

Grading on a curve |This article needs additional citations for verification. (October 2014)| |Some of this article's listed sources may not be reliable. (October 2014)| ||The neutrality of this article is disputed. (September 2015)| In education, grading on a curve (also referred to as curved grading, bell curving, or using grading curves) is a statistical method of assigning grades designed to yield a pre-determined distribution of grades among the students in a class. The term "curve" refers to the bell curve, the graphical representation of the probability density of the normal distribution (also called the Gaussian distribution), but this method does not necessarily use any specific frequency distribution. One method of applying a curve uses three steps: - Numeric scores (or possibly scores on a sufficiently fine-grained ordinal scale) are assigned to the students. The absolute values are less relevant, provided that the order of the scores corresponds to the relative performance of each student within the course. - These scores are converted to percentiles (or some other system of quantiles). - The percentile values are transformed to grades according to a division of the percentile scale into intervals, where the interval width of each grade indicates the desired relative frequency for that grade. For example, if there are five grades in a particular university course, A, B, C, D, and F, where A is reserved for the top 20% of students, B for the next 30%, C for the next 30%-40%, and D or F for the remaining 10%-20%, then scores in the percentile interval from 0% to 20% will receive a grade of D or F, scores from 21% to 50% will receive a grade of C, scores from 51% to 80% receive a grade of B, and scores from 81% to 100% will achieve a grade of A. Consistent with the example illustrated above, a grading curve allows academic institutions to ensure the distribution of students across certain grade point average (GPA) thresholds. As many professors establish the curve to target a course average of a C, the corresponding grade point average equivalent would be a 2.0 on a standard 4.0 scale employed at most North American universities. Similarly, a grade point average of 3.0 on a 4.0 scale would indicate that the student is within the top 20% of the class. Grading curves serve to attach additional significance to these figures, and the specific distribution employed may vary between academic institutions. The ultimate objective of grading curves is to minimize or eliminate the influence of variation between different instructors of the same course, ensuring that the students in any given class are assessed relative to their peers. This also circumvents problems associated with utilizing multiple versions of a particular examination, a method often employed where test administration dates vary between class sections. Regardless of any difference in the level of difficulty, real or perceived, the grading curve ensures a balanced distribution of academic results.

https://en.wikipedia.org/wiki/Grading_curve

Mohandas Karamchand Gandhi (1869-1948) was the most important Indian political and spiritual leader of the 20th century. Gandhi's influence was so great that his methods were later adopted by many political activists around the world, including American civil rights leaders, such as Martin Luther King Jr. Gandhi was born into a middle-class Hindu family, in the city of Porbander, a small town on the western coast of India. At the age of 13, Gandhi entered into an arranged marriage with a 10-year-old girl named Kasturba. (They were to remain married their entire lives.) In 1888, at the age of 19, Gandhi traveled to England to study law. After three years, he became a lawyer and returned to India, and after a year of practicing law unsuccessfully, he was offered a job by an Indian businessman with interests in South Africa. In 1892, at the age of 23, Gandhi traveled to South Africa, where he was to remain for over 20 years. At the time, the Indians in South Africa, mostly Hindus, had no legal rights. The European colonialists did not consider Hindus to be full human beings and referred to them as "coolies". Gandhi became a leader of the Indian community and, over the years, developed a political movement based on the methods of non-violent civil disobedience, which he called "satyagraha". Around 1905, Gandhi gave up Western ways and, for the rest of his life, followed the traditional Hindu precepts of austerity and self-denial. He dressed simply, in a loin cloth and shawl, and had no other material possessions. In 1915, at the age of 46, Gandhi returned to India, where he spent a year traveling widely and then the next few years, helping to settle many local disputes. His success lead to him being admired throughout the country, so much so that India's most well-known writer, Rabindranath Tagore, gave Gandhi the title Mahatma ("Great Soul"). Gandhi himself, however, repudiated the honor, even though, within the Hindu culture, being called "Mahatma" is a symbol of enormous respect. At the time Gandhi was born, India was a heterogeneous region, a British colony consisting of more than 500 different "native states", that is, kingdoms and principalities. (Gandhi himself was born in the state of Kathiawar.) The native states were allowed a certain degree of local autonomy, but the country as a whole was controlled by strict British authority. Soon after his return to India, Gandhi dedicated himself to the goal of Indian independence. From 1920-1922, he led a "non-cooperation movement", in which he called upon Indians to stop cooperating with the British, to become self-reliant, and to withdraw from British organizations. In 1922, the British authorities imprisoned Gandhi on charges of sedition (that is, inciting rebellion). In 1925, Gandhi was released due to ill health but, over his lifetime, he was to be imprisoned many times. Gandhi became a social reformer, working tirelessly to enhance Hindu-Muslim relations, as he slowly led his country into independence. Over the years, he founded many newspapers, which he used to further his ideals. (A little known fact is that Gandhi is one of the principal figures in the history of Indian journalism.) Gandhi developed satyagraha into a national movement, stressing passive resistance, nonviolent disobedience, boycotts and, on occasion, hunger strikes. He became so well-known and respected, that he gained influence with both the general public and the British rulers. For example, in 1939, by a combination of fasting and satyagraha, Gandhi was able to compel several states, that were ruled by princes, to grant democratic reforms. Not only could he unify the many diverse elements of the Indian National Congress, he was able to force political concessions from the British by threatening to fast until death. After World War II, Gandhi was involved in the deliberations that led to India's independence. The same deliberations, however, also led to partition of India into two countries: modern-day India (primarily for Hindus) and Pakistan (for Muslims). Gandhi strongly opposed this partition, which ultimately resulted in the death of about 1 million people and the dislocation of over 11 million people. Although Gandhi was a man of faith, he did not found a church, nor did he create any specific dogma for his followers. On January 30, 1948, just after India attained its independence from Britain, Gandhi was assassinated at the age of 78. The killer was a Hindu fanatic working as part of a conspiracy that blamed Gandhi for the partition of the country. Although Gandhi was a man of faith, he did not found a church, nor did he create any specific dogma for his followers. Gandhi believed in the unity of all mankind under one god, and preached Hindu, Muslim and Christian ethics. As a youth, he was neither a genius nor a child prodigy. Indeed, he suffered from extreme shyness. However, he approached life as a very long series of small steps towards his goals, which he pursued relentlessly. By the time he died, India had become an independent country, free of British rule, in fact, the largest democracy in the world, mostly Hindu with a sizable Muslim minority. Today, Gandhi is remembered not only as a political leader, but as a moralist who appealed to the universal conscience of mankind. As such, he changed the world. © All contents Copyright 2016, Harley Hahn

http://www.harley.com/people/mohandas-gandhi.html

Skip to main content Get your brand new Wikispaces Classroom now and do "back to school" in style. Pages and Files Feedback on TRAILS Getting Started With TRAILS Helpful Tips When Using TRAILS Resources for Educators TRAILS Scholarly Publication TRAILS- Free Information Literacy Assessments Promote Your Page Too Please add any and all lesson plans here. Lesson Plans from Our Previous Wiki Develop, Use, and Revise Search Strategies Identify Potential Sources Great Lessons on Evaluating Sources from Our Previous Wiki Evaluating Web Sites Assignment via Pensacola Catholic High School Website Evaluation Form The Ten Finger Method to Evaluate Websites MANY activities, lessons, and tips for evaluating internet resources - includes misleading that can be used as a classroom activity. Give some students reliable sources and give others unreliable/false once. Have them do research and present to the class their findings on their source. Do not tell them that some are not reliable. This highlights the further need for the evaluation of internet sources. Lessons and Resources for Evaluating Website Content - includes additional hoax Mrs. V's Lessons Adventures of Cyber Bee Grades 9-12 Evaluation Activities Grades 6-8 Identifying Good Sites Grades 4-5 Rating Websites Grades 2-3 Finding Good Sites Grades K-1 What's a Good Site? Kathy Schrock's Guide to Website Evaluation - Very good including step by step instructions. Evaluating Internet Research Sources - Practical ways to look at evaluating resources. Evaluation Wizard for Websites Using Information Responsibly: Lessons for Teaching Students How to Use Information Responsibly Plagiarism Lessons for Grades 9-11 The Purdue Online Writing Lab Cyber Bee - Citing Electronic Resources Resource Citation Including Email, ListServs, and Instant Messenger MANY Citation Resources Including Elementary, Middle, and High School Students help on how to format text Turn off "Getting Started"

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|This article is part of a series on the politics and government of |Titles and honours| |Precedent and law| The Roman Constitution was an uncodified set of guidelines and principles passed down mainly through precedent. The Roman constitution was not formal or even official, largely unwritten and constantly evolving. Having those characteristics, it was therefore more like the British common law system than a statutory law system like the written United States Constitution, even though the constitution's evolution through the years was often directed by passage of new laws and repeal of older ones.. Concepts that originated in the Roman constitution live on in both forms of government to this day. Examples include checks and balances, the separation of powers, vetoes, filibusters, quorum requirements, term limits, impeachments, the powers of the purse, and regularly scheduled elections. Even some lesser used modern constitutional concepts, such as the bloc voting found in the electoral college of the United States, originate from ideas found in the Roman constitution. Over the years, the Roman constitution continuously evolved. By 573 BC, the Constitution of the Roman Kingdom had given way to the Constitution of the Roman Republic. By 27 BC, the Constitution of the Roman Republic had given way to the Constitution of the Roman Empire. By 300 AD, the Constitution of the Roman Empire had given way to the Constitution of the Late Roman Empire. The actual changes, however, were quite gradual. Together, these four constitutions formed four epochs in the continuous evolution of one master constitution. The Roman senate was the most permanent of all of Rome's political institutions. It was probably founded before the first king of Rome ascended the throne. It survived the fall of the Roman Kingdom in 510 BC, the fall of the Roman Republic in 27 BC, and the fall of the Roman Empire in 476 AD. It was, in contrast to many modern institutions named 'Senate', not a legislative body. The power of the senate waxed and waned throughout its history. During the days of the kingdom, it was little more than an advisory council to the king. The last king of Rome, the tyrant Lucius Tarquinius Superbus, was overthrown following a coup d'état that was planned in the senate. During the early republic, the senate was politically weak. During these early years, the executive magistrates were quite powerful. The transition from monarchy to constitutional rule was probably more gradual than the legends suggest. Thus, it took a prolonged weakening of these executive magistrates before the senate was able to assert its authority over those magistrates. By the middle republic, the senate reached the apex of its republican power. This occurred because of the convergence of two factors. The plebeians had recently achieved full political enfranchisement. Therefore, they were not as aggressive as they had been during the early republic in pushing for radical reforms. In addition, the period was marked by prolonged warfare against foreign enemies. The result was that both the popular assemblies and the executive magistrates deferred to the collective wisdom of the senate. The late republic saw a decline in the senate's power. This decline began following the reforms of the radical tribunes Tiberius and Gaius Gracchus. The declining influence of the senate during this era, in large part, was caused by the class struggles that had dominated the early republic. The end result was the overthrow of the republic, and the creation of the Roman Empire. The senate of the very early Roman Empire was as weak as it had been during the late republic. However, after the transition from republic to empire was complete, the senate arguably held more power than it had held at any previous point. All constitutional powers (legislative, executive and judicial) had been transferred to the senate. However, unlike the senate of the republic, the senate of the empire was dominated by the emperor. It was through the senate that the emperor exercised his autocratic powers. By the late principate, the senate's power had declined into near-irrelevance. It never again regained the power that it had held before that point. Much of the surviving literature from the imperial period was written by senators. To a large degree, this demonstrates the strong cultural influence of the senate, even during the late empire. The institution survived the fall of the Empire in the West, and even enjoyed a modest revival as imperial power was reduced to a government of Italy only. The senatorial class was severely affected by the Gothic wars. The first Roman assembly, the Comitia Curiata, was founded during the early kingdom. Its only political role was to elect new kings. Sometimes, the king would submit his decrees to it for ratification. During the early republic, the Comitia Curiata was the only legislative assembly with any power. Shortly after the founding of the republic, however, the Comitia Centuriata and the Comitia Tributa became the predominant legislative assemblies. Most modern legislative assemblies are bodies consisting of elected representatives. Their members typically propose and debate bills. These modern assemblies use a form of representative democracy. In contrast, the assemblies of the Roman Republic used a form of direct democracy. The Roman assemblies were bodies of ordinary citizens, rather than elected representatives. In this regard, bills voted on (called plebiscites) were similar to modern popular referenda. Unlike many modern assemblies, Roman assemblies were not bicameral. That is to say that bills did not have to pass both major assemblies in order to be enacted into law. In addition, no other branch had to ratify a bill (rogatio) in order for it to become law (lex). Members also had no authority to introduce bills for consideration; only executive magistrates could introduce new bills. This arrangement is also similar to what is found in many modern countries. Usually, ordinary citizens cannot propose new laws for their enactment by a popular election. Unlike many modern assemblies, the Roman assemblies also had judicial functions. After the founding of the empire, the powers of the assemblies were transferred to the senate. When the senate elected magistrates, the results of those elections would be read to the assemblies. Occasionally, the emperor would submit laws to the Comitia Tributa for ratification. The assemblies ratified laws up until the reign of the emperor Domitian. After this point, the assemblies simply served as vehicles through which citizens would organize. During the years of the Roman Kingdom, the king (rex) was the only executive magistrate with any power. He was assisted by two quaestors, whom he appointed. He would often appoint other assistants for other tasks. When he died, an interrex would preside over the senate and assemblies, until a new king was elected. Under the Constitution of the Roman Republic, the "executive branch" was composed of both ordinary as well as extraordinary magistrates. Each ordinary magistrate would be elected by one of the two major Legislative Assemblies of the Roman Republic. The principal extraordinary magistrate, the dictator, would be appointed upon authorization by the Senate of the Roman Republic. Most magistrates were elected annually for a term of one year. The terms for all annual offices would begin on New Year's Day, and end on the last day of December. The two highest ranking ordinary magistrates, the consuls and praetors, held a type of authority called imperium (Latin for "command"). Imperium allowed a magistrate to command a military force. Consuls held a higher grade of imperium than praetors. Consuls and praetors, as well as censors and curule aediles, were regarded as "curule magistrates". They would sit on a curule chair, which was a symbol of state power. Consuls and praetors where attended by bodyguards called lictors. The lictors would carry fasces. The fasces, which consisted of a rod with an embedded axe, were symbols of the coercive power of the state. Quaestors were not curule magistrates, and had little real power. Plebeian tribunes were not officially "magistrates", since they were elected only by the plebeians. Since they were considered to be the embodiment of the People of Rome, their office and their person were considered sacrosanct. It was considered to be a capital offense to harm a tribune, to attempt to harm a tribune, or to attempt to obstruct a tribune in any way. All other powers of the tribunate derived from this sacrosanctity. The tribunes were assisted by plebeian aediles. In an emergency, a dictator would be appointed. A newly appointed dictator would usually select a deputy, known as the "Magister Equitum" ("Master of the Horse"). Both the dictator and his master of the horse were extraordinary magistrates, and they both held imperium. In practice, the dictator functioned as a consul without any constitutional checks on his power. After 202 BC, the dictatorship fell into disuse. During emergencies, the senate would pass the senatus consultum ultimum ("ultimate decree of the senate"). This suspended civil government, and declared (something analogous to) martial law. It would declare "videant consules ne res publica detrimenti capiat" ("let the consuls see to it that the state suffer no harm"). In effect, the consuls would be vested with dictatorial powers. After the fall of the republic, the old magistracies (dictators, consuls, praetors, censors, aediles, quaestors and tribunes) were either outright abandoned, or simply lost all powers. The emperor became the master of the state. The founding of the empire was tantamount to a restoration of the old monarchy. The chief executive became the unchallenged power in the state, the senate became a powerless advisory council, and the assemblies became irrelevant. The legacy of the Roman constitution The Roman constitution was one of the few constitutions to exist before the 18th century. None of the others are as well known to us today. And none of the others governed such a vast empire for so long. Therefore, the Roman constitution was used as a template, often the only one, when the first constitutions of the modern era were being drafted. And because of this, many modern constitutions share a similar, even identical, superstructure (such as a separation of powers and checks and balances) as did the Roman constitution. - Abbott, Frank Frost (1901). A History and Description of Roman Political Institutions. Elibron Classics (ISBN 0-543-92749-0). - Byrd, Robert (1995). The Senate of the Roman Republic. U.S. Government Printing Office, Senate Document 103-23. - Cicero, Marcus Tullius (1841). The Political Works of Marcus Tullius Cicero: Comprising his Treatise on the Commonwealth; and his Treatise on the Laws. Translated from the original, with Dissertations and Notes in Two Volumes. By Francis Barham, Esq. London: Edmund Spettigue. Vol. 1. - Lintott, Andrew (1999). The Constitution of the Roman Republic. Oxford University Press (ISBN 0-19-926108-3). - Polybius (1823). The General History of Polybius: Translated from the Greek. By James Hampton. Oxford: Printed by W. Baxter. Fifth Edition, Vol 2. - Taylor, Lily Ross (1966). Roman Voting Assemblies: From the Hannibalic War to the Dictatorship of Caesar. The University of Michigan Press (ISBN 0-472-08125-X). - Byrd, 161 - Ihne, Wilhelm. Researches Into the History of the Roman Constitution. William Pickering. 1853. - Johnston, Harold Whetstone. Orations and Letters of Cicero: With Historical Introduction, An Outline of the Roman Constitution, Notes, Vocabulary and Index. Scott, Foresman and Company. 1891. - Mommsen, Theodor. Roman Constitutional Law. 1871-1888 - Tighe, Ambrose. The Development of the Roman Constitution. D. Apple & Co. 1886. - Von Fritz, Kurt. The Theory of the Mixed Constitution in Antiquity. Columbia University Press, New York. 1975. - The Histories by Polybius - Cambridge Ancient History, Volumes 9–13. - A. Cameron, The Later Roman Empire, (Fontana Press, 1993). - M. Crawford, The Roman Republic, (Fontana Press, 1978). - E. S. Gruen, "The Last Generation of the Roman Republic" (U California Press, 1974) - F. Millar, The Emperor in the Roman World, (Duckworth, 1977, 1992). - A. Lintott, "The Constitution of the Roman Republic" (Oxford University Press, 1999) - Cicero's De Re Publica, Book Two - Rome at the End of the Punic Wars: An Analysis of the Roman Government; by Polybius Secondary source material - Considerations on the Causes of the Greatness of the Romans and their Decline, by Montesquieu - The Roman Constitution to the Time of Cicero - What a Terrorist Incident in Ancient Rome Can Teach Us

https://en.wikipedia.org/wiki/Roman_Constitution

Beacon Lesson Plan Library Leon County Schools After reading the novel FREAK THE MIGHTY students will be able to describe and illustrate the setting of the novel, explain character development through production of a graphic organizer, and identify the elements of the plot. The student describes or illustrates the setting in a literary text. The student explains character development in a literary text. The student creates a graphic organizer that represents the complex elements of a plot in a literary text. -FREAK THE MIGHTY (Rodman Philbrick,Scholastic, 2001) -Legal size white paper -One copy of novel for each student -One copy of the rubric for each student (SEE ASSOCIATED FILES) 1. Order enough copies of FREAK THE MIGHTY for each student (books can be ordered from Amazon.com or Scholastic) 2. Copy One-Pager form onto transparency. 3. Set up overhead projector 4. Assemble necessary materials (paper, markers, crayons, rulers) 5. Copy assessment rubric for distribution to students (see associated file) Students should have read the novel FREAK THE MIGHTY before initiating this lesson. 1. Review story elements through a class discussion. Discussion points include: plot, setting, and character development. Discuss with the students examples of these elements based on the novel FREAK THE MIGHTY. Possible questions could be as follows: PLOT: What is plot? What are some examples of the plot from FREAK THE MIGHTY? What is the main conflict? How is the conflict resolved? SETTING: What is setting? Where does the story take place? How do we know? What are the clues in the novel that help us determine the setting of the story? If you had to illustrate the setting, what do you think it would look like? CHARACTER DEVELOPMENT: Who are the main characters? What do we know about our main characters? How do we know? What are the experiences they go through? What do they look like? How do the different characters deal with the conflicts in the novel? 2. Address any questions or concerns that the students may have. Discuss the following: Re-define the term conflict. Re-define the term resolution. How would the novel be different if the characters personalities were swapped? What do you think would happen next if you could write the next chapter? 3. Through a class discussion review the ideas of a road map, flow chart, and sequencing chart. Students work individually to develop graphic organizers to help them map out the plot of the story. Students work on scrap paper to develop ideas. Let students be as creative as possible. Periodically point out exemplary examples/ideas of graphic organizers as you observe their work. 4. Allow students to share their ideas with their classmates in either pairs or groups of four to five. Groups report back to class. 5. Review the elements discussed today. Tell the students tomorrow they will be working on demonstrating their knowledge of these elements through designing their own One Pagers (show example on overhead projector). Days 2 & 3 1. Review the elements of plot, setting, and character development. Address any questions that the students may have. 2. Introduce the One Pager idea to the students (transparency / see associated file). 3. Pass out the scoring rubric to students (see associated file). 4. Instruct students on what is expected of them. They are to complete the One Pager using the transparency and the rubric as a guide. 5. Pass out legal size paper and supplies to students. 6. Instruct students that they will have two days to complete this activity. 7. When students finish this activity, have them turn in their One Pagers with their scoring rubric (see associated file). 8. Score students using the rubric and return the One Pagers to students. Assessment of studentsí One Pagers will be based on the rubric attached in Associated Files. Teachers' guide to FREAK THE MIGHTY including discussion questions. This site also contains links to the author and the movie “The Mighty” produced by Miramax, 1999.Teachers' guide to FREAK THE MIGHTY

http://www.beaconlearningcenter.com/lessons/lesson.asp?ID=2238

|This article needs additional citations for verification. (January 2016)| A family of particular transformations may be continuous (such as rotation of a circle) or discrete (e.g., reflection of a bilaterally symmetric figure, or rotation of a regular polygon). Continuous and discrete transformations give rise to corresponding types of symmetries. Continuous symmetries can be described by Lie groups while discrete symmetries are described by finite groups (see Symmetry group). These two concepts, Lie and finite groups, are the foundation for the fundamental theories of modern physics. Symmetries are frequently amenable to mathematical formulations such as group representations and can, in addition, be exploited to simplify many problems. Arguably the most important example of a symmetry in physics is that the speed of light has the same value in all frames of reference, which is known in mathematical terms as Poincare group, the symmetry group of special relativity. Another important example is the invariance of the form of physical laws under arbitrary differentiable coordinate transformations, which is an important idea in general relativity. - 1 Symmetry as invariance - 2 Local and global symmetries - 3 Continuous symmetries - 4 Discrete symmetries - 5 Mathematics of physical symmetry - 6 Mathematics - 7 See also - 8 References - 9 External links Symmetry as invariance Invariance is specified mathematically by transformations that leave some quantity unchanged. This idea can apply to basic real-world observations. For example, temperature may be constant throughout a room. Since the temperature is independent of position within the room, the temperature is invariant under a shift in the measurer's position. Similarly, a uniform sphere rotated about its center will appear exactly as it did before the rotation. The sphere is said to exhibit spherical symmetry. A rotation about any axis of the sphere will preserve how the sphere "looks". Invariance in force The above ideas lead to the useful idea of invariance when discussing observed physical symmetry; this can be applied to symmetries in forces as well. For example, an electric field due to a wire is said to exhibit cylindrical symmetry, because the electric field strength at a given distance r from the electrically charged wire of infinite length will have the same magnitude at each point on the surface of a cylinder (whose axis is the wire) with radius r. Rotating the wire about its own axis does not change its position or charge density, hence it will preserve the field. The field strength at a rotated position is the same. Suppose some configuration of charges (may be non-stationary) produce an electric field in some direction, then rotating the configuration of the charges (without disturbing the internal dynamics that produces the particular field) will lead to a net rotation of the direction of the electric field. These two properties are interconnected through the more general property that rotating any system of charges causes a corresponding rotation of the electric field. In Newton's theory of mechanics, given two bodies, each with mass m, starting from rest at the origin and moving along the x-axis in opposite directions, one with speed v1 and the other with speed v2 the total kinetic energy of the system (as calculated from an observer at the origin) is 1⁄2m(v12 + v22) and remains the same if the velocities are interchanged. The total kinetic energy is preserved under a reflection in the y-axis. The last example above illustrates another way of expressing symmetries, namely through the equations that describe some aspect of the physical system. The above example shows that the total kinetic energy will be the same if v1 and v2 are interchanged. Local and global symmetries Symmetries may be broadly classified as global or local. A global symmetry is one that holds at all points of spacetime, whereas a local symmetry is one that has a different symmetry transformation at different points of spacetime; specifically a local symmetry transformation is parameterised by the spacetime co-ordinates. Local symmetries play an important role in physics as they form the basis for gauge theories. The two examples of rotational symmetry described above - spherical and cylindrical - are each instances of continuous symmetry. These are characterised by invariance following a continuous change in the geometry of the system. For example, the wire may be rotated through any angle about its axis and the field strength will be the same on a given cylinder. Mathematically, continuous symmetries are described by continuous or smooth functions. An important subclass of continuous symmetries in physics are spacetime symmetries. Continuous spacetime symmetries are symmetries involving transformations of space and time. These may be further classified as spatial symmetries, involving only the spatial geometry associated with a physical system; temporal symmetries, involving only changes in time; or spatio-temporal symmetries, involving changes in both space and time. - Time translation: A physical system may have the same features over a certain interval of time ; this is expressed mathematically as invariance under the transformation for any real numbers t and a in the interval. For example, in classical mechanics, a particle solely acted upon by gravity will have gravitational potential energy when suspended from a height above the Earth's surface. Assuming no change in the height of the particle, this will be the total gravitational potential energy of the particle at all times. In other words, by considering the state of the particle at some time (in seconds) and also at , say, the particle's total gravitational potential energy will be preserved. - Spatial translation: These spatial symmetries are represented by transformations of the form and describe those situations where a property of the system does not change with a continuous change in location. For example, the temperature in a room may be independent of where the thermometer is located in the room. - Spatial rotation: These spatial symmetries are classified as proper rotations and improper rotations. The former are just the 'ordinary' rotations; mathematically, they are represented by square matrices with unit determinant. The latter are represented by square matrices with determinant −1 and consist of a proper rotation combined with a spatial reflection (inversion). For example, a sphere has proper rotational symmetry. Other types of spatial rotations are described in the article Rotation symmetry. - Poincaré transformations: These are spatio-temporal symmetries which preserve distances in Minkowski spacetime, i.e. they are isometries of Minkowski space. They are studied primarily in special relativity. Those isometries that leave the origin fixed are called Lorentz transformations and give rise to the symmetry known as Lorentz covariance. - Projective symmetries: These are spatio-temporal symmetries which preserve the geodesic structure of spacetime. They may be defined on any smooth manifold, but find many applications in the study of exact solutions in general relativity. - Inversion transformations: These are spatio-temporal symmetries which generalise Poincaré transformations to include other conformal one-to-one transformations on the space-time coordinates. Lengths are not invariant under inversion transformations but there is a cross-ratio on four points that is invariant. Mathematically, spacetime symmetries are usually described by smooth vector fields on a smooth manifold. The underlying local diffeomorphisms associated with the vector fields correspond more directly to the physical symmetries, but the vector fields themselves are more often used when classifying the symmetries of the physical system. Some of the most important vector fields are Killing vector fields which are those spacetime symmetries that preserve the underlying metric structure of a manifold. In rough terms, Killing vector fields preserve the distance between any two points of the manifold and often go by the name of isometries. A discrete symmetry is a symmetry that describes non-continuous changes in a system. For example, a square possesses discrete rotational symmetry, as only rotations by multiples of right angles will preserve the square's original appearance. Discrete symmetries sometimes involve some type of 'swapping', these swaps usually being called reflections or interchanges. - Time reversal: Many laws of physics describe real phenomena when the direction of time is reversed. Mathematically, this is represented by the transformation, . For example, Newton's second law of motion still holds if, in the equation , is replaced by . This may be illustrated by recording the motion of an object thrown up vertically (neglecting air resistance) and then playing it back. The object will follow the same parabolic trajectory through the air, whether the recording is played normally or in reverse. Thus, position is symmetric with respect to the instant that the object is at its maximum height. - Spatial inversion: These are represented by transformations of the form and indicate an invariance property of a system when the coordinates are 'inverted'. Said another way, these are symmetries between a certain object and its mirror image. - Glide reflection: These are represented by a composition of a translation and a reflection. These symmetries occur in some crystals and in some planar symmetries, known as wallpaper symmetries. C, P, and T symmetries - Every particle is replaced with its antiparticle. This is C-symmetry (charge symmetry); - Everything appears as if reflected in a mirror. This is P-symmetry (parity symmetry); T-symmetry is counterintuitive (surely the future and the past are not symmetrical) but explained by the fact that the Standard model describes local properties, not global ones like entropy. To properly reverse the direction of time, one would have to put the big bang and the resulting low-entropy state in the "future." Since we perceive the "past" ("future") as having lower (higher) entropy than the present (see perception of time), the inhabitants of this hypothetical time-reversed universe would perceive the future in the same way as we perceive the past. These symmetries are near-symmetries because each is broken in the present-day universe. However, the Standard Model predicts that the combination of the three (that is, the simultaneous application of all three transformations) must be a symmetry, called CPT symmetry. In the 4 dimensional matrix description of P,T is through a diagonal matrix, the negative identity, as well as C. Hence CPT is the identity operator. CP violation, the violation of the combination of C- and P-symmetry, is necessary for the presence of significant amounts of baryonic matter in the universe. CP violation is a fruitful area of current research in particle physics. ||This section may contain misleading parts. (June 2015)| A type of symmetry known as supersymmetry has been used to try to make theoretical advances in the standard model. Supersymmetry is based on the idea that there is another physical symmetry beyond those already developed in the standard model, specifically a symmetry between bosons and fermions. Supersymmetry asserts that each type of boson has, as a supersymmetric partner, a fermion, called a superpartner, and vice versa. Supersymmetry has not yet been experimentally verified: no known particle has the correct properties to be a superpartner of any other known particle. If superpartners exist they must have masses greater than current particle accelerators can generate. Mathematics of physical symmetry Continuous symmetries are specified mathematically by continuous groups (called Lie groups). Many physical symmetries are isometries and are specified by symmetry groups. Sometimes this term is used for more general types of symmetries. The set of all proper rotations (about any angle) through any axis of a sphere form a Lie group called the special orthogonal group . (The 3 refers to the three-dimensional space of an ordinary sphere.) Thus, the symmetry group of the sphere with proper rotations is . Any rotation preserves distances on the surface of the ball. The set of all Lorentz transformations form a group called the Lorentz group (this may be generalised to the Poincaré group). Discrete symmetries are described by discrete groups. For example, the symmetries of an equilateral triangle are described by the symmetric group . An important type of physical theory based on local symmetries is called a gauge theory and the symmetries natural to such a theory are called gauge symmetries. Gauge symmetries in the Standard model, used to describe three of the fundamental interactions, are based on the SU(3) × SU(2) × U(1) group. (Roughly speaking, the symmetries of the SU(3) group describe the strong force, the SU(2) group describes the weak interaction and the U(1) group describes the electromagnetic force.) Also, the reduction by symmetry of the energy functional under the action by a group and spontaneous symmetry breaking of transformations of symmetric groups appear to elucidate topics in particle physics (for example, the unification of electromagnetism and the weak force in physical cosmology). Conservation laws and symmetry The symmetry properties of a physical system are intimately related to the conservation laws characterizing that system. Noether's theorem gives a precise description of this relation. The theorem states that each continuous symmetry of a physical system implies that some physical property of that system is conserved. Conversely, each conserved quantity has a corresponding symmetry. For example, the isometry of space gives rise to conservation of (linear) momentum, and isometry of time gives rise to conservation of energy. The following table summarizes some fundamental symmetries and the associated conserved quantity. |translation in time |translation in space |rotation in space |Discrete symmetry||P, coordinate inversion||spatial parity| |C, charge conjugation||charge parity| |T, time reversal||time parity| |CPT||product of parities| |Internal symmetry (independent of |U(1) gauge transformation||electric charge| |U(1) gauge transformation||lepton generation number| |U(1) gauge transformation||hypercharge| |U(1)Y gauge transformation||weak hypercharge| |U(2) [ U(1) × SU(2) ]||electroweak force| |SU(2) gauge transformation||isospin| |SU(2)L gauge transformation||weak isospin| |P × SU(2)||G-parity| |SU(3) "winding number"||baryon number| |SU(3) gauge transformation||quark color| |SU(3) (approximate)||quark flavor| |S(U(2) × U(3)) [ U(1) × SU(2) × SU(3) ] Continuous symmetries in physics preserve transformations. One can specify a symmetry by showing how a very small transformation affects various particle fields. The commutator of two of these infinitesimal transformations are equivalent to a third infinitesimal transformation of the same kind hence they form a Lie algebra. for a general field, . Without gravity only the Poincaré symmetries are preserved which restricts to be of the form: where M is an antisymmetric matrix (giving the Lorentz and rotational symmetries) and P is a general vector (giving the translational symmetries). Other symmetries affect multiple fields simultaneously. For example local gauge transformations apply to both a vector and spinor field: where are generators of a particular Lie group. So far the transformations on the right have only included fields of the same type. Supersymmetries are defined according to how the mix fields of different types. Another symmetry which is part of some theories of physics and not in others is scale invariance which involve Weyl transformations of the following kind: If the fields have this symmetry then it can be shown that the field theory is almost certainly conformally invariant also. This means that in the absence of gravity h(x) would restricted to the form: with D generating scale transformations and K generating special conformal transformations. For example N=4 super-Yang-Mills theory has this symmetry while General Relativity doesn't although other theories of gravity such as conformal gravity do. The 'action' of a field theory is an invariant under all the symmetries of the theory. Much of modern theoretical physics is to do with speculating on the various symmetries the Universe may have and finding the invariants to construct field theories as models. In string theories, since a string can be decomposed into an infinite number of particle fields, the symmetries on the string world sheet is equivalent to special transformations which mix an infinite number of fields. - Conservation law - Conserved current - Covariance and contravariance - Fictitious force - Galilean invariance - Gauge theory - General covariance - Harmonic coordinate condition - Inertial frame of reference - Lie group - List of mathematical topics in relativity - Lorentz covariance - Noether's theorem - Poincaré group - Special relativity - Spontaneous symmetry breaking - Standard model - Standard model (mathematical formulation) - Symmetry breaking - Wheeler–Feynman Time-Symmetric Theory - G. Kalmbach H.E.: Quantum Mathematics: WIGRIS. RGN Publications, Delhi, 2014. - Leon Lederman and Christopher T. Hill (2005) Symmetry and the Beautiful Universe. Amherst NY: Prometheus Books. - Schumm, Bruce (2004) Deep Down Things. Johns Hopkins Univ. Press. - Victor J. Stenger (2000) Timeless Reality: Symmetry, Simplicity, and Multiple Universes. Buffalo NY: Prometheus Books. Chpt. 12 is a gentle introduction to symmetry, invariance, and conservation laws. - Anthony Zee (2007) Fearful Symmetry: The search for beauty in modern physics, 2nd ed. Princeton University Press. ISBN 978-0-691-00946-9. 1986 1st ed. published by Macmillan. - Brading, K., and Castellani, E., eds. (2003) Symmetries in Physics: Philosophical Reflections. Cambridge Univ. Press. - -------- (2007) "Symmetries and Invariances in Classical Physics" in Butterfield, J., and John Earman, eds., Philosophy of Physic Part B. North Holland: 1331-68. - Debs, T. and Redhead, M. (2007) Objectivity, Invariance, and Convention: Symmetry in Physical Science. Harvard Univ. Press. - John Earman (2002) "Laws, Symmetry, and Symmetry Breaking: Invariance, Conservations Principles, and Objectivity." Address to the 2002 meeting of the Philosophy of Science Association. - G. Kalmbach H.E.: Quantum Mathematics: WIGRIS. RGN Publications, Delhi, 2014 - Mainzer, K. (1996) Symmetries of nature. Berlin: De Gruyter. - Mouchet, A. "Reflections on the four facets of symmetry: how physics exemplifies rational thinking". European Physical Journal H 38 (2013) 661 hal.archives-ouvertes.fr:hal-00637572 - Thompson, William J. (1994) Angular Momentum: An Illustrated Guide to Rotational Symmetries for Physical Systems. Wiley. ISBN 0-471-55264-X. - Bas Van Fraassen (1989) Laws and symmetry. Oxford Univ. Press. - Eugene Wigner (1967) Symmetries and Reflections. Indiana Univ. Press.

https://en.wikipedia.org/wiki/Symmetry_in_physics

Lesson 16: Writing to a text file In the previous lesson, we learned to read from a text file. In this lesson, we will learn to write to a text file. The two methods are very similar, but there is one very important difference: You must have write permissions to the file. This means that the file will have to be located in a folder where you have the necessary permissions. If you work locally on your own computer, you can set the permissions yourself: right-click on the folder and choose "Properties". With most web hosts, you will normally have one folder with write permissions. It's often called something like "cgi-bin", "log", "databases" or something similar. If your web host permits it, you might also be able to set permissions yourself. Usually you can simply right-click on a folder in your FTP client and choose "properties" or "permissions" or something similar. The screendumps below shows how it's done in FileZilla. Read more on your web host's support pages. Note that it is the text file that needs to be in the folder with write permissions - not the PHP file. Open the text file for writing In the same way as when reading from a text file, the fopen function is used for writing, but this time we set the mode to "w" (writing) or "a" (appending). The difference between writing and appending is where the 'cursor' is located - either at the beginning or at the end of the text file. The examples in this lesson use an empty text file called textfile.txt. But you can also create your own text file if you like. First, let us try to open the text file for writing: <?php // Open the text file $f = fopen("textfile.txt", "w"); // Close the text file fclose($f); ?> Example 1: Write a line to the text file To write a line, we must use the function fwrite, like this: <html> <head> <title>Writing to a text file</title> </head> <body> <?php // Open the text file $f = fopen("textfile.txt", "w"); // Write text line fwrite($f, "PHP is fun!"); // Close the text file fclose($f); // Open file for reading, and read the line $f = fopen("textfile.txt", "r"); echo fgets($f); fclose($f); ?> </body> </html> Since we opened the file for writing, the line is added at the top, and thus overwrites the existing line. If we instead open the file appending, the line is added at the bottom of the text file, which then will increase by one line each time it's written to. Example 2: Adding a text block to a text file Of course, it is also possible to add an entire text block, instead of just a single line, like this: <html> <head> <title>Write to a text file</title> </head> <body> <h1>Adding a text block to a text file:</h1> <form action="myfile.php" method='post'> <textarea name='textblock'></textarea> <input type='submit' value='Add text'> </form> <?php // Open the text file $f = fopen("textfile.txt", "w"); // Write text fwrite($f, $_POST["textblock"]); // Close the text file fclose($f); // Open file for reading, and read the line $f = fopen("textfile.txt", "r"); // Read text echo fgets($f); fclose($f); ?> </body> </html> In the next lessons, we look at another way of storing data: databases. |Related topics in the PHP Forum||Replies||Views| |how to create a 'secure'login system using php?||0||173| |PHP Code for retrieve and write data from a table||2||5492| |Can I use PHP in CSS style file?||6||32772| |PHP Code for retrieve and write||1||5017| |How write a SCHEDULING SCRIPT(autoposting) WITH PHP||2||4142|

http://html.net/tutorials/php/lesson16.php

January 14, 2013 Study Reveals First Ever Images Of Early Tetrapod Backbone And How It Helped In Land Evolution [Watch Video: 3D X-Ray Images Of Early Tetrapod Backbone] Lawrence LeBlond for redOrbit.com - Your Universe OnlineUsing high-energy X-rays and a new data extraction protocol, an international consortium of scientists have for the first time rendered a 3D model of a prehistoric tetrapod backbone. The new reconstruction has shed new light on how the early animals moved once they made it onto land. One of the main creatures studied was a fierce-looking ichthyostega that lived from 374 — 359 million years ago and was a transitional species between aquatic and terrestrial animals. The 3D model showed researchers that these new land dwellers moved much like modern seals do. The researchers believe ichthyostega was more of a shallow water predator, navigating swamps and ponds in search of food, occasionally making landfall to perhaps feed. The researchers think the animal dragged itself across flat ground, using its front legs to push up and forward. Results of the new study have been published in this week´s issue of the journal Nature. The international study was led by Stephanie E. Pierce from The Royal Veterinary College in London and Jennifer A. Clack from the University of Cambridge. Other members of the team hailed from Sweden and France. Tetrapods are four-limbed vertebrates. In our modern world, animals such as amphibians, reptiles, birds, and mammals are all tetrapods. Early tetrapods, such as the ichthyostega, made short excursions across shallow bodies of water and perhaps even shorter jaunts over land, using their underdeveloped limbs for primitive locomotion. Just how these early tetrapods transitioned from a life at sea to land-dwelling has been a hotly debated topic among paleontologists and evolution biologists for decades. Not only do all tetrapods have four limbs, but they have backbones (vertebral column) as well. These vertebrates also include fish, from which tetrapods evolved. The backbone is formed from vertebrae connected in a row–from head to tail. But unlike the backbone of modern tetrapods, in which each vertebra is composed of only one bone, early tetrapods had vertebrae made up of multiple parts. “For more than 100 years, early tetrapods were thought to have vertebrae composed of three sets of bones - one bone in front, one on top, and a pair behind,” said Pierce. “But, by peering inside the fossils using synchrotron X-rays we have discovered that this traditional view literally got it back-to-front.” “The results of this study force us to re-write the textbook on backbone evolution in the earliest limbed animals,” Discovery News quoted Pierce as saying. To make their analysis, the team relied on the European Synchrotron Radiation Facility (ESRF) in France to scan three fossil fragments of early tetrapods. Using the X-ray scanner, details began to emerge of the fossil bones buried deep inside the rock matrix. Although the rock obscured most of the X-rays, the team was able to decipher the readings using a detailed data extraction method. "Without the new method, it would not have been possible to reveal the elements of the spine in three dimensions with a resolution of 30 micrometres" noted study coauthor Sophie Sanchez from University of Uppsala and ESRF. Between the X-ray images and the data extraction tools, the team discovered that what they believed to be the first bone (the intercentrum) was actually the last in the series. The team said that this revelation brings new insight into how the vertebral structure plays out for the functional evolution of the tetrapod backbone. "By understanding how each of the bones fit together we can begin to explore the mobility of the spine and test how it may have transferred forces between the limbs during the early stages of land movement,” noted Pierce. Aside from the backbone discovery, the team also found that the ichthyostega also had an unusual assortment of previously unknown skeletal formations including a string of bones extending down the middle of its chest. “These chest bones turned out to be the earliest evolutionary attempt to produce a bony sternum. Such a structure would have strengthened the ribcage of Ichthyostega, permitting it to support its body weight on its chest while moving about on land,” Clack explained. In continuing their research, the team said the next phase will be to further investigate how the backbone aided in the locomotion of these early tetrapodous animals. Image below shows an artist's impression of an Ichthyostega Tetrapod, with the cut-out showing the 3-D reconstruction of two vetrebrae from the study. Credit: Julia Molnar

http://www.redorbit.com/news/science/1112763058/tetrapod-backbone-early-land-evolution-011413/

People often wonder how delicate arches and finely balanced pillars of stone stand up to the stress of holding up their own immense weight. Actually, new research suggests, it’s that stress that helps pack individual grains of sand together and slows erosion of the formations. In lab experiments, scientists dropped small blocks of loosely consolidated sandstone into water—and watched them completely fall apart as the water dissolved minerals holding the grains together. But when the scientists placed weights on top of the sandstone samples before submersing them, disintegration ceased once stress in the eroding column rose to a certain threshold that packed the sand grains into a strong, rocklike material, the researchers report online today in Nature Geoscience. In other tests, weight-induced stress similarly protected samples against complete erosion from simulated rainfall. At large scale in the real world, stress transmitted through arches and pillars to their bases (in landforms such as Delicate Arch in Utah’s Arches National Park, shown) slows down—but doesn’t stop—natural sculpting due to wind and water, the researchers say. Bits of the landform that don’t bear weight are among the first to wear away, which helps explain why arches are often unusually smooth. Cracks, fissures, and soft layers in rock formations influence the shapes these natural sculptures take as they evolve.

http://www.sciencemag.org/news/2014/07/what-keeps-stone-arches-falling-down

Alphabet Letter S Impact Poster Preschool Lesson Plan Activities Alphabet > Letter S > Impact Letter S Poster Note: This activity counts as two activities due to the steps involved. Version for 2 & 3 year olds: *Step 1: Adult draw upper case and lower case letter on the large piece of paper (ahead of time) and print and cut the four letter pictures. Proceed with the steps below. Version for 3 & 4’s: *Step 1: Adult draws large dots to form the letters (include arrows to show direction). Child will connect the dots with paint. Proceed with the steps below. *This activity is better accomplished when the poster is placed on a wall, poster board or easel. *Encourage child to paint the letter. Explain the proper direction of “writing” the letter. Join child in the painting process. *While the paint dries, discuss some facts about the letter pictures and conduct a short activity or craft for each image, ask child to find the hidden S's in the images, emphasize on the letter sound. * Child will glue one or two pictures a day until all eight letter cards are on the poster. *Display the poster at the child’s eye level in a place he will see regularly (kitchen wall, child’s room, family room) at least for a week. *Remember to make a picture of child with his completed first letter project for his scrapbook. *Large piece of paper or tape tape four standard size to make a poster *Paint brush or sponge. Print and cut these letter pictures for the poster Letter S Impact flash cards |To view updates to these activities visit: http://www.first-school.ws/activities/alpha/s/impactsposter.htm|

http://first-school.ws/activities/alpha/s/impactsposter.htm

Super Plastic Both Attracts and Repels Water An odd new material could be a boon in dry regions with limited access to clean water. A new, practical method for making surfaces with patterns of areas that strongly attract and strongly repel water could lead to a highly efficient method for capturing clean water. This versatile material could also find uses in fabricating new types of devices for medical tests and chemical synthesis. Scientists have reported numerous applications of water-attracting (superhydrophilic) and water-repelling (superhydrophobic) surfaces, including fog-free eyeglasses and windshields, and self-cleaning cloth and glass. Now a group of researchers in MIT’s materials science and engineering department has combined those opposing characteristics on a single surface, by using a simple and versatile fabrication process. [For images of this new dual-quality material, click here.] Robert Cohen, Michael Rubner, and colleagues started by assembling a nano-structured film made of alternating layers of positively and negatively charged polymers and silica nanoparticles. The film’s structure and a coating of waxy fluorinated silane cause water to bead on it, forming near-perfect spheres that easily roll off. To add the superhydrophilic regions (to which water droplets cling), the researchers applied a naturally hydrophilic polymer to selected areas. In dry regions of the world, without easy access to clean water, such a material could be used for collecting water. In this application, the hydrophilic areas of the material would attract moisture in the air, collecting water drops that accumulate, until they spill over into the hydrophobic regions and roll into a collecting channel. Currently, in countries with limited access to clean water, the inhabitants typically use large polypropylene fiber meshes to harvest water from fog. The new technology “would provide a more than tenfold increase in water capture compared to the inefficient nets that are used currently,” says Andrew Parker, a biologist at Oxford University and the Natural History Museum in London, who has studied the desert beetle that inspired the MIT work. If the new material “could be added simply to the roofs of houses in areas subjected to desert fogs,” says Parker, “then a water supply could be gained with little effort.” Rubner’s lab is also taking the technique further. “When we harvest water, we have chemistry built into the hydrophilic area so that it has an antibacterial agent to kill off bacteria and other things that cause harm,” Rubner says. This decontaminates the water as it accumulates so that the collected water is safe for use. Applying this technique, the researchers have been able to kill common harmful bacteria in four minutes, he says. The coating could also find uses in biomedical applications to make microfluidic chips. Typically, microfluidic devices contain enclosed micrometer-wide channels etched into silicon, glass, or plastic plates. Then pressure or electric fields drive tiny volumes of fluids, typically nanoliters, along these channels for diagnostic tests and genetics research. For instance, to test for the presence of a certain protein in blood you could take blood in one channel and direct it to another channel containing a chemical reagent that identifies the protein. Compared with conventional microfluidics, a microfluidic chip based on the new surface would have the advantage of easier mixing, Rubner says. Right now, the chips need pumps and valves that move the liquid around to induce mixing. “In our case you can mix the liquids by just controlling the amount of liquid you put on the surface,” he says. With a pipette, you could add precise amounts of fluid into two hydrophilic grooves placed close to each other. As you add more fluid, the droplets bulge out at the edge of the grooves because of the surrounding hydrophobic area. Eventually, the bulging surfaces touch and mix. Being able to confine liquids to a small region could provide densely packed reaction sites with more control over the reaction, he says, since adjacent drops won’t mix unless they are forced to. While the exact uses of this new material are still uncertain, it opens up many possibilities, says Kenneth Wynne, a chemical engineering professor at Virginia Commonwealth University. “Patterning ultra-hydrophilic patches on a ultra-hydrophobic surface in this way is new and useful,” he says.

https://www.technologyreview.com/s/405883/super-plastic-both-attracts-and-repels-water/

Pharisees (fârˈĭsēz) [key], one of the two great Jewish religious and political parties of the second commonwealth. Their opponents were the Sadducees, and it appears that the Sadducees gave them their name, perushim, Hebrew for "separatists" or "deviants." The Pharisees began their activities during or after the Hasmonean revolt (c.166–142 B.C.). The Pharisees upheld an interpretation of Judaism that was in opposition to the priestly Temple cult. They stressed faith in the one God; the divine revelation of the law both written and oral handed down by Moses through Joshua, the elders, and the prophets to the Pharisees; and eternal life and resurrection for those who keep the law. Pharisees insisted on the strict observance of Jewish law, which they began to codify. While in agreement on the broad outlines of Jewish law, the Pharisees encouraged debate on its fine points, and according to one view, practiced the tradition of zuggot, or pairs of scholars with opposing views. They developed the synagogue as an alternative place of worship to the Temple, with a liturgy consisting of biblical and prophetic readings, and the repetition of the shma, the basic creed of Judaism. In addition, they supported the separation of the worldly and the spiritual spheres, ceding the former to the secular rulers. Though some supported the revolt against Rome in A.D. 70, most did not. One Pharisee was Yohanan ben Zakkai, who fled to Jamnia, where he was instrumental in developing post-Temple Judaism. By separating Judaism from dependence on the Temple cult, and by stressing the direct relation between the individual and God, the Pharisees laid the groundwork for normative rabbinic Judaism. Their influence on Christianity was substantial as well, despite the passages in the New Testament which label the Pharisees "hypocrites" or "offspring of the vipers." St. Paul was originally a Pharisee. After the fall of the Temple (A.D. 70), the Pharisees became the dominant party until c.135. See L. Finkelstein, The Pharisees: The Sociological Background of Their Faith (3d ed., 2 vol., 1963); A. Finkel, The Pharisees and the Teacher of Nazareth (1964); L. Baeck, Pharisees (1947, repr. 1966); J. Neusner, From Politics to Piety (1973) and The Pharisees (1985). The Columbia Electronic Encyclopedia, 6th ed. Copyright © 2012, Columbia University Press. All rights reserved.

http://www.factmonster.com/encyclopedia/society/pharisees.html

Voltage drop describes how the supplied energy of a voltage source is reduced as electric current moves through the passive elements (elements that do not supply voltage) of an electrical circuit. Voltage drops across internal resistances of the source, across conductors, across contacts, and across connectors are undesired; supplied energy is lost (dissipated). Voltage drops across loads and across other active circuit elements are desired; supplied energy performs useful work. For example, an electric space heater may have a resistance of ten ohms, and the wires which supply it may have a resistance of 0.2 ohms, about 2% of the total circuit resistance. This means that approximately 2% of the supplied voltage is lost in the wire itself. Excessive voltage drop may result in unsatisfactory operation of, and damage to, electrical and electronic equipment. National and local electrical codes may set guidelines for the maximum voltage drop allowed in electrical wiring, to ensure efficiency of distribution and proper operation of electrical equipment. The maximum permitted voltage drop varies from one country to another. In electronic design and power transmission, various techniques are employed to compensate for the effect of voltage drop on long circuits or where voltage levels must be accurately maintained. The simplest way to reduce voltage drop is to increase the diameter of the conductor between the source and the load, which lowers the overall resistance. In power distribution systems, a given amount of power can be transmitted with less voltage drop if a higher voltage is used. More sophisticated techniques use active elements to compensate for the undesired voltage drop. Voltage drop in direct-current circuits: resistance Consider a direct-current circuit with a nine-volt DC source; three resistors of 67 ohms, 100 ohms, and 470 ohms; and a light bulb—all connected in series. The DC source, the conductors (wires), the resistors, and the light bulb (the load) all have resistance; all use and dissipate supplied energy to some degree. Their physical characteristics determine how much energy. For example, the DC resistance of a conductor depends upon the conductor's length, cross-sectional area, type of material, and temperature. If the voltage between the DC source and the first resistor (67 ohms) is measured, the voltage potential at the first resistor will be slightly less than nine volts. The current passes through the conductor (wire) from the DC source to the first resistor; as this occurs, some of the supplied energy is "lost" (unavailable to the load), due to the resistance of the conductor. Voltage drop exists in both the supply and return wires of a circuit. If the voltage across each resistor is measured, the measurement will be a significant number. That represents the energy used by the resistor. The larger the resistor, the more energy used by that resistor, and the bigger the voltage drop across that resistor. Ohm's Law can be used to verify voltage drop. In a DC circuit, voltage equals current multiplied by resistance. . Also, Kirchhoff's circuit laws state that in any DC circuit, the sum of the voltage drops across each component of the circuit is equal to the supply voltage. Voltage drop in alternating-current circuits: impedance In alternating-current circuits, opposition to current flow does occur because of resistance (just as in direct-current circuits). Alternating current circuits also present a second kind of opposition to current flow: reactance. This "total" opposition (resistance "plus" reactance) is called impedance. The impedance in an alternating-current circuit depends on the spacing and dimensions of the elements and conductors, the frequency of the alternating current, and the magnetic permeability of the elements, the conductors, and their surroundings. The voltage drop in an AC circuit is the product of the current and the impedance (Z) of the circuit. Electrical impedance, like resistance, is expressed in ohms. Electrical impedance is the vector sum of electrical resistance, capacitive reactance, and inductive reactance. It is expressed by the formula , analogous to Ohm's law for direct-current circuits. - Utility brownout - Voltage divider - Electrical distribution - Electrical resistance - Kirchhoff's voltage law - Electrical conduction - Ground loop (electricity) - Power cable - Mesh analysis - Electrical Principles for the Electrical Trades (Jim Jennesson) 5th edition

https://en.wikipedia.org/wiki/Voltage_drop

This lesson will look at the partisan political issues which emerged in the election of 1864 around Abraham Lincoln's role as a wartime president. Through an examination of primary documents, students will focus on Lincoln's suspension of habeas corpus, the Emancipation Proclamation, his decision to arm the freed slaves, his refusal to accept a compromise peace with the South, and the election of 1864. Popular sovereignty allowed the settlers of a federal territory to decide the slavery question without interference from Congress. This lesson plan will examine how the Kansas–Nebraska Act of 1854 affected the political balance between free and slave states and explore how its author, Stephen Douglas, promoted its policy of popular sovereignty in an effort to avoid a national crisis over slavery in the federal territories. In reviewing events, documentary evidence, and biographical information, students come to understand the complex nature of political decision-making in the United States. In this lesson, they consider the momentous questions facing the country during the Reconstruction debate by weighing the many factors that went into the solutions offered. Students also think critically as they consider whether and how other solutions might have played out. In this lesson, students examine the development of new constitutions in the reconstructed South. They also consider the political and social realities created by a dramatically changed electorate. In gaining a firmer grasp of the causes for the shifting alliances of this time, students see how far-reaching the consequences of the Civil War and Reconstruction era were and how much these events continue to shape our collective destiny today. In this lesson students will learn about Abraham Lincoln the individual and the President. By examining Alexander Gardner's February 5, 1865 photograph and reading a short biography of Lincoln, students will consider who the man on the other side of the lens was. Students will demonstrate their understanding by writing an "I Am" Poem and creating their own multimedia portrait of Lincoln. The focus of this lesson is the Robert Gould Shaw and the Massachusetts 54th Regiment Memorial by Augustus Saint-Gaudens. Students will put themselves in the shoes of the men of the Massachusetts 54th Regiment as they read, write, pose, and then create a comic strip about these American heroes. The newly re-elected Abraham Lincoln sought to unite the American people by interpreting the waning conflict as a divine judgment upon both sides of the war. This lesson will examine Lincoln's Second Inaugural Address to determine how he sought to reunite a divided country through a providential interpretation of the Civil War. This lesson will examine the economic, military and diplomatic strengths and weaknesses of the North and South on the eve of the Civil War. In making these comparisons students will use maps and read original documents to decide which side, if any, had an overall advantage at the start of the war. Abraham Lincoln felt that the attempt of seven states to leave the American union peacefully was, in fact, a total violation of law and order. This lesson will examine Lincoln's First Inaugural Address to understand why he thought his duty as president required him to treat secession as an act of rebellion and not a legitimate legal or constitutional action by disgruntled states. This unit explores the political thought of Abraham Lincoln on the subject of American union. For him, the union was not just a structure to govern the national interests of American states; it also represented a consensus about the future of freedom in America—a future where slavery would eventually be eliminated and liberty protected as the birthright of every human being. Students will examine Lincoln's three most famous speeches—the Gettysburg Address and the First and Second Inaugural Addresses—in addition to a little known fragment on the Constitution, union, and liberty to see what they say regarding the significance of union to the prospects for American self-government. Although Lincoln did not attend high school or college, he possessed a logical and inquisitive mind that found clarity in working out legal and political problems on paper. One fragment he wrote after the 1860 presidential election addressed how the Constitution and union were informed by the ideals of the Declaration of Independence. Lincoln wrote that while America's prosperity was dependent upon the union of the states, "the primary cause" was the principle of "Liberty to all." He believed this central ideal of free government embraced all human beings, and concluded that the American revolution would not have succeeded if its goal was "a mere change of masters." For Lincoln, union meant a particular kind of government of the states, one whose equality principle "clears the path for all—gives hope to all—and, by consequence, enterprize, and industry to all." As president of the United States, Lincoln used his First and Second Inaugural Addresses to explore the meaning of the American union in the face of a divided country. Upon assuming the presidency for the first time, he spoke at length about the nature of union, why secession was antithetical to self-government, and how the federal constitution imposed a duty upon him to defend the union of the states from rebellious citizens. When he was reelected four years later, and as the Civil War drew to a close, Lincoln transcended both Northern triumphalism and Southern defiance by offering a providential reading of the war and emancipation in hopes of reuniting the country. In his most famous speech, delivered upon the dedication of a national cemetery at the battlefield in Gettysburg, Pennsylvania, Lincoln gave a brief but profound meditation on the meaning of the Civil War and American union. With the Emancipation Proclamation as a new and pivotal development of the federal war effort, Lincoln sought to explain why the war to preserve the Union had to become a war to secure the freedom of former slaves. The nation would need to experience "a new birth of freedom" so that "government of the people, by the people, for the people, shall not perish from the earth." Upon completing this unit, students should have a better understanding of why Lincoln revered the union of the American states as "the last best, hope of earth." If your students lack experience in dealing with primary sources, you might use one or more preliminary exercises to help them develop these skills. The Learning Page at the American Memory Project of the Library of Congress includes a set of such activities. Another useful resource is the Digital Classroom of the National Archives, which features a set of Document Analysis Worksheets. Each lesson in this unit is designed to stand alone; taken together they present a robust portrait of how Lincoln viewed the American union. If there is not sufficient time to use all four lessons in the unit, either the first or third lesson convey Lincoln's understanding of the American union as a means to securing "Liberty to all"—with the first lesson focusing on the principled connection between the Declaration of Independence and the U.S. Constitution, and the third lesson addressing the practical connection between the Union war effort, the freedom of the newly emancipated slaves, and the preservation of American self-government. Adding the second lesson would show why Lincoln's understanding of the union and Constitution obliged the president to defend the nation from secession. Adding the fourth lesson would explore how Lincoln thought that only a common memory of the war as the chastening of God to both sides for the national (not Southern) sin of slavery could restore national unity.

http://edsitement.neh.gov/category/subject-areas/history-and-social-studies/us/civil-war-and-reconstruction-1850-1877?page=22

Where does our food come from? How is the climate changing what is on our plates? And, with a global population set to hit nine billion by 2050, how will we make sure everyone has enough to eat? These are just some of the questions pupils can explore through Food for Thought, Oxfam's new global citizenship resource for schools, now available on the Guardian Teacher Network. Can You Beat the System is a role-play task that encourages primary and secondary pupils to put themselves in the shoes of small-scale farmers working in less economically developed countries. Teams must work together to "produce" crops in the face of challenges from the weather, governments and traders. The aim is to explore factors affecting people's ability to grow food while highlighting inequalities in the global food system. Similar themes are covered in the primary lesson Farming Snakes and Ladders and the secondary lesson Farming Heroes . Both activities encourage pupils to identify the challenges faced by small-scale farmers and to suggest ways in which these can be overcome. The benefits of small-scale farming to the local community and wider world are also considered. Diet and Climate Change is an activity for primary and secondary pupils that explores the impact of climate change on the diet and income of an Ethiopian family who make their living by farming. Pupils consider the ways families can adapt to problems such as drought, and who should be helping them to do this. Grow Island is a role-play activity that examines issues of fairness and sustainability related to land ownership. Primary pupils are encouraged to think about the importance of land and why demand for it can be very high. For secondary pupils, Geography Mystery in Tanzania is a group work activity that looks at land purchases specifically for bio-fuels. Pupils are asked to consider who might benefit from this sort of investment and what some of the longer-term problems might be. The Power-Shift is a differentiated activity for primary and secondary pupils that aims to boost understanding of who's who in the global food system. Pupils consider the different groups in society who are able to make things fairer and the relative power each one has to bring about change. The activity is set in the context of something that pupils would like to change about their own school. Throughout the Food for Thought activities, pupils are encouraged to learn, think and take action as active global citizens. The action planning guide supports pupils in choosing a course of action, planning how to implement it and evaluating its success. The resources are supported by a Food for Thought wall chart that can be used to track pupils' learning and includes a teachers' guide to all of the activities. The Guardian Teacher Network has almost 100,000 pages of lesson plans and interactive materials. To see and share for yourself go to teachers.theguardian.com. There are also hundreds of jobs on the site, contact us for a free trial of your first advert: schoolsjobs.theguardian.com.

http://www.theguardian.com/education/2012/jan/16/oxfam-food-for-thought-resources

El país nórdico lidera el informe PISA con una enseñanza gratuita que pone en Primaria a los profesores más preparados Los niños finlandeses de hoy estarán el día de mañana entre los mejores profesionales ... It often occurs that we get confused with the concept of prior knowledge and its relationship to construction of new learning. It would only seem logical to always find out what the students know before delivering a class or a course of any discipline. However, the difference resides in what information we would be looking for and the purpose of retrieving that data. When we work with CLIL projects it is necessary to spend enough time in the process of exploring previous knowledge through different tools. The links that students can make to their personal experience and lives, the hypotheses and ideas they may have incidentally acquired about a certain topic will all contribute to set their always curious minds to work. Lev Vygotsky said " Learning always proceeds from the known to the new. Good teaching will recognize and build on this connection." Some tips to explore prior knowledge: -Use various tools individually or in groups such as: incomplete phrases or sentences, brainstorming, short multiple choice questionnaires, graphic organizers, cartoons, short videos, pictures, parts of stories and others. -Accept all the opinions without judging or correcting, stating that you are in an exploratory stage and that all ideas will be welcomed. -Keep a record of students' ideas to use at a later stage. -Refrain from correcting or indicating the right response. -Use your observations and collected information to decide on the project's future path. Good CLIL lessons should initiate by favoring risk taking to express ideas through drawings, writings and brainstorming allowing for different views and tolerating wrong or hilarious answers avoiding any judgment. It will be throughout the process of experiencing the unit/project that the students together with appropriate teacher's interventions and class discussions will be able to reflect on their own ideas. Teacher's tolerance, observation and confidence in students' possibilities are of crucial importance to set the atmosphere of high challenge and high support classrooms. Many CLIL projects or units would fit into a constructivist perspective if they were seriously "meaning oriented". One of the most common errors of some publications that present themselves under the "CLIL" umbrella is that they don't offer real problems or questions to be solved by the students. In those cases, information is just correlated around a certain "topic". Arriving to integration through a good leading question is one of the first important steps to make when planning a CLIL didactic unit or project. Jerome Bruner said: "The art of asking provoking questions is at least as important as that of providing clear answers [...], and the art of setting those questions to good use and keeping them alive is as important as the first two." Here are some tips to come up with a good question: -Avoid simple “yes-no” questions -The question will need reasoning and some research to be answered -It will relate to curricular guidelines and to students´ lives -It will motivate students to read, write, think and speak Can the world feed 10 billion people? Do revolutions always work? Do all animals have hearts? Why do animals travel? Why did humans lose their fur? Constructivism provides a strong rationale for content-based curricula such as CLIL, since it is holistically oriented and meaning seeking based. Then, LET'S START OUR CLIL PROJECTS WITH A GOOD QUESTION!!!! Sharing your scoops to your social media accounts is a must to distribute your curated content. Not only will it drive traffic and leads through your content, but it will help show your expertise with your followers. How to integrate my topics' content to my website? Integrating your curated content to your website or blog will allow you to increase your website visitors’ engagement, boost SEO and acquire new visitors. By redirecting your social media traffic to your website, Scoop.it will also help you generate more qualified traffic and leads from your curation work. Distributing your curated content through a newsletter is a great way to nurture and engage your email subscribers will developing your traffic and visibility. Creating engaging newsletters with your curated content is really easy.

http://www.scoop.it/t/teaching-english-from-a-constructivist-perspective/?tag=Education

a relative clause that modifies a noun or pronoun, as the clause that I told you about in This is the book that I told you about and who saw us in It was she who saw us. An adjective clause = a relative pronoun or relative adverb + subject + verb OR a relative pronoun or relative adverb + verb. Adjectives modify nouns and pronouns, giving a description or more information. An adjective clause is simply a group of words with a subject and a verb that ... An adjective clause is a dependent clause that functions as an adjective in the sentence. Adjective clauses can also be called relative clauses. An adjective ... On these occasions we use subordination to indicate that one part of a sentence is secondary (or subordinate) to another part. For example, to emphasize that father sets his unicorn traps at night, we can turn the first main clause into an adjective clause My father, ... More » (Dependent clauses are also called subordinate clauses.) There are three basic types of dependent clauses: adjective clauses, adverb clauses, and noun ... Adjective clauses, like adverb clauses, are introduced by dependent signals. ... ( Again, the adjective clause is underlined and modifies the subject "students."). An adjective clause usually comes after the noun it modifies and is made up of ... See the definition of Adjective Clause in Grammar Monster's list of grammar ... An adjective clause is a group of words with a subject and verb that modifies a noun in a sentence. In this lesson, we will learn how adjective... In addition to subject-pattern adjective clauses, there are also object-pattern ones . They have that name because in them, the relative pronoun replaces the ...

http://www.ask.com/web?qsrc=6&o=102140&oo=102140&l=dir&gc=1&qo=popularsearches&ad=dirN&q=Adjective+Clause

Harmonic series (music) Pitched musical instruments are often based on an approximate harmonic oscillator such as a string or a column of air, which oscillates at numerous frequencies simultaneously. At these resonant frequencies, waves travel in both directions along the string or air column, reinforcing and canceling each other to form standing waves. Interaction with the surrounding air causes audible sound waves, which travel away from the instrument. Because of the typical spacing of the resonances, these frequencies are mostly limited to integer multiples, or harmonics, of the lowest frequency, and such multiples form the harmonic series (see harmonic series (mathematics)). The musical pitch of a note is usually perceived as the lowest partial present (the fundamental frequency), which may be the one created by vibration over the full length of the string or air column, or a higher harmonic chosen by the player. The musical timbre of a steady tone from such an instrument is determined by the relative strengths of each harmonic. Partial, harmonic, fundamental, inharmonicity, and overtone A "complex tone" (the sound of a note with a timbre particular to the instrument playing the note) "can be described as a combination of many simple periodic waves (i.e., sine waves) or partials, each with its own frequency of vibration, amplitude, and phase." (See also, Fourier analysis.) A partial is any of the sine waves (or "simple tones", as Ellis calls them when translating Helmholtz) of which a complex tone is composed. A harmonic is any member of the harmonic series, an ideal set of frequencies that are positive integer multiples of a common fundamental frequency. The fundamental is also considered a harmonic because it is 1 times itself. A harmonic partial is any real partial component of a complex tone that matches (or nearly matches) an ideal harmonic. An inharmonic partial is any partial that does not match an ideal harmonic. Inharmonicity is a measure of the deviation of a partial from the closest ideal harmonic, typically measured in cents for each partial. Many pitched acoustic instruments are designed to have partials that are close to being whole-number ratios with very low inharmonicity; therefore, in music theory, and in instrument design, it is convenient, although not strictly accurate, to speak of the partials in those instruments' sounds as "harmonics", even though they have some inharmonicity. Other pitched instruments, especially certain percussion instruments, such as marimba, vibraphone, tubular bells, and timpani, contain mostly inharmonic partials, yet may give the ear a good sense of pitch because of a few strong partials that resemble harmonics. Unpitched, or indefinite-pitched instruments, such as cymbals, gongs, or tam-tams make sounds (produce spectra) that are rich in inharmonic partials (make "noise") and give no impression of implying any particular pitch. An overtone is any partial except the lowest partial. The term overtone does not imply harmonicity or inharmonicity and has no other special meaning other than to exclude the fundamental. It is the relative strengths of the different overtones that gives an instrument its particular timbre, tone color, or character. When writing or speaking of overtones and partials numerically, care must be taken to designate each correctly to avoid any confusion of one for the other, so the second overtone may not be the third partial, because it is second sound in series. Some electronic instruments, such as theremins and synthesizers, can play a pure frequency with no overtones (a sine wave). Synthesizers can also combine pure frequencies into more complex tones, such as to simulate other instruments. Certain flutes and ocarinas are very nearly without overtones. Frequencies, wavelengths, and musical intervals in example systems The simplest case to visualise is a vibrating string, as in the illustration; the string has fixed points at each end, and each harmonic mode divides it into 1, 2, 3, 4, etc., equal-sized sections resonating at increasingly higher frequencies. Similar arguments apply to vibrating air columns in wind instruments, although these are complicated by having the possibility of anti-nodes (that is, the air column is closed at one end and open at the other), conical as opposed to cylindrical bores, or end-openings that run the gamut from no flare (bell), cone flare (bell), or exponentially shaped flares (bells). In most pitched musical instruments, the fundamental (first harmonic) is accompanied by other, higher-frequency harmonics. Thus shorter-wavelength, higher-frequency waves occur with varying prominence and give each instrument its characteristic tone quality. The fact that a string is fixed at each end means that the longest allowed wavelength on the string (which gives the fundamental frequency) is twice the length of the string (one round trip, with a half cycle fitting between the nodes at the two ends). Other allowed wavelengths are 1/2, 1/3, 1/4, 1/5, 1/6, etc. times that of the fundamental. Theoretically, these shorter wavelengths correspond to vibrations at frequencies that are 2, 3, 4, 5, 6, etc., times the fundamental frequency. Physical characteristics of the vibrating medium and/or the resonator it vibrates against often alter these frequencies. (See inharmonicity and stretched tuning for alterations specific to wire-stringed instruments and certain electric pianos.) However, those alterations are small, and except for precise, highly specialized tuning, it is reasonable to think of the frequencies of the harmonic series as integer multiples of the fundamental frequency. The harmonic series is an arithmetic series (1×f, 2×f, 3×f, 4×f, 5×f, ...). In terms of frequency (measured in cycles per second, or hertz (Hz) where f is the fundamental frequency), the difference between consecutive harmonics is therefore constant and equal to the fundamental. But because human ears respond to sound nonlinearly, higher harmonics are perceived as "closer together" than lower ones. On the other hand, the octave series is a geometric progression (2×f, 4×f, 8×f, 16×f, ...), and people hear these distances as "the same" in the sense of musical interval. In terms of what one hears, each octave in the harmonic series is divided into increasingly "smaller" and more numerous intervals. The second harmonic, whose frequency is twice of the fundamental, sounds an octave higher; the third harmonic, three times the frequency of the fundamental, sounds a perfect fifth above the second harmonic. The fourth harmonic vibrates at four times the frequency of the fundamental and sounds a perfect fourth above the third harmonic (two octaves above the fundamental). Double the harmonic number means double the frequency (which sounds an octave higher). Harmonics and tuning If the harmonics are transposed into the span of one octave, some of them are approximated by the notes of what the West has adopted as the chromatic scale based on the fundamental tone. The Western chromatic scale has been modified into twelve equal semitones, which is slightly out of tune with many of the harmonics, especially the 7th, 11th, and 13th harmonics. In the late 1930s, composer Paul Hindemith ranked musical intervals according to their relative dissonance based on these and similar harmonic relationships. Below is a comparison between the first 31 harmonics and the intervals of 12-tone equal temperament (12TET), transposed into the span of one octave. Tinted fields highlight differences greater than 5 cents (1/20th of a semitone), which is the human ear's "just noticeable difference" for notes played one after the other (smaller differences are noticeable with notes played simultaneously). |Harmonic||12tET Interval||Note||Variance cents| |17||minor second||C♯, D♭||+5| |19||minor third||D♯, E♭||−2| |25||minor sixth||G♯, A♭||−27| |7||14||28||minor seventh||A♯, B♭||−31| The frequencies of the harmonic series, being integer multiples of the fundamental frequency, are naturally related to each other by whole-numbered ratios and small whole-numbered ratios are likely the basis of the consonance of musical intervals (see just intonation). This objective structure is augmented by psychoacoustic phenomena. For example, a perfect fifth, say 200 and 300 Hz (cycles per second), causes a listener to perceive a combination tone of 100 Hz (the difference between 300 Hz and 200 Hz); that is, an octave below the lower (actual sounding) note. This 100 Hz first-order combination tone then interacts with both notes of the interval to produce second-order combination tones of 200 (300 – 100) and 100 (200 – 100) Hz and all further nth-order combination tones are all the same, being formed from various subtraction of 100, 200, and 300. When one contrasts this with a dissonant interval such as a tritone (not tempered) with a frequency ratio of 7:5 we get, for example, 700 – 500 = 200 (1st order combination tone) and 500 – 200 = 300 (2nd order). The rest of the combination tones are octaves of 100 Hz so the 7:5 interval actually contains 4 notes: 100 Hz (and its octaves), 300 Hz, 500 Hz and 700 Hz. Note that the lowest combination tone (100 Hz) is a 17th (2 octaves and a major third) below the lower (actual sounding) note of the tritone. All the intervals succumb to similar analysis as has been demonstrated by Paul Hindemith in his book The Craft of Musical Composition, although he rejected the use of harmonics from the 7th and beyond. Timbre of musical instruments |This section needs additional citations for verification. (November 2011)| The relative amplitudes (strengths) of the various harmonics primarily determine the timbre of different instruments and sounds, though onset transients, formants, noises, and inharmonicities also play a role. For example, the clarinet and saxophone have similar mouthpieces and reeds, and both produce sound through resonance of air inside a chamber whose mouthpiece end is considered closed. Because the clarinet's resonator is cylindrical, the even-numbered harmonics are less present. The saxophone's resonator is conical, which allows the even-numbered harmonics to sound more strongly and thus produces a more complex tone. The inharmonic ringing of the instrument's metal resonator is even more prominent in the sounds of brass instruments. Human ears tend to group phase-coherent, harmonically-related frequency components into a single sensation. Rather than perceiving the individual partials–harmonic and inharmonic, of a musical tone, humans perceive them together as a tone color or timbre, and the overall pitch is heard as the fundamental of the harmonic series being experienced. If a sound is heard that is made up of even just a few simultaneous sine tones, and if the intervals among those tones form part of a harmonic series, the brain tends to group this input into a sensation of the pitch of the fundamental of that series, even if the fundamental is not present. Variations in the frequency of harmonics can also affect the perceived fundamental pitch. These variations, most clearly documented in the piano and other stringed instruments but also apparent in brass instruments, are caused by a combination of metal stiffness and the interaction of the vibrating air or string with the resonating body of the instrument. David Cope (1997) suggests the concept of interval strength, in which an interval's strength, consonance, or stability (see consonance and dissonance) is determined by its approximation to a lower and stronger, or higher and weaker, position in the harmonic series. See also: Lipps–Meyer law. Thus, an equal-tempered perfect fifth ( play (help·info)) is stronger than an equal-tempered minor third ( play (help·info)), since they approximate a just perfect fifth ( play (help·info)) and just minor third ( play (help·info)), respectively. The just minor third appears between harmonics 5 and 6 while the just fifth appears lower, between harmonics 2 and 3. |Wikimedia Commons has media related to Harmonic series.| - Fourier series - Klang (music) - Otonality and Utonality - Piano acoustics - Scale of harmonics - Stretched tuning - Undertone series - IEV 1994, sound: http://www.electropedia.org/iev/iev.nsf/display?openform&ievref=801-21-01 - Ibid, fundamental: http://www.electropedia.org/iev/iev.nsf/display?openform&ievref=801-30-01 - William Forde Thompson (2008). Music, Thought, and Feeling: Understanding the Psychology of Music. p. 46. ISBN 978-0-19-537707-1. - John R. Pierce (2001). "Consonance and Scales". In Perry R. Cook. Music, Cognition, and Computerized Sound. MIT Press. ISBN 978-0-262-53190-0. - Martha Goodway and Jay Scott Odell (1987). The Historical Harpsichord Volume Two: The Metallurgy of 17th- and 18th- Century Music Wire. Pendragon Press. ISBN 978-0-918728-54-8. - Riemann by Shedlock (1876). Dictionary of Music. Augener & Co., London. p. 143. let it be understood, the second overtone is not the third tone of the series, but the second. - Juan G. Roederer (1995). The Physics and Psychophysics of Music. p. 106. ISBN 0-387-94366-8. - Fonville, John. 1991. "Ben Johnston's Extended Just Intonation: A Guide for Interpreters", p.121. Perspectives of New Music 29, no. 2 (Summer): 106–37. - Hindemith, Paul (1942). The Craft of Musical Composition: Book 1—Theoretical Part,[page needed]. Translated by Arthur Mendel (London: Schott & Co; New York: Associated Music Publishers. ISBN 0901938300). . - Cope, David (1997). Techniques of the Contemporary Composer, p. 40–41. New York, New York: Schirmer Books. ISBN 0-02-864737-8. - Interaction of reflected waves on a string is illustrated in a simplified animation - A Web-based Multimedia Approach to the Harmonic Series - Importance of prime harmonics in music theory - "Addendum to 'The Devolution of the Shepherd Trumpet and It's Seminal Importance in Music History'" (link at bottom of page). Describes how the harmonic series is the basis of European folk song melodies - Octave Frequency Sweep, Consonance & Dissonance - The combined oscillation of a string with several of its lowest harmonics can be seen clearly in an interactive animation from Edward Zobel's "Zona Land" (requires plugin).

https://en.wikipedia.org/wiki/Overtone_series

This week in history: The Fall of Constantinople had profound consequences On May 29, 1453 — 560 years ago this week — Constantinople fell to the Ottoman Turks. The fall of this great city signaled the end of the Byzantine Empire, the medieval incarnation of the Roman Empire, and saw the armies of Islam spread into Europe from Asia for the first time. In A.D. 330, the Roman Emperor Constantine founded the city of Constantinople on the Greek village of Byzantine to be the new imperial capital. Sitting on the Bosporus strait, which connects Europe and Asia, the new city was more easily defended than Rome, and it was a Christian city to reflect the emperor's religious preference. Like Rome, Constantinople had seven hills divided into 14 districts. For centuries, the city stood as the center of imperial power, even after the fall of the Western Roman Empire in A.D. 476. Historians refer to this medieval incarnation of the empire as Byzantine. The Franks and the Italians of the time referred to its inhabitants simply as “the Greeks.” The inhabitants themselves, however, continued to refer to themselves as Romans, and saw their emperors as the literal successors to Augustus, Marcus Aurelius and Constantine. Containing impressive city walls, Constantinople was virtually impervious to attack, such as when an army of Goths approached the city after the battle of Adrianople in A.D. 378. After the rise of Islam, the Byzantine empire lost much of its territory in the Middle East and North Africa, but the city of Constantinople proved an impervious rock upon which wave after wave of Muslim armies couldn't break. As Constantinople held the line against Islam in the East, modern Western civilization developed in France and Western Europe. Though the Franks had defeated Islamic armies from Spain, the loss of Byzantine to Islam may well have seen the creation of a Muslim Europe. Toward the end of the Middle Ages, however, Byzantine power was waning considerably. Practicing Orthodox Christianity, Constantinople had fallen to Catholic knights during the Fourth Crusade in 1204, ushering in nearly 60 years of Catholic rule before an Orthodox emperor was able to retake the throne. The mid-14th century saw the Black Death claim the lives of perhaps half the city's population. By the early 15th century, the Islamic Ottoman Turks had conquered virtually all of present day Turkey, and the Byzantine empire was a shadow of its former self, consisting of a few scattered territories and islands outside of Constantinople itself. In 1451, Mehmed II succeeded his father to become the Ottoman sultan. In his book “1453: The Holy War for Constantinople and the Clash of Islam and the West,” historian Roger Crowley described the 19-year-old ruler: “The man whom the Renaissance later presented as a monster of cruelty and perversion was a mass of contradictions. He was astute, brave and highly impulsive — capable of deep deception, tyrannical cruelty and acts of sudden kindness. He was moody and unpredictable, a bisexual who shunned close relationships, never forgave an insult, but who came to be loved for his pious foundations.” Upon becoming sultan, Mehmed immediately began a new building program for his navy, and soon set about plans to do something that the many sultans before him couldn't — the conquest of Constantinople. In early 1453, he took an army of somewhere between 100,000 and 200,000 Ottoman troops into Byzantine territory, and on April 6 began major siege operations against the city. Constantine XI proved to be the last of the Byzantine emperors. Having ruled since 1449, Constantine knew the empire's defenses alone, including more than 12 miles of walls, were not enough to repel a determined Ottoman siege or assault. - Over 100 missing, 14 dead as strong quake... - 5 things to know about the Super Bowl and sex... - Rubio faces barrage of attacks in Republican... - Search for possible survivors of midair crash... - US adds just 151k jobs in January; 4.9... - Police think 6 Chicago deaths result of... - Pope's sex abuse panel tells survivor to take... - Lindsey Vonn wins women's World Cup downhill - US adds just 151k jobs in January; 4.9... 28 - Cruz tops Trump in Iowa; Clinton,... 22 - Obama pushes for every child to learn... 21 - Cruz, Rubio eye NH momentum as Dems... 15 - Donald Trump calls for Iowa election... 14 - Americans are packing classes in how to... 12 - Kentucky Sen. Rand Paul quits 2016 GOP... 12 - Once a bromance, now a brawl: Trump and... 10

http://www.deseretnews.com/article/865580842/This-week-in-history-The-Fall-of-Constantinople-had-profound-consequences.html

Assessment for Learning This Page was created by Bryan Funk (2009) This Page was edited by Kari Duffy (January 2010) - 1 Assessment for Learning - 2 Six Big Strategies that Matter - 3 Assessment for Learning vs. Assessment of Learning - 4 Role of the Student - 5 Role of the Teacher - 6 Feedback for the Student - 7 Criterion-Referenced Assessment - 8 Setting and Using Criteria - 9 Planning with Assessment in Mind - 10 Stop Motion - 11 References Assessment for Learning Assessment for learning focuses on engaging students in classroom assessment in support of their own learning and informing teachers about what to do next to help students to progress. Assessment for learning is assessment for improvement not assessment for accountability as can be the case with summative assessments (Stiggins, 2002). The keys to Assessment for Learning (AFL) is to use a variety of assessment tools and methods in order to provide ongoing evidence to students, teachers and parents that demonstrates how well each student is mastering the identified outcomes. This evidence is used to provide descriptive feedback to the students and to enable the teacher to differentiate the instruction to meet the needs of individual students or groups. Black and William clearly indicate that formative assessment (AFL) will raise performance standards and improve overall student success (Black & William, 1998). In his book, Talk About Assessment: Strategies and Tools to Improve Learning, Damian Cooper (2007) defines Assessment for Learning as "assessment designed primarily to promote learning. Early drafts, first tries, and practice assignments are all examples of assessment for learning", and describes Assessment of Learning as "assessment designed primarily to determine student achievement at a given point in time. Report card grades should be comprised of data from assessments of learning". Cooper's (2007) first chapter introduces the educator to Eight Big Ideas about assessment: Big Idea 1 Assessment serves different purposes at different times: it may be used to find out what students already know and can do; it may be used to help students improve their learning; or it may be used to let students and their parents know how much they have learned in a prescribed period of time. Big Idea 2 Assessment must be planned and purposeful. Big Idea 3 Assessments must be balanced, including oral, performance, and written tasks, and be flexible in order to improve learning for all students. Big Idea 4 Assessment and instruction are inseparable because effective assessment informs learning. Big Idea 5 For assessment to be helpful to students, it must inform them in words, not numerical scores or letter grades, what they have done well, what they have done poorly, and what they need to do next in order to improve. Big Idea 6 Assessment is a collaborative process that is most effective when it involves self-, peer, and teacher assessment. Big Idea 7 Performance standards are an essential component of effective assessment. Big Idea 8 Grading and reporting student achievement is a caring, sensitive process that requires teachers' professional judgement. Six Big Strategies that Matter Black and Wiliam's work led to the development of five performance strategies for assessment for learning; these five have since been re-structured by Dr. Linda Kaser and Dr. Judy Halbert into six big strategies that matter and are described below (Koehn, 2008 p. 2). 1. Providing learners with clarity about and understanding of the learning intentions of the work being done (learners are presented the learning intentions at the beginning of the lesson, throughout the lesson, and refer to the learning intentions in their reflections and responses so teachers can see that connections between tasks and what is supposed to be learned are made) 2. Providing to and co-developing with learners the criteria for success (what will the finished task look like, how will you share your understandings with others?) 3. Providing ongoing descriptive feedback that moves learning forward for each learner (using feed forward in language the students understand; how can the next task improve upon the previous?) 4. Designing and using thoughtful classroom questions to lead discussions that generate evidence of learning (allowing the students to participate and interact amongst each other in meaningful oral discussion – talk is student to student(s), not a dialogue between teacher and one student) 5. Putting learners to work as learning/teaching resources for each other using self and peer assessment (student coaching, students understanding learning intentions so well that they can teach a younger student or peer) 6. Doing everything we can think of to make sure that learners have ownership of their own learning (empowering each student to succeed). Assessment for Learning vs. Assessment of Learning Gregory, Cameron, and Davies (1997) outline some distinct differences between Assessment for Learning and Assessment of Learning. Educators are using these terms to help distinguish between the teacher's role as a learning coach versus the teacher's role of judging the extent of a student's achievement in relation to an established standard. This assessment is considered summative and is done at the end. 1. Assessment for learning is the big deal, while assessment of learning is the done deal. 2. Assessment for learning is formative, while assessment of learning is summative. 3. Assessment for learning is supportive, while assessment of learning measures. 4. Assessment for learning uses descriptions, while assessment of learning uses scores. 5. Assessment for learning happens day by day, moment by moment, while assessment of learning happens at the end. The assertion is that neither one is better than the other, but both need to be used within a students learning so that the student is able to understand not only the work that is being asked of them, but also how their own learning occurs. Assessment for learning is intended to be both diagnostic and formative to help students improve their learning. (chart from Anne Davies website: http://www.annedavies.com/assessment_for_learning_tr_tjb.html) Role of the Student Students are involved in identifying achievement expectations from the beginning of the learning by studying exemplars of strong and weak work. It is very important that students have a clear understanding of the learning intentions and expected outcomes of the work they are being asked to do. Assessment for learning will make the student’s learning visible and will enable both the teacher and learner to reflect and adjust the learning process. The learners play an important role in developing and understanding the scaffolding they will be climbing as they approach those outcomes. Students partner with their teacher to continuously monitor their current level of attainment in relation to agreed-upon expectations so they can set goals for what to learn next and thus play a role in managing their own progress. Students are asked to communicate evidence of learning to one another, to their teacher, and to their families, and they do so along the entire learning journey, not only at the end. the learning, students are inside the assessment process, watching themselves grow, feeling in control of their success, and believing that continued success is within reach if they keep trying. Role of the Teacher Assessment for learning not only provides reflective feedback to guide the learning process, but empowers students to control and dictate the direction of their learning. Purposeful use of AFL will enable students to experience metacognition where they engage and reflect on their learning experience. “Intelligent thought involves 'metacognition' or self monitoring of learning and thinking" (Shepard, 2000. p. 8). Although much of assessment for learning is about empowering the student to understand and take control of their learning, the teacher plays a critical role in chosing the appropriate assessments and using them to differentiate instruction to meet the individual needs of the students. The teacher is responsible for aligning instruction with the targeted outcomes, selecting and adapting materials and resources and identifying specific learning requirements of students or groups or students. Once the teacher has collected begun collecting information through assessment their focus becomes creating differentiated teaching strategies and learning opportunities in order to assist individual students move forward in their learning. This is partially accomplished by providing immediate feedback and direction to students and then completing the cycle again. In Rethinking Classroom Assessment: Assessment for Learning, Assessment as Learning, Assessment of Learning there are four critical questions that the teacher must ask when planning for assessment for learning: Why am I assessing? If the intent of assessment is to enhance student learning teachers use assessment for learning to uncover what students believe to be true and to learn more about the connections students are making, their prior knowledge, preconceptions, gaps, and learning styles. This information is used to inform and differntiate instruction to build on what students already know and to challenge students when their are problems inhibiting progression to the next stages of learning. Teachers use this information to provide their students with descriptive feedback that will further their learning and not as a sumamtive assessment or to report a grade. What am I assessing? Assessment for learning requires ongoing assessment of the outcomes that comprise the intended learning. In most cases these are the curriculum outcomes. Teachers create assessments that will expose students’ thinking and skills in relation to the intended learning, and the common preconceptions. What assessment method should I use? When planning assessment for learning, the teacher must think about what assessment is designed to expose, and must decide which assessment approaches are most likely to give detailed information about what each student is thinking and learning. The methods need to incorporate a variety of ways for students to demonstrate their learning. For example, having students complete tasks orally or through visual representation allow those who are struggling with reading or writing to demonstrate their learning. How can I use the information? The information collected in assessment for learning is used to report to the student and by offering descriptive, on time feedback and to provide the teacher with information to allow for changes in instruction for individual students or groups of students. Feedback for the Student Black and Wiliam (1998) suggested that feedback was a key component in assessment for learning. Cooper (2007) and Davies (2000a) assert that the quality of the feedback matters, as well as the timing. Quality feedback is descriptive feedback. Descriptive feedback makes it clear for the learning what is working and what needs to be worked on. Allowing students to adjust or change what they are doing through descriptive feedback, students are more likely to be successful (Davies, 2000a). In her book Making Classroom Assessment Work, Anne Davies (2000b) tells us that descriptive feedback that supports learning, and: - comes during, as well as after, the learning - is easily understandable and related directly to the learning - is specific - so performance can improve - involves choice on the part of the learner as to what and how to receive feedback - is part of an ongoing conversation about learning - is in comparison to models, exemplars or descriptions - is about the performance or the work, not the person Feedback for learning is an integral part of the teaching process. It is the vital link between the teacher’s assessment of a student’s learning and the action following that assessment. Immediate feedback is key maximizing student learning. Cooper (2007) and Gregory et al (1997) provide many examples of feedback that is easy and fast. In their book Setting and Using Criteria, Gregory et al (1997) provide ten quick ways that a teacher can give immediate feedback to guide the student's learning, without putting a mark in the gradebook. One such example is Met/Not Yet/I Noticed. This technique gives the student immediate feedback when the criteria is set up in a ruberic and the teacher is simply checking off Met or Not Yet and giving descriptive feedback in the I Noticed column. Descriptive feedback makes explicit connections between students’ thinking and the learning that is expected. It addresses misinterpretations and lack of understanding. Feedback should help identify the next steps and an example of what good work looks like (Davies, 2000b). Feedback for learning will support or challenge an idea that a student holds. It allows the teacher to provide recognition for achievement and growth, and to give precise directions for improvement. Good descriptive feedback should also cause students to think about, and allow them to respond to, the teacher's or peer's suggestions. With a criterion-referenced standard the student's performance is measured against a predetermined set of performance indicators. We commonly see this type of assessment outside of the school setting. For example, when a coach is teaching a new skill to an athlete and a driver's road test. Performance standards and ruberics are becoming more and more common is the educational setting, as teachers see the merit in allowing the students to know what the criteria is before they begin the task (Cooper, 2007). Another technique that is often employed is including students in setting the criteria. This increases student buy-in and makes them accountable to the standards they have set themselves. Assessment for learning provides information about what students can already do and what they are not able to do (Gregory, et al 1997). It is measured using a predetermined set of exemplars. This must be shared with the students before they begin the work. After the students have done their work, their tasks can be measured using the criteria that was set our for that task. This provides the necessary information for the teacher to create the next steps in instruction. Students will be in various places in their learning and by using assessment for learning the teacher can compare the student progress with the intended objectives and adjust the pace of instruction, the resources, or the amount of work required in order to lesson confusion and frustration on the part of the student (Davis, 2000b). By focusing on what the student does know and moving forward the learner is being supported rather than criticized. By knowing the criteria ahead of time, frustration is decreased and the student has a sense of ownership over their own work. Using performance standards and exemplars also decreases the frustration level of the students because they are able to understand how the criteria is being presented (Gregory, et al, 1997). Setting and Using Criteria In Setting And Using Criteria, Gregory et al (1997) outline some effective approaches to establishing criteria in the classroom. Teachers can set criteria on their own, or include students in deciding what will be valued for that particular task. Criteria should always come out of the learning outcomes set for that grade. Criteria should be set for projects and assignments and does not have to be set for day to day tasks. By setting the criteria, the teacher is outlining how the task will be valued or judged. Gregory et al (1997) have found that the "following four step process for creating criteria with students encourages student participation, understanding, and ownership". Step One: Brainstorm Step Two: Sort and catagorize Step Three: Make and post a T-chart Step Four: Add, revise, refine This process can be done with the student's using something very familiar to start with, so that they learn how to set criteria. After some amount of practice, the students become faster at pulling out the major ideas or skills that you want them to know or do. Planning with Assessment in Mind Grant Wiggins and Jay McTighe (1998) coined the idea of backward design planning with their book, "Understanding by Design". The focus of their work was to state that if learning was to be effective for the students, the teacher must begin with the final destination in mind, and that the programs or activities must be 'backward in design' (Wiggins & McTighe, 1998). Designing curriculum this way has been described as backward because teachers traditionally start curriculum planning with interesting activities and textbooks in mind, rather than starting with the big ideas or goals they want the students to master (Wiggins & McTighe, 1998). Teachers should be clear about what learning targets or goals will be set for the students and what formative and culminating assessments will be used to provide evidence that the students have mastered those targets or goals. The students need to be informed what the assessments will be along the way and for the final culmination, so that they have a clear sense of what their goals need to be. Students should also be given the reasons why each assessment will be looked at, so that they will understand what is being asked of them and when (Wiggins & McTighe, 1998). Teachers begin with the end in mind, and set the task to reflect the learning. Teachers should inform the students about the big ideas and essential questions, the performance requirements, and the evaluative criteria at the beginning of the unit or course. The students should be able to describe these goals (big ideas and essential questions) of the unit or course. This helps to ensure that the students are aware of the expectations and optimal learning takes place. "To begin with the end in mind means to start with a clear understanding of your destination. It means to know where you’re going so that you better understand where you are now so that the steps you take are always in the right direction." (Covey, 1990) Created by Kirsten O'Coin (2016) Black, P., & Wiliam, D. (1998). Inside the black box: Raising standards through classroom assessment [Electronic version]. Phi Delta Kappan, 80(2). 139-44. 32 Cooper, D. (2010). Talk About Assessment: High School Strategies and Tools. Nelson/Thomson, Canada Ltd., Toronto. Cooper, D. (2007). Talk about assessment: Strategies and tools to improve learning. Nelson/Thomson, Canada Ltd., Toronto. Covey, S. (1990). The Seven Habits of Highly Effective People. New York: Fireside. Davies, A. (2000a). Feed Back...Feed Forward: Using assessment to boost literacy learning. Online Journal © 2003 Classroom Connections International - www.connect2learning.com (Originally published in Primary Leadership. Vol.2 No. 3 Spring Issue (2000) p.53-55). Davies, A. (2000b). Making Classroom Assessment Work. Courtney, BC: Classroom Connections. Gregory, K., Cameron, C. & Davies, A. (2000). Self-Assessment and Goal-Setting. Merville, BC: Connections Publishing. Gregory, K., Cameron, C. & Davies, A. (1997). Setting and Using Criteria. Merville, BC: Connections Publishing. Koehn. (2008). Together is better (BCTF Teacher Inquirer). Retrieved February 18, 2009 from Website: http://bctf.ca/uploadedFiles/Publications/TeacherInquirer/archive/2008-09/2008-10/Koehn.pdf 53 Rethinking Classroom Assessment: Assessment for Learning, Assessment as Learning, Assessment of Learning - new publication of the Western and Northern Canadian Protocol (WNCP). Retrieved March 1, 2009 from: http://www.aac.ab.ca/public/rethinking.pdf Shepard, L. A. (2000). The role of assessment in a learning culture [Electronic version]. Educational Researcher, 29(7). 4-14. 32 Stiggins, R. (2002). "Assessment Crisis: The Absence of Assessment FOR Learning.” Phi Delta Kappan, 83 (10), 758–765 Stiggins, R. (2005). "From Formative Assessment to Assessment FOR Learning: A Path to Success in Standards-Based Schools." Phi Delta Kappan, 87,(4) Wiggins, G. (1993). Assessing Student Performance: Exploring the purpose and limits of testing. San Fransisco, CA: Jossey-Bass. Wiggins, G. & McTighe, J. (1998) Understanding by Design. Alexandria, VA:ASCD.

http://etec.ctlt.ubc.ca/510wiki/Assessment_for_Learning

Earth and Mars are two of the rocky terrestrial planets that orbit within the inner Solar System. In some ways they are very similar, but in other ways, they couldn’t be more different. Let’s take a look at Earth and Mars and consider their similarities and differences. The origin of Earth and Mars is the same for all of the planets in the Solar System. Scientists believe that the Sun, the planets and everything in the Solar System all formed at the same time within the solar nebula; a giant cloud of cold molecular hydrogen. The Sun formed in the middle of the Solar System, and its rapid rotation created a flattened disk of material surrounding it. It was from within this disk that Earth and Mars and all the planets formed. Both Earth and Mars are terrestrial planets, made up of rock and metal. We know quite a bit about the internal structure of Earth, but can only guess at the structure of Mars. Scientists think that both planets have an inner core of metal surrounded by a mantle of rock. A thin crust covers the mantle. Since Mars is a smaller world, it’s believed that it cooled faster, so it’s crust is thicker. Mars lacks a magnetic field, so it probably doesn’t have a rotating liquid metal core like we have here on Earth. Unlike Earth, Mars has no plate tectonics. Perhaps it was because of this rapid cooling, but the crust on Mars is one thick shell surrounding the entire planet. While plate tectonics are constantly resurfacing planet Earth, we can see that the surface of Mars is ancient, hammered by thousands of meteorite impacts. The lack of plate tectonics also allowed hotspots in the mantle to remain in the same position for billions of years. This was how the largest volcanoes in the Solar System, such as Olympus Mons, could get so big. Mars is small. At only 6,792 km across, it’s about half the diameter of Earth, and has only 10% of the Earth’s mass. This means that it has a much lower surface gravity than Earth. If you could stand on the surface of Mars, you would only experience about 1/3rd the gravity you have on Earth. You would be able to jump 3 times higher than you could on Earth. Earth and Mars also share water. But here on Earth, water is everywhere; the oceans account for 71% of the surface area of the planet. Mars looks dry and dusty, but there are vast deposits of water ice at the planet’s northern and southern poles. It’s thought that Mars had large quantities of water on its surface billions of years ago; there’s evidence of flooding and ancient river valleys. But that water is long gone. What water remains is locked as ice underneath the ground. Of course, the biggest difference between Earth and Mars is life. Here on Earth, life is everywhere – you can see the green forests from space. Mars looks dry and dusty, and no spacecraft sent to Mars have found any life. Scientists think there could be life hiding underground with water, or inside rocks, but nothing has been found so far. NASA Solar System Exploration: Earth-Mars Comparison Chart NASA: Earth and Mars

http://www.universetoday.com/22677/earth-and-mars/

Black Hole Science To begin to understand Hawking's contributions, we must look at the tiny, subatomic particles that comprise everything in our universe. For instance, particle pairs are constantly appearing and disappearing together. In every pair, there's a particle and an antiparticle with the opposite properties, like a proton and its corresponding antiparticle, the antiproton. Without being interrupted, these particles simultaneously appear, cancel each other out and disappear as fast as they arrived. But Hawking wondered what would happen to these particles if they were face-to-face with a black hole. In the 1970s, Hawking put forth the idea that a black hole likely sucks in one particle -- usually the antiparticle -- while allowing the other particle to escape. According to Hawking's theory, it's this leftover particle at the entrance of a black hole that ends up being emitted as a type of radiation called Hawking radiation. He hypothesizes that since the other particle falls into the black hole, it helps reduce the mass of the black hole by that incremental amount. Over time, black holes decrease in mass and collapse, resulting in a huge explosion that spits out matter throughout space. In a sense, Hawking's work suggests there's a lot more going on in black holes and the areas around them. Since researchers have to measure the presence of black holes indirectly, it's been difficult to confirm or counter Hawking's theory. One group created a miniature black hole of sorts in a lab and observed that Hawking radiation could be real [source: Shiga]. Still, others think the theory can only be confirmed by evidence from a real black hole. Then followed the information paradox -- the debate surrounding what happened to qualities of matter once inside a black hole. Hawking wasn't afraid to be bold in drawing conclusions. In 1997, he made a bet with colleague John Preskill, arguing that information is permanently lost once it falls into a black hole [source: Hogan]. But in 2004, Hawking admitted that information isn't lost or channeled into another universe, but rather it seeps back into the existing universe in distorted form. Hawking publicly admitting defeat confirms his attitude toward science as a field that's constantly adding to and correcting itself. His work on black holes might not seem particularly groundbreaking at first glance, but it spurred conversation that might not have taken place otherwise What does Hawking think about the origins of the universe? Find out on the next page.

http://science.howstuffworks.com/dictionary/famous-scientists/physicists/stephen-hawking2.htm

Library of Congress In its broadest sense, impeachment is the process by which public officials may be removed from office on the basis of their conduct. Strictly speaking, it is the decision by a legislature to accuse an official of one or more offenses that warrant removal according to constitutional standards. A vote to impeach then triggers a trial based on those charges. The most famous impeachment proceedings have involved presidents, but every state has its own procedures. Most follow the federal model in general, but vary widely in their specifics. At the federal level, impeachment starts in the House of Representatives, where members may initiate resolutions to impeach a sitting president. The House Judiciary Committee decides if a resolution merits a formal impeachment inquiry. A simple majority vote in the full House can launch a formal inquiry. The House Judiciary Committee conducts an investigation to determine if allegations against a president warrant charges, or articles of impeachment. If a simple majority of the full House votes to charge a president with at least one article of impeachment, that indictment will move to the Senate for trial. At that point, the president has been "impeached" by the House. House members act as or appoint congressional prosecutors. The chief justice of the Supreme Court presides over the trial in the Senate chamber. A two-thirds vote is required to convict and remove from office. The U.S. Constitution states that, "The President, Vice President and all civil Officers of the United States, shall be removed from Office on Impeachment for, and Conviction of, Treason, Bribery, or other high Crimes and Misdemeanors." (Article II, Section ). The House of Representatives has impeached two Presidents: Andrew Johnson and Bill Clinton. Johnson was charged in 1868 with eight articles of impeachment, but was acquitted by a single vote in the Senate trial (depicted in the above engraving). Bill Clinton was charged with four articles of impeachment by the House in 1998, but was acquitted by the Senate early the next year. Richard Nixon resigned in 1974 before a final vote in the full House could send him to trial on three articles of impeachment. Each state constitution outlines a unique impeachment procedure, including variations on the list of impeachable offenses, protocol for an impeachment trial and the body responsible for an initial investigation. According to the Associated Press, seven governors in U.S. history have been removed from office following impeachment proceedings. The National Conference of State Legislatures said that a longer list would include states that have investigated governors for alleged offenses, voted to impeach a governor ahead of a trial, or held trials that resulted in acquittal. The only governor to be removed from office in the last 80 years was Gov. Evan Mecham of Arizona, who was convicted in 1988 of obstructing justice and misusing $80,000 in state money that he was charged with funneling to his car dealership to keep it afloat. In January 2008, the Illinois House of Representatives voted 114-1 to impeach Gov. Rod Blagojevich for abuse of power in connection with the federal investigation that had led to his arrest the month before. Mr. Blagojevich was charged with trying to sell the Senate seat vacated by Barack Obama and with seeking to extort campaign contributions in return for official actions, including providing reimbursement to a hospital. Following the process that has been generally adopted by state legislatures in recent decades, the Illinois House created a special investigative committee, which made a recommendation in favor of impeachment to the entire body. In all states except Alaska, Nebraska and Oregon, the House votes on articles of impeachment ahead of a trial. In Alaska, the process is reversed, according to The Book of States. That state's Senate must impeach a governor by voting on impeachment articles in order to initiate a trial in the House. Nebraska is the only state with a unicameral legislature. Without a state House, the Nebraska Senate votes to impeach before passing articles on to the state Supreme Court for a trial. Oregon is the only state without constitutional provisions for impeachment of a governor or other executive and judicial officers, according to the NCSL. Those officials may be removed from office, but not by the state's legislature. State courts in Oregon may try public officials for criminal offenses, but the procedure depends upon the jurisdiction of a crime. Nebraska convicted and dismissed from office its first governor in 1871, but has not fully impeached any governor since. Alaska has never removed a governor from office. — Rebecca Cathcart, Jan 12, 2008 ARTICLES ABOUT IMPEACHMENT Facing an impeachment effort, Gov. Paul LePage said that he would skip the traditional State of the State speech to a joint session of the Legislature. January 13, 2016, Wednesday Anger over the stuttering economy and political scandals is fueling calls for the ouster of Ms. Rousseff, who has three years to go in her presidential term. December 14, 2015, Monday A pesar de una crisis económica y varios escándalos de corrupción, el juicio político de Dilma Rousseff acapara la atención de los brasileros. December 4, 2015, Friday A faltering economy, an environmental disaster and street protests take a back seat as a powerful lawmaker, Eduardo Cunha — himself facing corruption charges — tries to oust the president. December 4, 2015, Friday The move against Ms. Rousseff by Eduardo Cunha, the House speaker, who is himself battling charges in a bribery scheme, opens a new phase of uncertainty in Brazil. December 3, 2015, Thursday Vice President Ahmed Adeeb will be charged with terrorism in what the government called a plot to kill the president, a minister said. November 6, 2015, Friday House Republicans may try to impeach John Koskinen, the commissioner of the Internal Revenue Service, though the specifics of any supposed impeachable offenses are vague. October 19, 2015, Monday A judicial panel’s suggestion that Judge Fuller, who resigned after an arrest, could still be impeached was seen as a broader message that the judiciary was embracing a harder line. September 19, 2015, Saturday Parliament in the Maldives voted overwhelmingly Tuesday to impeach Vice President Mohamed Jameel after accusing him of dereliction of duties. July 22, 2015, Wednesday The ousted prime minister voiced confidence in her innocence in connection with a plan to prop up rice prices, which her opponents called a waste of money. May 20, 2015, Wednesday Ahead of a summer break, House Republicans have moved forward with plans to sue President Obama. John A. Boehner, the House speaker, said he intended to file a lawsuit accusing President Obama of failing to carry out laws passed by Congress. For the past year, the 32-year-old mayor of Ellisville, Mo., has been at war with its City Council, which is now trying to remove him from office. Gov. Rod Blagojevich proclaimed his innocence before the Illinois Senate voted to impeach him. (Video: MSNBC)

http://topics.nytimes.com/top/reference/timestopics/subjects/i/impeachment/index.html?query=Ahmadinejad,%20Mahmoud&field=per&match=exact

Freedom of the press Freedom of the press or freedom of the media is the freedom of communication and expression through mediums including various electronic media and published materials. While such freedom mostly implies the absence of interference from an overreaching state, its preservation may be sought through constitutional or other legal protections. With respect to governmental information, any government may distinguish which materials are public or protected from disclosure to the public based on classification of information as sensitive, classified or secret and being otherwise protected from disclosure due to relevance of the information to protecting the national interest. Many governments are also subject to sunshine laws or freedom of information legislation that are used to define the ambit of national interest. The United Nations' 1948 Universal Declaration of Human Rights states: "Everyone has the right to freedom of opinion and expression; this right includes freedom to hold opinions without interference, and impart information and ideas through any media regardless of frontiers" This philosophy is usually accompanied by legislation ensuring various degrees of freedom of scientific research (known as scientific freedom), publishing, press and printing the depth to which these laws are entrenched in a country's legal system can go as far down as its constitution. The concept of freedom of speech is often covered by the same laws as freedom of the press, thereby giving equal treatment to spoken and published expression. - 1 Relationship to self-publishing - 2 Status of press freedom worldwide - 3 History - 4 Implications of new technologies - 5 Organizations for press freedom - 6 See also - 7 Notes - 8 References - 9 External links Relationship to self-publishing Freedom of the press is construed as an absence of interference by outside entities, such as a government or religious organization, rather than as a right for authors to have their works published by other people. This idea was famously summarized by the 20th century American journalist, A. J. Liebling, who wrote, "Freedom of the press is guaranteed only to those who own one". Freedom of the press gives the printer or publisher exclusive control over what the publisher chooses to publish, including the right to refuse to print anything for any reason. If the author cannot reach a voluntary agreement with a publisher to produce the author's work, then the author must turn to self-publishing. Status of press freedom worldwide Beyond legal definitions, several non-governmental organizations use other criteria to judge the level of press freedom around the world: - Reporters Without Borders considers the number of journalists murdered, expelled or harassed, and the existence of a state monopoly on TV and radio, as well as the existence of censorship and self-censorship in the media, and the overall independence of media as well as the difficulties that foreign reporters may face. - The Committee to Protect Journalists (CPJ) uses the tools of journalism to help journalists by tracking press freedom issues through independent research, fact-finding missions, and firsthand contacts in the field, including local working journalists in countries around the world. CPJ shares information on breaking cases with other press freedom organizations worldwide through the International Freedom of Expression Exchange, a global e-mail network. CPJ also tracks journalist deaths and detentions. CPJ staff applies strict criteria for each case; researchers independently investigate and verify the circumstances behind each death or imprisonment. - Freedom House likewise studies the more general political and economic environments of each nation in order to determine whether relationships of dependence exist that limit in practice the level of press freedom that might exist in theory. So the concept of independence of the press is one closely linked with the concept of press freedom. Worldwide press freedom index Every year, Reporters Without Borders establishes a ranking of countries in terms of their freedom of the press. The worldwide Press Freedom Index list is based on responses to surveys sent to journalists that are members of partner organisations of the RWB, as well as related specialists such as researchers, jurists and human rights activists. The survey asks questions about direct attacks on journalists and the media as well as other indirect sources of pressure against the free press, such as non-governmental groups. RWB is careful to note that the index only deals with press freedom, and does not measure the quality of journalism. In 2011–2012, the countries where press was the most free were Finland, Norway and Germany, followed by Estonia, Netherlands, Austria, Iceland, and Luxembourg. The country with the least degree of press freedom was Eritrea, followed by North Korea, Turkmenistan, Syria, Iran, and China. Freedom of the Press Freedom of the Press is a yearly report by US-based non-governmental organization Freedom House, measuring the level of freedom and editorial independence enjoyed by the press in every nation and significant disputed territories around the world. Levels of freedom are scored on a scale from 1 (most free) to 100 (least free). Depending on the basics, the nations are then classified as "Free", "Partly Free", or "Not Free". According to Reporters Without Borders, more than a third of the world's people live in countries where there is no press freedom. Overwhelmingly, these people live in countries where there is no system of democracy or where there are serious deficiencies in the democratic process. Freedom of the press is an extremely problematic problem/concept for most non-democratic systems of government since, in the modern age, strict control of access to information is critical to the existence of most non-democratic governments and their associated control systems and security apparatus. To this end, most non-democratic societies employ state-run news organizations to promote the propaganda critical to maintaining an existing political power base and suppress (often very brutally, through the use of police, military, or intelligence agencies) any significant attempts by the media or individual journalists to challenge the approved "government line" on contentious issues. In such countries, journalists operating on the fringes of what is deemed to be acceptable will very often find themselves the subject of considerable intimidation by agents of the state. This can range from simple threats to their professional careers (firing, professional blacklisting) to death threats, kidnapping, torture, and assassination. Reporters Without Borders reports that, in 2003, 42 journalists lost their lives pursuing their profession and that, in the same year, at least 130 journalists were in prison as a result of their occupational activities. In 2005, 63 journalists and 5 media assistants were killed worldwide. Examples include: - The Lira Baysetova case in Kazakhstan. - The Georgiy R. Gongadze case in Ukraine - In Nepal, Eritrea, and mainland China, journalists may spend years in jail simply for using the "wrong" word or photo. Regions closed to foreign reporters - Chechnya, Russia - Ogaden, Ethiopia - Jammu & Kashmir, India - Waziristan, Pakistan - Agadez, Niger - North Korea Central, Northern and Western Europe has a long tradition of freedom of speech, including freedom of the press. After World War II, Hugh Baillie, the president of United Press wire service based in the U.S., promoted freedom of news dissemination. In 1966 he called for an open system of news sources and transmission, and minimum of government regulation of the news. His proposals were aired at the Geneva Conference on Freedom of Information in 1948, but were blocked by the Soviets and by France. Until 1694, England had an elaborate system of licensing; the most recent was seen in the Licensing of the Press Act 1662. No publication was allowed without the accompaniment of a government-granted license. Fifty years earlier, at a time of civil war, John Milton wrote his pamphlet Areopagitica. In this work Milton argued forcefully against this form of government censorship and parodied the idea, writing "when as debtors and delinquents may walk abroad without a keeper, but unoffensive books must not stir forth without a visible jailer in their title." Although at the time it did little to halt the practice of licensing, it would be viewed later a significant milestone as one of the most eloquent defences of press freedom. Milton's central argument was that the individual is capable of using reason and distinguishing right from wrong, good from bad. In order to be able to exercise this ration right, the individual must have unlimited access to the ideas of his fellow men in “a free and open encounter." From Milton's writings developed the concept of the open marketplace of ideas, the idea that when people argue against each other, the good arguments will prevail. One form of speech that was widely restricted in England was seditious libel, and laws were in place that made criticizing the government a crime. The King was above public criticism and statements critical of the government were forbidden, according to the English Court of the Star Chamber. Truth was not a defense to seditious libel because the goal was to prevent and punish all condemnation of the government. Locke contributed to the lapse of the Licensing Act in 1695, whereupon the press needed no license. Still, many libels were tried throughout the 18th century, until "the Society of the Bill of Rights" led by John Horne Tooke and John Wilkes organised a campaign to publish Parliamentary Debates. This culminated in three defeats of the Crown in the 1770 cases of Almon, of Miller and of Woodfall, who all had published one of the Letters of Junius, and the unsuccessful arrest of John Wheble in 1771. Thereafter the Crown was much more careful in the application of libel; for example, in the aftermath of the Peterloo Massacre, Burdett was convicted, whereas by contrast the Junius affair was over a satire and sarcasm about the non-lethal conduct and policies of government. In Britain's American colonies, the first editors discovered their readers enjoyed it when they criticized the local governor; the governors discovered they could shut down the newspapers. The most dramatic confrontation came in New York in 1734, where the governor brought John Peter Zenger to trial for criminal libel After the publication of satirical attacks. The defense lawyers argued that according to English common law, truth was a valid defense against libel. The jury acquitted Zenger, who became the iconic American hero for freedom of the press. The result was an emerging tension between the media and the government. By the mid-1760s, there were 24 weekly newspapers in the 13 colonies, and the satirical attack on government became common features in American newspapers. John Stuart Mill in 1869 in his book On Liberty approached the problem of authority versus liberty from the viewpoint of a 19th-century utilitarian: The individual has the right of expressing himself so long as he does not harm other individuals. The good society is one in which the greatest number of persons enjoy the greatest possible amount of happiness. Applying these general principles of liberty to freedom of expression, Mill states that if we silence an opinion, we may silence the truth. The individual freedom of expression is therefore essential to the well-being of society. Mill wrote: - If all mankind minus one, were of one opinion, and one, and only one person were of the contrary opinion, mankind would be no more justified in silencing that one person, than he, if he had the power, would be justified in silencing mankind. Between September 4, 1770 and October 7, 1771 the kingdom of Denmark–Norway had the most unrestricted freedom of press of any country in Europe. This occurred during the regime of Johann Friedrich Struensee, whose second act was to abolish the old censorship laws. However, due to the great amount of mostly anonymous pamphlets published that was critical and often slanderous towards Struensee's own regime, he reinstated some restrictions regarding the freedom of press a year later, October 7, 1771. After the Italian unification in 1861, the Albertine Statute of 1848 was adopted as the constitution of the Kingdom of Italy. The Statute granted the freedom of the press with some restrictions in case of abuses and in religious matters, as stated in Article 28: - "The press shall be free, but the law may suppress abuses of this freedom. However, Bibles, catechisms, liturgical and prayer books shall not be printed without the prior permission of the Bishop." After the abolition of the monarchy in 1946 and the abrogation of the Statute in 1948, the Constitution of the Republic of Italy guarantees the freedom of the press, as stated in Article 21, Paragraphs 2 and 3: - "The press may not be subjected to any authorisation or censorship. Seizure may be permitted only by judicial order stating the reason and only for offences expressly determined by the law on the press or in case of violation of the obligation to identify the persons responsible for such offences." The Constitution allows the warrantless confiscation of periodicals in cases of absolute urgency, when the Judiciary cannot timely intervene, on the condition that a judicial validation must be obtained within 24 hours. Article 21 also gives restrictions against those publications considered offensive by public morality, as stated in Paragraph 6: - "Publications, performances, and other exhibits offensive to public morality shall be prohibited. Measures of preventive and repressive measure against such violations shall be established by law." Nazi Germany (1933–1945) In 1933 Freedom of the Press was suppressed in Hitler's Germany by the Reichstag Fire Decree of President Paul Von Hindenburg, just as Adolf Hitler was coming to power. Hitler largely suppressed freedom of the press through Joseph Goebbels' Ministry of Public Enlightenment and Propaganda. The Ministry acted as a central control-point for all media, issuing orders as to what stories could be run and what stories would be suppressed. Anyone involved in the film industry—from directors to the lowliest assistant—had to sign an oath of loyalty to the Nazi Party, due to opinion-changing power Goebbels perceived movies to have. (Goebbels himself maintained some personal control over every single film made in Nazi Europe.) Journalists who crossed the Propaganda Ministry were routinely imprisoned. Freedom of speech was first confirmed in Poland by the King Casimir the Great in 1347 in Wiślica Statutes. First polish newspapers were written, later replaced by print. The oldest Polish regular newspaper was pro-royal Merkuriusz Polski Ordynaryjny, founded in 1661. Later laws expanded freedom of the press further (i.e. March Constitution). One of the world's first freedom of the press acts was introduced in Sweden in 1766, mainly due to classical liberal member of parliament Anders Chydenius. Excepted and liable to prosecution was only vocal opposition to the King and the Church of Sweden. The Act was largely rolled back after King Gustav's coup d'état in 1772, restored after the overthrowing of his son, Gustav IV of Sweden in 1809, and fully recognized with the abolishment of the king's prerogative to cancel licenses in the 1840s. Congress shall make no law respecting an establishment of religion, or prohibiting the free exercise thereof; or abridging the freedom of speech, or of the press... The Indian Constitution, while not mentioning the word "press", provides for "the right to freedom of speech and expression" (Article 19(1) a). However this right is subject to restrictions under sub clause (2), whereby this freedom can be restricted for reasons of "sovereignty and integrity of India, the security of the State, friendly relations with foreign States, public order, preserving decency, preserving morality, in relation to contempt, court, defamation, or incitement to an offense". Laws such as the Official Secrets Act and Prevention of Terrorist Activities Act (PoTA) have been used to limit press freedom. Under PoTA, person could be detained for up to six months for being in contact with a terrorist or terrorist group. PoTA was repealed in 2006, but the Official Secrets Act 1923 continues. For the first half-century of independence, media control by the state was the major constraint on press freedom. Indira Gandhi famously stated in 1975 that All India Radio is "a Government organ, it is going to remain a Government organ..." With the liberalization starting in the 1990s, private control of media has burgeoned, leading to increasing independence and greater scrutiny of government. It ranks poorly at 140th rank out of 179 listed countries in the Press Freedom Index 2013 released by Reporters Without Borders (RWB). Analytically India's press freedom, as could be deduced by the Press Freedom Index, has constantly reduced since 2002, when it culminated in terms of apparent freedom, achieving a rank of 80 among the reported countries. Implications of new technologies Many of the traditional means of delivering information are being slowly superseded by the increasing pace of modern technological advance. Almost every conventional mode of media and information dissemination has a modern counterpart that offers significant potential advantages to journalists seeking to maintain and enhance their freedom of speech. A few simple examples of such phenomena include: - Satellite television versus terrestrial television: Whilst terrestrial television is relatively easy to manage and manipulate, satellite television is much more difficult to control as journalistic content can easily be broadcast from other jurisdictions beyond the control of individual governments. An example of this in the Middle East is the satellite broadcaster Al Jazeera. This Arabic-language media channel operates out of Qatar, whose government is relatively liberal with respect to many of its neighboring states. As such, its views and content are often problematic to a number of governments in the region and beyond. However, because of the increased affordability and miniaturisation of satellite technology (e.g. dishes and receivers) it is simply not practicable for most states to control popular access to the channel. - Web-based publishing (e.g., blogging) vs. traditional publishing: Traditional magazines and newspapers rely on physical resources (e.g., offices, printing presses) that can easily be targeted and forced to close down. Web-based publishing systems can be run using ubiquitous and inexpensive equipment and can operate from any global jurisdiction. To get control over web publications, nations and organisations are using geolocation and geolocation software. - Voice over Internet protocol (VOIP) vs. conventional telephony: Although conventional telephony systems are easily tapped and recorded, modern VOIP technology can employ low-cost strong cryptography to evade surveillance. As VOIP and similar technologies become more widespread they are likely to make the effective monitoring of journalists (and their contacts and activities) a very difficult task for governments. Naturally, governments are responding to the challenges posed by new media technologies by deploying increasingly sophisticated technology of their own (a notable example being China's attempts to impose control through a state-run internet service provider that controls access to the Internet) but it seems that this will become an increasingly difficult task as journalists continue to find new ways to exploit technology and stay one step ahead of the generally slower-moving government institutions that attempt to censor them. In May 2010, U.S. President Barack Obama signed legislation intended to promote a free press around the world, a bipartisan measure inspired by the murder in Pakistan of Daniel Pearl, the Wall Street Journal reporter, shortly after the September 11 attacks in 2001. The legislation, called the Daniel Pearl Freedom of the Press Act, requires the United States Department of State to expand its scrutiny of news media restrictions and intimidation as part of its annual review of human rights in each country. In 2012 the Obama Administration collected communication records from 20 separate home and office lines for Associated Press reporters over a two-month period, possibly in an effort to curtail government leaks to the press. The surveillance caused widespread condemnation by First Amendment experts and free press advocates, and led 50 major media organizations to sign and send a letter of protest to American attorney general Eric Holder. Organizations for press freedom - Article 19 - Canadian Journalists for Free Expression - The Committee to Protect Journalists - Electronic Frontier Foundation - Freedom House - Index on Censorship - Inter American Press Association - International Freedom of Expression Exchange - International Press Institute - Media Legal Defence Initiative - OSCE Representative on Freedom of the Media - Reporters Without Borders - Student Press Law Center - World Association of Newspapers and News Publishers - World Press Freedom Committee - Worldwide Governance Indicators |Part of a series on| |Censorship by country| - Article 10 of the European Convention on Human Rights - Areopagitica: a speech of Mr John Milton for the liberty of unlicensed printing to the Parliament of England - Chilling effect (term) - Cohen v. Cowles Media Co. — a ruling in the USA that a reporter's promise of a source's confidentiality may be enforced in court. - Declaration of Windhoek (1991) - Editorial independence - Free Speech, "The People’s Darling Privilege" - First Amendment to the United States Constitution - Freedom of speech - Freedom of the Press Act (1766) - Freedom of the Press (report) - Freedom of the press in the Russian Federation - Freedom of the press in the United States - Freedom of the press in Ukraine - Free speech in the media during the 2011 Libyan civil war - Gag order - International Freedom of Expression Exchange — "The largest online archive of information on press freedom violations", dating back to 1995 and covering more than 120 countries. - Journalism ethics and standards - Journaliste en danger - Journalistic standards - List of indices of freedom - Media blackout - Media transparency - News embargo - Section Two of the Canadian Charter of Rights and Freedoms - Photography is Not a Crime - Prior restraint - State media - Tunisia Monitoring Group - Virginia Declaration of Rights - World Press Freedom Day on May 3 - Worldwide Press Freedom Index - John Peter Zenger - Powe, L. A. Scot (1992). The Fourth Estate and the Constitution: Freedom of the Press in America. University of California Press. ISBN 9780520913165. - "Press Freedom Index 2014", Reporters Without Borders, 11 May 2014 - Press Freedom Index 2011-2012", Reporters Without Borders - "Description: Reporters Without Borders". The Media Research Hub. Social Science Research Council. 2003. Retrieved 23 September 2012. - Freedom House (2005). "Press Freedom Table (Press Freedom vs. Democracy ranks)". Freedom of the Press 2005. UK: World Audit. Retrieved 23 September 2012. - "Editor's daughter killed in mysterious circumstances", International Freedom of Expression Exchange (IFEX), 2 July 2002 - "Ukraine remembers slain reporter", BBC News, 16 September 2004 - "Do journalists have the right to work in Chechnya without accreditation?". Moscow Media Law and Policy Center. March 2000. Retrieved 2008-09-06. - "India praises McCain-Dalai Lama meeting". Washington, D.C.: WTOPews.com. July 27, 2008. Retrieved 2008-09-06. - Landay, Jonathan S. (March 20, 2008). "Radical Islamists no longer welcome in Pakistani tribal areas". McClatchy Washington Bureau. Retrieved 2008-09-06. - Eleonora W. Schoenebaum, ed. (1978), Political Profiles: The Truman Years, pp 16–17, Facts on File Inc., ISBN 9780871964533. - "British Press Freedom Under Threat", Editorial, New York Times, 14 November 2013. Retrieved 19 November 2013. - Alison Olson, "The Zenger Case Revisited: Satire, Sedition and Political Debate in Eighteenth Century America", Early American Literature, vol.35 no.3 (2000), pp: 223-245. - John Stuart Mill (1867). On Liberty. p. 10. - Laursen, John Christian (January 1998). "David Hume and the Danish Debate about Freedom of the Press in the 1770s". Journal of the History of Ideas 59 (1): 167–172. doi:10.1353/jhi.1998.0004. JSTOR 3654060. - "Lo Statuto Albertino" (PDF). The official website of the Presidency of the Italian Republic. - "The Italian Constitution" (PDF). The official website of the Presidency of the Italian Republic. - Jonathon Green and Nicholas J. Karolides, eds. (2009). Encyclopedia of Censorship. Infobase Publishing. pp. 194–96. - Skąd się wzięły gazety? - "The Freedom of the Press Act", Sveriges Riksdag - "The Swedish tradition of freedom of press" - "The World's First Freedom of Information Act (Sweden/Finland 1766)" - freedominfo.org, "Sweden" - "The Prevention of Terrorism Act 2002". - "Freedom of the Press". PUCL Bulletin (People's Union for Civil Liberties). July 1982. - "Press Freedom Index 2013". Reporters Without Borders. - "U.S. to Promote Press Freedom". New York Times. 17 May 2010. - Hicken, Jackie (15 May 2013). "Journalists push back against Obama administration for seizure of Associated Press records". Deseret News. Retrieved 16 May 2013. - Savage, Charlie; Leslie Kaufman (13 May 2013). "Phone Records of Journalists Seized by U.S.". The New York Times. Retrieved 16 May 2013. - Gant, Scott (2007). We're All Journalists Now: The Transformation of the Press and Reshaping of the Law in the Internet Age. New York: Free Press. ISBN 0-7432-9926-4. - Gardner, Mary A. The Inter American Press Association: Its Fight for Freedom of the Press, 1926–1960 (University of Texas Press, 2014) - George, Cherian. Freedom from the Press: Journalism and State Power in Singapore (2012) - McDonald, Blair (2015). "Freedom of Expression Revisited: Citizenship and Journalism in the Digital Era". Canadian Journal of Communication 40 (1). - Molnár, Peter, ed. Freedom of Speech and Freedom of Information Since the Fall of the Berlin Wall (Central European University Press, 2014) - Nord, Lars W., and Torbjörn Von Krogh. "The Freedom of The Press or The Fear Factor? Analysing Political Decisions and Non-Decisions in British Media Policy 1990-2012." Observatorio (OBS*) (2015) 9#1 pp: 01-16. - Starr, Paul (2004). The Creation of the Media: Political Origins of Modern Communications. New York: Basic Books. ISBN 0-465-08193-2. - Stockmann, Daniela. Media Commercialization and Authoritarian Rule in China (2012) - Thierer, Adam & Brian Anderson (2008). A Manifesto for Media Freedom. New York: Encounter Books. ISBN 1-59403-228-9. - Wilke, Jürgen (2013). Censorship and Freedom of the Press. Leibniz Institute of European History (IEG). |Look up freedom of the press in Wiktionary, the free dictionary.| |Wikimedia Commons has media related to Freedom of the press.| |Wikiquote has quotations related to: Freedom of the press| - Media Freedom Navigator Media Freedom Indices at a Glance - Risorse Etiche Publish and translate articles of independent journalists - the ACTivist Magazine - Paradox of media freedom in Pakistan - South East Europe Media Organisation - Banned Magazine, the journal of censorship and secrecy. - News and Free Speech — Newspaper Index Blog - Press Freedom - OSCE Representative on Freedom of the Media - MANA — the Media Alliance for New Activism - International Freedom of Expression Exchange — Monitors press freedom around the world - IPS Inter Press Service Independent news on press freedom around the world - The Reporters Committee for Freedom of the Press - Reporters Without Borders - Doha Center for Media Freedom - World Press Freedom Committee - Student Press Law Center - Union syndicale des journalistes CFDT - Mapping media freedom in Europe

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Each fraction can be reduced to itís simplest form, in which the numerator and denominator are as small as possible. To simplify a fraction you have to first find the greatest common factor of the numerator and denominator. You then have to divide the numerator and denominator by the biggest possible integer (whole number). For example if you have 12/18. The numerator (12) and the denominator (18) are both divisible by 6 (so the largest number that goes exactly into both 12 and 18 is 6 and therefore the greatest common factor is 6). Therefore if you divide the numerator and denominator by 6, you can write the fraction as (12/6)/(18/6) = 2/3. You can also simplify a fraction in steps. So 12/18 can also be written as 6/9 (the numerator and denominator can both be divided by 2). However, both 6 and 9 can also be divided by 3, so you can write 6/9 or even 2/3. 2/3 is the simplest form. More info Simplify fractions

http://www.fractioncalculator-online.com/simplify-fractions/

Influenza, or flu, is a respiratory infection caused by a variety of flu viruses. The most familiar aspect of the flu is the way it can “knock you off your feet” as it sweeps through entire communities. The flu differs in several ways from the common cold, a respiratory infection also caused by viruses. For example, people with colds rarely get fevers or headaches or suffer from the extreme exhaustion that flu viruses cause. The Centers for Disease Control and Prevention (CDC) estimates that 10 to 20 percent of Americans come down with the flu during each flu season, which typically lasts from November to March. Children are two to three times more likely than adults to get sick with the flu, and children frequently spread the virus to others. Although most people recover from the illness, CDC estimates that in the United States more than 100,000 people are hospitalized and about 36,000 people die from the flu and its complications every year. Flu outbreaks usually begin suddenly and occur mainly in the late fall and winter. The disease spreads through communities creating an epidemic. During the epidemic, the number of cases peaks in about 3 weeks and subsides after another 3 or 4 weeks. Half of the population of a community may be affected. Because schools are an excellent place for flu viruses to attack and spread, families with school-age children have more infections than other families, with an average of one-third of the family members infected each year. IMPORTANCE OF FLU Besides the rapid start of the outbreaks and the large numbers of people affected, the flu is an important disease because it can cause serious complications. Most people who get the flu get better within a week (although they may have a lingering cough and tire easily for a while longer). For elderly people, newborn babies, and people with certain chronic illnesses, however, the flu and its complications can be life-threatening. You can get the flu if someone around you who has the flu coughs or sneezes. You can get the flu simply by touching a surface like a telephone or door knob that has been contaminated by a touch from someone who has the flu. The viruses can pass through the air and enter your body through your nose or mouth. If you’ve touched a contaminated surface, they can pass from your hand to your nose or mouth. You are at greatest risk of getting infected in highly populated areas, such as in crowded living conditions and in schools. If you get infected by the flu virus, you will usually feel symptoms 1 to 4 days later. You can spread the flu to others before your symptoms start and for another 3 to 4 days after your symptoms appear. The symptoms start very quickly and may include - Body aches - Dry cough - Sore throat - Stuffy nose Typically, the fever begins to decline on the second or third day of the illness. The flu almost never causes symptoms in the stomach and intestines. The illness that some call “stomach flu” is not influenza. Usually, health care providers diagnose the flu on the basis of whether it is epidemic in the community and whether the person’s complaints fit the current pattern of symptoms. Health care providers rarely use laboratory tests to identify the virus during an epidemic. Health officials, however, monitor certain U.S. health clinics and do laboratory tests to determine which type of flu virus is responsible for the epidemic. The main way to keep from getting flu is to get a yearly flu vaccine. You can get the vaccine at your doctor’s office or a local clinic, and in many communities at workplaces, supermarkets, and drugstores. You must get the vaccine every year because it changes. Scientists make a different vaccine every year because the strains of flu viruses change from year to year. Nine to 10 months before the flu season begins, they prepare a new vaccine made from inactivated (killed) flu viruses. Because the viruses are killed, they cannot cause infection. The vaccine preparation is based on the strains of the flu viruses that are in circulation at the time. It includes those A and B viruses (see section below on types of flu viruses) expected to circulate the following winter. Sometimes, an unpredicted new strain may appear after the vaccine has been made and distributed to doctor’s offices and clinics. Because of this, even if you do get the flu vaccine, you still may get infected. If you do get infected, however, the disease usually is milder because the vaccine still will give you some protection.

http://blackdoctor.org/492/the-flu/

Posted by Jman on Thursday, April 25, 2013 at 2:12pm. ******I WILL GIVE MY ANSWERS AFTER I POST THE QUESTIONS****** 1. Momentum is a measure of ____. how hard it is to stop an object the amount of matter in an object the tendency of an object to change its motion the amount of force acting on an object Use the following equation to answer question 2: acceleration (in meters/second2) = net force (in newtons) mass (in kilograms) 2. A 300-N force acts on a 25-kg object. The acceleration of the object is ____. 3. The statement "for every action, there is an equal but opposite reaction" is a statement of ____. the law of conservation of momentum Newton's first law Newton's second law Newton's third law 4. A fixed, single pulley that is used to lift a block does which one of the following? doubles the force required to lift the block decreases the force required to lift the block makes the block easier to lift by changing the direction of the force needed to lift it decreases the force required and changes the direction of the force required 5. A heat engine ____. changes mechanical energy to thermal energy changes mechanical energy into electrical energy changes thermal energy into mechanical energy changes kinetic energy into thermal energy 6. Inertia _____. depends on direction depends on momentum resists a change in motion of an object both a and b 7. When a toy truck collides into a toy car, the momentum of _____ is the same before and after the collision. the truck multiplied by the car the truck plus the car 8. Air resistance _____ as you move faster. remains the same 9. A slanted surface used to raise an object is a(n) ____. (1 point) 10. A bar that is free to pivot about a fixed point is a _____. 11. When two or more simple machines work together they are called a(n) _____. 12. In a diesel engine, the fuel is ignited by _____. ☺☺☺☺☺THAT'S ALL THE Q's☺☺☺☺☺ Additional Details added 2 mins ago - Science Check answer - Jman, Thursday, April 25, 2013 at 2:16pm Answer This Question More Related Questions - science - Words to use: Accelerate - Sliding - Brake - Friction - Slipping - ... - calculus help check and help - if an object with mass m is dropped from rest, ... - phsyical science - the weight of an object is a) the force of earths gravity b) ... - Physical science - Is it possible for an object to have more kinetic energy but ... - physics - A 3 kg object traveling 6m/s east has a perfectly elastic collilsion ... - Physics - Object 1 (mass=5.5 kg) traveling in the +x direction at a speed of 12... - OPERATING SYSTEM - Hi, I need help with these questions. I have tried to find ... - Physics URGENT - Object 1 (mass=6.4 kg) traveling in the +x direction at a speed... - physics - I had wrote that was because I was told that in the beginning of the ... - physics - a moving object collides with an object initially at rest. is it ...

http://www.jiskha.com/display.cgi?id=1366913531

What Is A Watershed? A watershed is defined as a topographically delineated area drained by a stream. It is also referred to as a catchment and extends ridge top to ridge top. The interactions of the terrestrial and aquatic environments within the watershed determine watershed health. Much of Oregon’s natural resource management is based on watersheds. Why? First, watersheds are natural boundaries. Watershed management organizes and guides land and other resource use to provide desired goods and services without adversely affecting soil and water resources. This type of management also recognizes the interrelationships among land use, soil, water and linkages between upland and downstream areas. Finally, watershed management is based on scientific method; local citizens identify and prioritize projects, develop action plans, implement projects, and evaluate.

http://www.calapooia.org/about/watershed-assessment/

One of the two most important types of volcanoes, shield volcanoes are large and broad and have relatively gentle slopes. Kilauea on the island of Hawaii is good example of a shield volcano. Eruptions on shield volcanoes are far less explosive than those on composite volcanoes. That is because the basalt lava that erupts from shield volcanoes contains less silica, SiO2, and is therefore less sticky (less viscous) and doesn't "plug up" the volcano. Because the lava is runnier, it travels further from the crater before it cools, causing the shield-like shape of the volcano as many eruptions build up over time. Image: The summit of La Cumbre, a shield volcano on Fernandina Island, Galapagos Islands, as seen from Earth orbit (credit: NASA/SPL) Brian Cox describes the Solar System's largest volcano. Professor Brian Cox describes the biggest volcano in the Solar System, Olympus Mons on Mars. Plate boundaries are places of chaos and mineral wealth. Professor Iain Stewart explains how hotspots are a good demonstration of Earth's system of tectonic plates. As the plates move across the Earth's surface, they interact with one another at plate boundaries, which are places where earthquakes and volcanoes are common. Typically, plate boundaries are also places of great mineral wealth. Lava flows from this spectacular volcano have built the island of Hawaii. Professor Iain Stewart explains how Mount Kilauea's eruptions of lava have built up the island of Hawaii over millions of years as a magma plume known as a hotspot rises up through the Earth's crust. A shield volcano is a type of volcano usually built almost entirely of fluid magma flows. They are named for their large size and low profile, resembling a warrior's shield lying on the ground. This is caused by the highly fluid lava they erupt, which travels farther than lava erupted from stratovolcanoes. This results in the steady accumulation of broad sheets of lava, building up the shield volcano's distinctive form. Shield volcanoes shape is due to the low-viscosity magma of their mafic lava.

http://www.bbc.co.uk/science/earth/surface_and_interior/shield_volcano

Federalism is a political concept describing the practice whereby a group of members are bound together by agreement or covenant (Latin: foedus, covenant) with a ... About the US Government System of Federalism and How it Works to Divide Powers Between State and Federal Government. Federalism. Federalism is one of the most important and innovative concepts in the U.S. Constitution, although the word never appears there. Federalism is the sharing ... federalism definition. A system of government in which power is divided between a national (federal) government and various regional governments. Although the federal system seems to strike a perfect balance of power between national and local needs, federations still have internal power struggles. Federalism. A principle of government that defines the relationship between the central government at the national level and its constituent units at the regional ... Federalism is a system of government in which the same territory is controlled by two levels of government. Generally, an overarching national government governs ... Seen and Heard. What made you want to look up federalism? Please tell us where you read or heard it (including the quote, if possible). Federalism is the theory or advocacy of federal principles for dividing powers between member units and common institutions. Unlike in a unitary state ... Federalism in the United States is the constitutional relationship between U.S. state governments and the federal government of the United States.

https://www.search.com/reference/Federalism

In observational astronomy an Einstein ring, also known as an Einstein-Chwolson ring or Chwolson ring, is the deformation of the light from a source (such as a galaxy or star) into a ring through gravitational lensing of the source's light by an object with an extremely large mass (such as another galaxy or a black hole). This occurs when the source, lens, and observer are all aligned. The first complete Einstein ring, designated B1938+666, was discovered by collaboration between astronomers at the University of Manchester and NASA's Hubble Space Telescope in 1998. Gravitational lensing is predicted by Albert Einstein's theory of general relativity. Instead of light from a source traveling in a straight line (in three dimensions), it is bent by the presence of a massive body, which distorts spacetime. An Einstein Ring is a special case of gravitational lensing, caused by the exact alignment of the source, lens, and observer. This results in a symmetry around the lens, causing a ring-like structure. - is the gravitational constant, - is the mass of the lens, - is the speed of light, - is the angular diameter distance to the lens, - is the angular diameter distance to the source, and - is the angular diameter distance between the lens and the source. Note that, over cosmological distances in general. The bending of light by a gravitational body was predicted by Albert Einstein in 1912, a few years before the publication of General Relativity in 1916 (Renn et al. 1997). The ring effect was first mentioned in academic literature by Orest Chwolson in 1924. Einstein remarked upon this effect in 1936 in a paper prompted by a letter by a Czech engineer, R W Mandl , but stated Of course, there is no hope of observing this phenomenon directly. First, we shall scarcely ever approach closely enough to such a central line. Second, the angle β will defy the resolving power of our instruments.— Science vol 84 p 506 1936 In this statement, β is the Einstein Radius currently denoted by (as in the expression above). However, Einstein was only considering the chance of observing Einstein rings produced by stars, which is low; however, the chance of observing those produced by larger lenses such as galaxies or black holes is higher since the angular size of an Einstein ring increases with the mass of the lens. Known Einstein rings Hundreds of gravitational lenses are currently known. About half a dozen of them are partial Einstein rings with diameters up to an arcsecond, although as either the mass distribution of the lenses is not perfectly axially symmetrical, or the source, lens, and observer are not perfectly aligned, we have yet to see a perfect Einstein ring. Most rings have been discovered in the radio range. The degree of completeness needed for an image seen through a gravitational lens to qualify as an Einstein ring is yet to be defined. The first Einstein ring was discovered by Hewitt et al. (1988), who observed the radio source MG1131+0456 using the Very Large Array. This observation saw a quasar lensed by a nearer galaxy into two separate but very similar images of the same object, the images stretched round the lens into an almost complete ring. These dual images are another possible effect of the source, lens, and observer not being perfectly aligned. The first complete Einstein ring to be discovered was B1938+666, which was found by King et al. (1998) via optical follow-up with the Hubble Space Telescope of a gravitational lens imaged with MERLIN. The galaxy causing the lens at B1938+666 is an ancient elliptical galaxy, and the image we see through the lens is a dark dwarf satellite galaxy, which we would otherwise not be able to see with current technology. In 2005, the combined power of the Sloan Digital Sky Survey (SDSS) with the Hubble Space Telescope was used in the Sloan Lens ACS (SLACS) Survey to find 19 new gravitational lenses, 8 of which showed Einstein rings, these are the 8 shown in the image to the right. As of 2009 this survey has found 85 confirmed gravitational lenses, there is not yet a number for how many show Einstein rings. This survey is responsible for most of the recent discoveries of Einstein rings in the optical range, following are some examples which were found: - FOR J0332-3557, discovered by Remi Cabanac et al. in 2005, notable for its high redshift which allows us to use it to make observations about the early universe. - The "Cosmic Horseshoe" is a partial Einstein ring which was observed through the gravitational lens of LRG 3-757, a distinctively large Luminous Red Galaxy. It was discovered in 2007 by V. Belokurov et al. - SDSSJ0946+1006, the "double Einstein ring" was discovered by Raphael Gavazzi and Tomasso Treu in 2008, notable for the presence of multiple rings observed through the same gravitational lens, the significance of which is explained in the next section on extra rings. Another example is the radio/X-Ray Einstein ring around PKS 1830-211, which is unusually strong in radio. It was discovered in X-Ray by Varsha Gupta et al. at the Chandra X-Ray observatory It is also notable for being the first case of a quasar being lensed by an almost face-on spiral galaxy. Using the Hubble Space Telescope, a double ring has been found by Raphael Gavazzi of the STScI and Tommaso Treu of the University of California, Santa Barbara. This arises from the light from three galaxies at distances of 3, 6, and 11 billion light years. Such rings help in understanding the distribution of dark matter, dark energy, the nature of distant galaxies, and the curvature of the universe. The odds of finding such a double ring are 1 in 10,000. Sampling 50 suitable double rings would provide astronomers with a more accurate measurement of the dark matter content of the universe and the equation of state of the dark energy to within 10 percent precision. To the right is a simulation depicting a zoom on a Schwarzschild black hole in front of the Milky Way. The first Einstein ring corresponds to the most distorted region of the picture and is clearly depicted by the galactic disc. The zoom then reveals a series of 4 extra rings, increasingly thinner and closer to the black hole shadow. They are easily seen through the multiple images of the galactic disk. The odd-numbered rings correspond to points which are behind the black hole (from the observer's position) and correspond here to the bright yellow region of the galactic disc (close to the galactic center), whereas the even-numbered rings correspond to images of objects which are behind the observer, which appear bluer since the corresponding part of the galactic disc is thinner and hence dimmer here. |Wikimedia Commons has media related to Einstein Rings.| - Drakeford, Jason; Corum, Jonathan; Overbye, Dennis (March 5, 2015). "Einstein’s Telescope - video (02:32)". New York Times. Retrieved December 27, 2015. - Overbye, Dennis (March 5, 2015). "Astronomers Observe Supernova and Find They’re Watching Reruns". New York Times. Retrieved March 5, 2015. - "A Bull's Eye for MERLIN and the Hubble". University of Manchester. 27 March 1998. - "ALMA at Full Stretch Yields Spectacular Images". ESO Announcement. Retrieved 22 April 2015. - Belokurov, V.; et al. (January 2009). "Two new large-separation gravitational lenses from SDSS". Monthly Notices of the Royal Astronomical Society 392 (1): 104–112. arXiv:0806.4188. Bibcode:2009MNRAS.392..104B. doi:10.1111/j.1365-2966.2008.14075.x. Retrieved 2015-10-14. - Loff, Sarah; Dunbar, Brian (10 February 2015). "Hubble Sees A Smiling Lens". NASA. Retrieved 10 February 2015. - "Discovery of the First "Einstein Ring" Gravitational Lens". NRAO. 2000. Retrieved 2012-02-08. - Browne, Malcolm W. (1998-03-31). "'Einstein Ring' Caused by Space Warping Is Found". The New York Times. Retrieved 2010-05-01. - Vegetti, Simona; et al. (January 2012). "Gravitational detection of a low-mass dark satellite at cosmological distance". Nature 481 (7381): 341–343. arXiv:1201.3643. Bibcode:2012Natur.481..341V. doi:10.1038/nature10669. Retrieved 16 July 2014. - Bolton, A; et al. "Hubble, Sloan Quadruple Number of Known Optical Einstein Rings". Hubblesite. Retrieved 2014-07-16. - Auger, Matt; et al. (November 2009). "The Sloan Lens ACS Survey. IX. Colors, Lensing and Stellar Masses of Early-type Galaxies". The Astrophysical Journal 705 (2): 1099–1115. arXiv:0911.2471. Bibcode:2009ApJ...705.1099A. doi:10.1088/0004-637X/705/2/1099. Retrieved 16 July 2014. - Cabanac, Remi; et al. (2005-04-27). "Discovery of a high-redshift Einstein ring". Astronomy and Astrophysics 436 (2): L21–L25. arXiv:astro-ph/0504585. Bibcode:2005A&A...436L..21C. doi:10.1051/0004-6361:200500115. Retrieved 2014-07-15. - Belokurov, V.; et al. (December 2007). "The Cosmic Horseshoe: Discovery of an Einstein Ring around a Giant Luminous Red Galaxy". The Astrophysical Journal 671 (1): L9–L12. arXiv:0706.2326. Bibcode:2007ApJ...671L...9B. doi:10.1086/524948. Retrieved 2014-07-15. - Gavazzi, Raphael; et al. (April 2008). "The Sloan Lens ACS Survey. VI: Discovery and Analysis of a Double Einstein Ring". The Astrophysical Journal 677 (2): 1046–1059. arXiv:0801.1555. Bibcode:2008ApJ...677.1046G. doi:10.1086/529541. Retrieved 2014-04-15. - "Montage of the SDP.81 Einstein Ring and the lensed galaxy". Retrieved 9 June 2015. - Mathur, Smita; Nair, Sunita (20 July 1997). "X-Ray Absorption toward the Einstein Ring Source PKS 1830-211". The Astrophysical Journal 484: 140–144. arXiv:astro-ph/9703015. Bibcode:1997ApJ...484..140M. doi:10.1086/304327. Retrieved 16 July 2014. - Gupta, Varsha. "Chandra Detection of AN X-Ray Einstein Ring in PKS 1830-211". ResearchGate.net. Retrieved 16 July 2014. - Courbin, Frederic (August 2002). "Cosmic alignment towards the radio Einstein ring PKS 1830-211 ?". The Astrophysical Journal 575 (1): 95–102. arXiv:astro-ph/0202026. Bibcode:2002ApJ...575...95C. doi:10.1086/341261. Retrieved 16 July 2014. - Langston, G.I.; et al. (May 1989). "MG 1654+1346 - an Einstein Ring image of a quasar radio lobe". Astronomical Journal 97: 1283–1290. Bibcode:1989AJ.....97.1283L. doi:10.1086/115071. Retrieved 16 July 2014. - "Hubble Finds Double Einstein Ring". Hubblesite.org. Space Telescope Science Institute. Retrieved 2008-01-26. - Cabanac, R. A.; et al. (2005). "Discovery of a high-redshift Einstein ring". Astronomy and Astrophysics 436 (2): L21–L25. arXiv:astro-ph/0504585. Bibcode:2005A&A...436L..21C. doi:10.1051/0004-6361:200500115. (refers to FOR J0332-3357) - Chwolson, O (1924). "Über eine mögliche Form fiktiver Doppelsterne". Astronomische Nachrichten 221 (20): 329. Bibcode:1924AN....221..329C. doi:10.1002/asna.19242212003. (The first paper to propose rings) - Einstein, Albert (1936). "Lens-like Action of a Star by the Deviation of Light in the Gravitational Field" (PDF). Science 84 (2188): 506–507. Bibcode:1936Sci....84..506E. doi:10.1126/science.84.2188.506. PMID 17769014. (The famous Einstein Ring paper) - Hewitt, J (1988). "Unusual radio source MG1131+0456 - A possible Einstein ring". Nature 333: 537. Bibcode:1988Natur.333..537H. doi:10.1038/333537a0. - Renn, Jurgen; Sauer, Tilman; Stachel, John (1997). "The Origin of Gravitational Lensing: A Postscript to Einstein's 1936 Science paper". Science 275 (5297): 184–186. Bibcode:1997Sci...275..184R. doi:10.1126/science.275.5297.184. PMID 8985006. - King, L (1998). "A complete infrared Einstein ring in the gravitational lens system B1938 + 666". MNRAS 295: 41. arXiv:astro-ph/9710171. Bibcode:1998MNRAS.295L..41K. doi:10.1046/j.1365-8711.1998.295241.x.

https://en.wikipedia.org/wiki/Einstein_ring

7.4.2 Nitrogen Compounds The N cycle is integral to functioning of the Earth system and to climate (Vitousek et al., 1997; Holland et al., 2005a). Over the last century, human activities have dramatically increased emissions and removal of reactive N to the global atmosphere by as much as three to five fold. Perturbations of the N cycle affect the atmosphere climate system through production of three key N-containing trace gases: N2O, ammonia (NH3) and NOx (nitric oxide (NO) + nitrogen dioxide (NO2)). Nitrous oxide is the fourth largest single contributor to positive radiative forcing, and serves as the only long-lived atmospheric tracer of human perturbations of the global N cycle (Holland et al., 2005a). Nitrogen oxides have short atmospheric lifetimes of hours to days (Prather et al., 2001). The dominant impact of NOx emissions on the climate is through the formation of tropospheric ozone, the third largest single contributor to positive radiative forcing (Sections 2.3.6, 7.4.4). Emissions of NOx generate indirect negative radiative forcing by shortening the atmospheric lifetime of CH4 (Prather 2002). Ammonia contributes to the formation of sulphate and nitrate aerosols, thereby contributing to aerosol cooling and the aerosol indirect effect (Section 7.5), and to increased nutrient supply for the carbon cycle (Section 7.5). Ammonium and NOx are removed from the atmosphere by deposition, thus affecting the carbon cycle through increased nutrient supply (Section 18.104.22.168.3). Atmospheric concentrations of N2O have risen 16%, from about 270 ppb during the pre-industrial era to 319 ppb in 2005 (Figure 7.16a). The average annual growth rate for 1999 to 2000 was 0.85 to 1.1 ppb yr–1, or about 0.3% per year (WMO, 2003). The main change in the global N2O budget since the TAR is quantification of the substantial human-driven emission of N2O (Table 7.7; Naqvi et al., 2000; Nevison et al., 2004; Kroeze et al., 2005; Hirsch et al., 2006). The annual source of N2O from the Earth’s surface has increased by about 40 to 50% over pre-industrial levels as a result of human activity (Hirsch et al., 2006). Human activity has increased N supply to coastal and open oceans, resulting in decreased O2 availability and N2O emissions (Naqvi et al., 2000; Nevison et al., 2004). Figure 7.16. (a) Changes in the emissions of fuel combustion NOx and atmospheric N2O mixing ratios since 1750. Mixing ratios of N2O provide the atmospheric measurement constraint on global changes in the N cycle. (b) Changes in the indices of the global agricultural N cycle since 1850: the production of manure, fertilizer and estimates of crop N fixation. For data sources see http://www-eosdis.ornl.gov/ (Holland et al., 2005b) and http://www.cmdl.noaa.gov/. Figure adapted from Holland et al. (2005c). Since the TAR, both top-down and bottom-up estimates of N2O have been refined. Agriculture remains the single biggest anthropogenic N2O source (Bouwman et al., 2002; Smith and Conen, 2004; Del Grosso et al., 2005). Land use change continues to affect N2O and NO emissions (Neill et al., 2005): logging is estimated to increase N2O and NO emissions by 30 to 350% depending on conditions (Keller et al., 2005). Both studies underscore the importance of N supply, temperature and moisture as regulators of trace gas emissions. The inclusion of several minor sources (human excreta, landfills and atmospheric deposition) has increased the total bottom-up budget to 20.6 TgN yr–1 (Bouwman et al., 2002). Sources of N2O now estimated since the TAR include coastal N2O fluxes of 0.2 TgN yr–1 (±70%; Nevison et al., 2004) and river and estuarine N2O fluxes of 1.5 TgN yr–1 (Kroeze et al., 2005). Box model calculations show the additional river and estuarine sources to be consistent with the observed rise in atmospheric N2O (Kroeze et al., 2005). Top-down estimates of surface sources use observed concentrations to constrain total sources and their spatial distributions. A simple calculation, using the present-day N2O burden divided by its atmospheric lifetime, yields a global stratospheric loss of about 12.5 ± 2.5 TgN yr–1. Combined with the atmospheric increase, this loss yields a surface source of 16 TgN yr–1. An inverse modelling study of the surface flux of N2O yields a global source of 17.2 to 17.4 TgN yr–1 with an estimated uncertainty of 1.4 (1 standard deviation; Hirsch et al., 2006). The largest sources of N2O are from land at tropical latitudes, the majority located north of the equator. The Hirsch et al. inversion results further suggest that N2O source estimates from agriculture and fertilizer may have increased markedly over the last three decades when compared with an earlier inverse model estimate (Prinn et al., 1990). Bottom-up estimates, which sum individual source estimates, are more evenly distributed with latitude and lack temporal variability. However, there is clear consistency between top-down and bottom-up global source estimates, which are 17.3 (15.8–18.4) and 17.7 (8.5–27.7) TgN yr–1, respectively. Concentrations of NOx and reduced nitrogen (NHx = NH3 + ammonium ion (NH4+)) are difficult to measure because the atmospheric lifetimes of hours to days instead of years generate pronounced spatial and temporal variations in their distributions. Atmospheric concentrations of NOx and NHx vary more regionally and temporally than concentrations of N2O. Total global NOx emissions have increased from an estimated pre-industrial value of 12 TgN yr–1 (Holland et al., 1999; Galloway et al., 2004) to between 42 and 47 TgN yr–1 in 2000 (Table 7.7). Lamarque et al. (2005a) forecast them to be 105 to 131 TgN yr–1 by 2100. The range of surface NOx emissions (excluding lightning and aircraft) used in the current generation of global models is 33 to 45 TgN yr–1 with small ranges for individual sources. The agreement reflects the use of similar inventories and parametrizations. Current estimates of NOx emissions from fossil fuel combustion are smaller than in the TAR. Since the TAR, estimates of tropospheric NO2 columns from space by the Global Ozone Monitoring Experiment (GOME, launched in 1995) and the SCanning Imaging Absorption SpectroMeter for Atmospheric CHartographY (SCIAMACHY, launched in 2002) (Richter and Burrows, 2002; Heue et al., 2005) provide constraints on estimates of NOx emissions (Leue et al., 2001). Martin et al. (2003a) use GOME data to estimate a global surface source of NOx of 38 TgN yr–1 for 1996 to 1997 with an uncertainty factor of 1.6. Jaeglé et al. (2005) partition the surface NOx source inferred from GOME into 25.6 TgN yr–1 from fuels, 5.9 TgN yr–1 from biomass burning and 8.9 TgN yr–1 from soils. Interactions between soil emissions and scavenging by plant canopies have a significant impact on soil NOx emissions to the free troposphere: the impact may be greatest in subtropical and tropical regions where emissions from fuel combustion are rising (Ganzeveld et al., 2002). Boersma et al. (2005) find that GOME data constrain the global lightning NOx source for 1997 to the range 1.1 to 6.4 TgN yr–1. Comparison of the tropospheric NO2 column of three state-of-the-art retrievals from GOME for the year 2000 with model results from 17 global atmospheric chemistry models highlights significant differences among the various models and among the three GOME retrievals (Figure 7.17, van Noije et al., 2006). The discrepancies among the retrievals (10 to 50% in the annual mean over polluted regions) indicate that the previously estimated retrieval uncertainties have a large systematic component. Top-down estimates of NOx emissions from satellite retrievals of tropospheric NO2 are strongly dependent on the choice of model and retrieval. Figure 7.17. Tropospheric column NO2 from (a) satellite measurements and (b) atmospheric chemistry models. The maps represent ensemble average annual mean tropospheric NO2 column density maps for the year 2000. The satellite retrieval ensemble comprises three state-of-the-art retrievals from GOME; the model ensemble includes 17 global atmospheric chemistry models. These maps were obtained after smoothing the data to a common horizontal resolution of 5° × 5° (adapted from van Noije et al., 2006). Knowledge of the spatial distribution of NOx emissions has evolved significantly since the TAR. An Asian increase in emissions has been compensated by a European decrease over the past decade (Naja et al., 2003). Richter et al. (2005; see also Irie et al., 2005) use trends for 1996 to 2004 observed by GOME and SCIAMACHY to deduce a 50% increase in NOx emissions over industrial areas of China. Observations of NO2 in shipping lanes from GOME (Beirle et al., 2004) and SCIAMACHY (Richter et al., 2004) give values at the low end of emission inventories. Data from GOME and SCIAMACHY further reveal large pulses of soil NOx emissions associated with rain (Jaeglé et al., 2004) and fertilizer application (Bertram et al., 2005). All indices show an increase since pre-industrial times in the intensity of agricultural nitrogen cycling, the primary source of NH3 emissions (Figure 7.16b and Table 7.7; Bouwman et al., 2002). Total global NH3 emissions have increased from an estimated pre-industrial value of 11 TgN yr–1 to 54 TgN yr–1 for 2000 (Holland et al., 1999; Galloway et al., 2004), and are projected to increase to 116 TgN yr–1 by 2050. Table 7.7. Global sources (TgN yr–1) of NOx, NH3 and N2O for the 1990s. |Source ||NOx ||NH3 ||N2O | |TARa ||AR4b ||TARa ||AR4a ||TARc ||AR4 | |Anthropogenic sources || || || || || || | |Fossil fuel combustion & industrial processes ||33 (20–24) ||25.6 (21–28) ||0.3 (0.1–0.5) ||2.5d ||1.3/0.7 (0.2–1.8) ||0.7 (0.2–1.8)d | |Aircraft ||0.7 (0.2–0.9) ||– e (0.5–0.8) ||- ||- ||- ||- | |Agriculture ||2.3f (0–4) ||1.6g ||34.2 (16–48) ||35g (16–48) ||6.3/2.9 (0.9–17.9) ||2.8 (1.7–4.8)g | |Biomass and biofuel burning ||7.1 (2–12) ||5.9 (6–12) ||5.7 (3–8) ||5.4d (3–8) ||0.5 (0.2–1.0) ||0.7 (0.2–1.0)g | |Human excreta ||– ||– ||2.6 (1.3–3.9) ||2.6g (1.3–3.9) ||– ||0.2g (0.1–0.3)h | |Rivers, estuaries, coastal zones ||– ||– ||– ||– ||– ||1.7 (0.5–2.9)i | |Atmospheric deposition ||– ||0.3g ||– ||– ||– ||0.6j (0.3–0.9)h | |Anthropogenic total ||43.1 ||33.4 ||42.8 ||45.5 ||8.1/4.1 ||6.7 | |Natural sources || || || || || || | |Soils under natural vegetation ||3.3f (3–8) ||7.3j (5–8) ||2.4 (1–10) ||2.4g (1–10) ||6.0/6.6 (3.3–9.9) ||6.6 (3.3–9.0)g | |Oceans ||– ||– ||8.2 (3–16) ||8.2g (3–16) ||3.0/3.6 (1.0–5.7) ||3.8 (1.8–5.8)k | |Lightning ||5 (2–12) ||1.1–6.4 (3–7) ||– ||– ||– ||– | |Atmospheric chemistry ||<0.5 ||– ||– ||– ||0.6 (0.3–1.2) ||0.6 (0.3–1.2)c | |Natural total ||8.8 ||8.4–13.7 ||10.6 ||10.6 ||9.6/10.8 ||11.0 | |Total sources ||51.9 (27.2–60.9) ||41.8–47.1 (37.4–57.7) ||53.4 (40–70) ||56.1 (26.8–78.4) ||17.7/14.9 (5.9–37.5) ||17.7 (8.5–27.7) | The primary sink for NHx and NOx and their reaction products is wet and dry deposition. Estimates of the removal rates of both NHx and NOx are provided by measurements of wet deposition over the USA and Western Europe to quantify acid rain inputs (Hauglustaine et al., 2004; Holland et al., 2005a; Lamarque et al., 2005a). Chemical transport models represent the wet and dry deposition of NOx and NHx and their reaction products. A study of 29 simulations with 6 different tropospheric chemistry models, focusing on present-day and 2100 conditions for NOx and its reaction products, projects an average increase in N deposition over land by a factor of 2.5 by 2100 (Lamarque et al., 2005b), mostly due to increases in NOx emissions. Nitrogen deposition rates over Asia are projected to increase by a factor of 1.4 to 2 by 2030. Climate contributions to the changes in oxidized N deposition are limited by the models’ ability to represent changes in precipitation patterns. An intercomparison of 26 global atmospheric chemistry models demonstrates that current scenarios and projections are not sufficient to stabilise or reduce N deposition or ozone pollution before 2030 (Dentener et al., 2006).

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